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# Cosmological magnetic fields
## I Introduction
Magnetic fields seem to be everywhere that we can look in the universe, from our own sun out to high-redshift Lyman-$`\alpha `$ systems. The fields we observe (based on synchrotron radiation and Faraday rotation) in galaxies and clusters have been amplified by gravitational collapse and possibly also by dynamo mechanisms. They are either primordial, i.e. originating in the early universe and already present at the onset of structure formation, or they are protogalactic, i.e. generated by battery mechanisms during the initial stages of structure formation. One way to distinguish these possibilities would be to detect or rule out the presence of fields coherent on cosmological scales during recombination via their imprint on the cosmic microwave background (CMB) radiation. The new generation of CMB observations (especially the MAP and Planck satellites) may be able to achieve this.
The origin, evolution and cosmological impact of magnetic fields represent a fascinating challenge to theorists. I will discuss some aspects of this challenge in the following sections. (See for some other recent reviews.) Some basic facts from magnetohydrodynamics will be useful for the discussion.
Maxwell’s equations are
$$_{[\mu }F_{\nu \alpha ]}=0,_\nu F^{\mu \nu }=J^\mu ,$$
(1)
where $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$ is the field tensor, $`A_\mu `$ is the four-potential, and $`J_\mu `$ is the four-current. The field tensor is observer-independent, while the electric and magnetic fields depend on the observer’s motion:
$$E_\mu =F_{\mu \nu }u^\nu ,B_\mu =\frac{1}{2}\epsilon _{\mu \nu \alpha }F^{\nu \alpha },$$
(2)
where $`u^\mu `$ is the observer’s four-velocity, and $`\epsilon _{\mu \nu \alpha }`$ is the covariant permutation tensor in the observers’ rest space.
Ohm’s law is
$$h_{\mu \nu }J^\nu =\sigma F_{\mu \nu }u^\nu ,$$
(3)
where $`h_{\mu \nu }=g_{\mu \nu }+u_\mu u_\nu `$ projects orthogonal to $`u^\mu `$ (and $`g_{\mu \nu }`$ is the metric). For most of the history of the universe, the conductivity $`\sigma `$ is extremely high. In the magnetohydrodynamic limit, we have $`\sigma \mathrm{}`$ while the current remains finite, so that $`E_\mu 0`$. Thus the electric field in the particle frame vanishes: $`F_{\mu \nu }u^\nu =0`$. In the observer’s frame, with four velocity $`\stackrel{~}{u}^\mu =u^\mu +v^\mu `$, where $`v^\mu `$ is the relative velocity ($`v_\mu u^\mu =0`$) and we neglect terms $`O(v^2)`$, the electric field is of course not zero, but given by
$$\stackrel{~}{E}_\mu =\epsilon _{\mu \nu \alpha }v^\nu B^\alpha .$$
(4)
In this limit, Maxwell’s equations may be written as :
$`\text{D}^\mu B_\mu `$ $`=`$ $`0,`$ (5)
$`\omega ^\mu B_\mu `$ $`=`$ $`\frac{1}{2}J_\mu u^\mu ,`$ (6)
$`\text{curl}B_\mu `$ $`=`$ $`h_{\mu \nu }J^\nu +\epsilon _{\mu \nu \alpha }B^\nu \dot{u}^\alpha ,`$ (7)
$`h_{\mu \nu }\dot{B}^\nu `$ $`=`$ $`\frac{2}{3}\mathrm{\Theta }B_\mu +\sigma _{\mu \nu }B^\nu +\epsilon _{\mu \nu \alpha }B^\nu \omega ^\alpha ,`$ (8)
where $`\text{D}_\mu `$ is the projected covariant derivative, and $`\text{curl}B_\mu =\epsilon _{\mu \nu \alpha }\text{D}^\nu B^\alpha `$ is the covariant spatial curl. The kinematic quantities are $`\mathrm{\Theta }`$ (the volume expansion of $`u^\mu `$-flowlines), $`\omega _\mu `$ (vorticity), $`\dot{u}_\mu `$ (four-acceleration), and $`\sigma _{\mu \nu }`$ (shear).
The key equation is (8), which is the induction equation in covariant form. When contracted with $`B^\mu `$, it leads to the conservation equation for magnetic energy density:
$$\dot{\rho }_{\mathrm{mag}}+\frac{4}{3}\mathrm{\Theta }\rho _{\mathrm{mag}}=\sigma _{\mu \nu }\pi ^{\mu \nu },$$
(9)
where
$$\rho _{\mathrm{mag}}=\frac{1}{2}B_\mu B^\mu ,\pi _{\mu \nu }=\frac{1}{3}B^\alpha B_\alpha h_{\mu \nu }B_\mu B_\nu ,$$
(10)
are the energy density and anisotropic stress of the magnetic field. Typically, the term on the right of equation (9) may be neglected, in which case $`\rho _{\mathrm{mag}}`$ obeys the same evolution equation as isotropic radiation, so that
$$r\frac{\rho _{\mathrm{mag}}}{\rho _{\mathrm{rad}}}=\text{ constant}.$$
(11)
In a Friedmann universe, where $`\mathrm{\Theta }=3H=3\dot{a}/a`$ and $`a`$ is the scale factor, we have from equation (9) that
$$a^2B=\text{ constant},$$
(12)
where $`B=(B_\mu B^\mu )^{1/2}`$. If we choose $`a=1`$ at the present time, then $`a^2B`$ is the comoving magnitude of the magnetic field. Observations show that galactic and cluster fields are at the micro-Gauss level.
Nucleosynthesis imposes limits based on the way in which a magnetic field affects the expansion rate, the reaction rates and the electron phase density :
$$a^2B10^7\mathrm{G},$$
(13)
on cosmological scales. We can understand this limit qualitatively by requiring that $`\rho _{\mathrm{mag}}<\rho _{\mathrm{rad}}`$ at nucleosynthesis, which gives the right order of magnitude.
The upper limit from the CMB on a large-scale field is much tighter :
$$a^2B10^9\mathrm{G}.$$
(14)
This field strength corresponds to an energy density
$$\mathrm{\Omega }_{\mathrm{mag}}\frac{\rho _{\mathrm{mag}}}{\rho _{\mathrm{crit}}}10^5\mathrm{\Omega }_{\mathrm{rad}},$$
(15)
so that, roughly speaking, magnetic fields cannot induce large-angle perturbations in the CMB above the observed level.
## II Magnetogenesis and amplification
Protogalactic magnetogenesis, i.e. the creation of magnetic fields during the process of structure formation, essentially relies upon battery-type mechanisms in which the gradients of electron number density $`n_\mathrm{e}`$ and pressure $`p_\mathrm{e}`$ are not aligned. Ohm’s law (3) is modified and leads to the modified induction equation (in Newtonian form)
$$\frac{\stackrel{}{B}}{t}=\stackrel{}{}\times (\stackrel{}{v}\times \stackrel{}{B})+\alpha \stackrel{}{n}_\mathrm{e}\times \stackrel{}{p}_\mathrm{e},$$
(16)
where $`\alpha `$ is a constant. It follows that if the gradient terms are non-aligned (as happens for example when shock waves develop in collapsing clouds), then nonzero $`B`$ can be generated. Very small fields are generated in this way, and typically require strong dynamo-type amplification in order to reach the currently observed levels.
A seed magnetic field, whether generated by battery mechanisms or already present in the form of a primordial field, is amplified adiabatically during gravitational collapse, simply by the fact that field lines are frozen into the plasma, and compression of the plasma results in compression of flux lines. This adiabatic compression is weak, with growth roughly given by
$$B\delta ^{2/3},$$
(17)
where $`\delta =\delta \rho /\rho `$ is the fractional over-density of the cloud \[this neglects the shear term in equation (9)\]. If the observed galactic fields ($`10^6`$ G) are the result only of adiabatic compression, then the seed field required could be up to $`10^9`$ G (comoving). This is at the level of the CMB limit on large-scales.
If the seed field is much weaker, then a stronger amplification is required – and the prime candidate mechanism for this is the galactic dynamo . This is based on differential rotation and turbulence, whereby small-scale magnetic fields are amplified via parametric resonance. The key issue of how efficient the dynamo is, has not been settled. There is therefore a large uncertainty in the amount of amplification that can be achieved, and thereby in the size of seed field that is necessary. In general qualitative terms, the seed field will be much less than that required for purely adiabatic compression. In terms of the $`r`$-factor in equation (11), a seed without dynamo amplification requires $`r10^{14}`$, whereas a seed with dynamo amplification could have $`r`$ as low as $`10^{34}`$ (this may be further reduced in the presence of a cosmological constant ).
Primordial magnetogenesis is the creation of magnetic fields in the early universe, before the process of structure formation. Many mechanisms have been proposed, based mainly on phase transitions before recombination, or on inflation. In phase transitions such as the QCD and EW transitions, local charge separation can arise, creating local currents that can generate (hyper-)magnetic fields . Other proposals include bubble-wall collisions, which produce phase gradients that can source gauge fields .
These mechanisms produce fields coherent on sub-Hubble scales. In order to generate super-Hubble scale fields, one requires inflationary models , or pre big bang models based on string theory , in which vacuum fluctuations of the field are amplified via the dilaton. Inflation stretches perturbations beyond the Hubble horizon and thus can in principle generate magnetic fields on large scales. There is however a problem in that vector perturbations are extremely small in standard models, essentially because the vector gauge field does not couple gravitationally to a conformally flat metric. One needs to break conformal invariance by new high-energy couplings of the photon (or to break gauge invariance). An example of such a coupling is provided by the Lagrangian for scalar electrodynamics:
$$=\frac{1}{4}F_{\mu \nu }F^{\mu \nu }(𝒟_\mu \varphi )^{}𝒟^\mu \varphi V(\varphi \varphi ^{}),$$
(18)
where $`\varphi `$ is the charged scalar field, and $`𝒟_\mu =_\mu \mathrm{i}eA_\mu `$ is the gauge-covariant derivative, with $`e`$ the coupling constant.
Inflation is often followed by a preheating period in which coherent oscillations of the inflaton produce parametric resonant amplification of perturbations. Since the inflaton is coherent on super-Hubble scales, this amplification can in principle affect super-Hubble scales, without in any way violating causality . Magnetic fields arising from inflationary fluctuations could thus in principle be amplified via preheating .
## III Magnetic fields and the CMB
In the absence of any preferred model of primordial magnetogenesis, and in view of the complexities of magnetohydrodynamics during structure formation (especially the dynamo mechanism), we need cosmological observational tests for deciding whether magnetogenesis is primordial or protogalactic. If magnetic fields could be detected in the voids between galactic clusters, this would be very strong evidence for a primordial origin.
The other key observational test is provided by the CMB. Dynamically significant magnetic fields present during recombination must be primordial. These primordial fields have various effects on the CMB.
In the absence of a magnetic field, the tightly coupled baryon-photon fluid undergoes longitudinal acoustic oscillations in density and velocity perturbations, with
$$\delta ,v\mathrm{exp}(\mathrm{i}kc_\mathrm{s}\eta ),$$
(19)
where $`c_\mathrm{s}`$ is the sound speed and $`\eta `$ is conformal time. A magnetic field splits these modes into 3 types:
(a) fast magnetosonic waves, which are like sound waves, but with increased speed,
$$c_\mathrm{s}^2c_\mathrm{s}^2+c_\mathrm{a}^2\mathrm{sin}^2\theta ,$$
(20)
where $`c_\mathrm{a}^2=\rho _{\mathrm{mag}}/\rho `$ is the Alfvén speed squared and $`\theta `$ is the angle between $`\stackrel{}{B}`$ and the propagation direction;
(b) slow magnetosonic waves, which have speed $`c_\mathrm{a}\mathrm{cos}\theta `$ and are partly transverse in velocity;
(c) incompressible Alfvén waves, whose speed is the same as the slow magnetosonic waves, and for which $`\delta =0`$.
The fast magnetosonic waves have a direct and simple, though small, effect on the acoustic peaks in CMB temperature anisotropies , based on the modification of the sound speed. The effect of Alfvén modes on CMB anisotropies has also been calculated .
Fast magnetosonic modes suffer diffusion damping just like the non-magnetized acoustic modes. The slow magnetosonic and Alfvén modes by contrast can be overdamped and survive on scales below the Silk scale . This could play an interesting role in structure formation.
In general, the dissipation of magnetized fluctuations injects non-thermal energy into the photon spectrum, which introduces chemical-potential and Compton distortions ($`\mu `$ and $`y`$ distortions) in the CMB blackbody. Upper limits on these distortions provided by the FIRAS experiment on COBE then place upper limits on the magnetic field strength :
$$a^2B10^8\text{G on scales }0.5600\mathrm{kpc}.$$
(21)
The anisotropic stress $`\pi _{\mu \nu }`$ induced by a magnetic field can source gravitational wave perturbations during recombination. This can be seen through the wave equation that governs the transverse traceless magnetic part of the Weyl tensor, $`H_{\mu \nu }`$, which provides a covariant description of gravitational waves :
$$\ddot{H}_{\mu \nu }+\text{D}^2H_{\mu \nu }=7H\dot{H}_{\mu \nu }+2\rho (1w)H_{\mu \nu }+2H\text{curl}\pi _{\mu \nu },$$
(22)
where $`\text{curl}\pi _{\mu \nu }=\epsilon _{\alpha \beta (\mu }\text{D}^\alpha \pi ^\beta _{\nu )}`$ is the covariant spatial tensor curl and $`w=p/\rho `$. In order to keep the tensor contribution to CMB temperature anisotropies within the observed limits, this places upper limits on the magnetic field .
Magnetic fields have an important effect on the polarization of the CMB via Faraday rotation . Linearly polarized radiation with frequency $`\nu `$ and wave vector $`\stackrel{}{e}`$, in a magnetized plasma with free electron density $`n_\mathrm{e}`$, has its plane of polarization rotated through an angle $`\phi `$, where
$$\frac{d\phi }{dt}\frac{n_\mathrm{e}}{\nu ^2}\stackrel{}{B}\stackrel{}{e}.$$
(23)
For a given line of sight $`\stackrel{}{e}`$, the polarization angle $`\phi `$ may be measured at different frequencies, thus providing in principle a measure of the magnetic field strength. The Planck experiment may be able to detect a field at the $`10^9`$ G level. An indirect effect of Faraday rotation is to depolarize the CMB on small angular scales, leading to a reduction in damping and thus a small increase in power in the temperature anisotropies .
Perhaps more significant than the small quantitative effects on polarization angle and on small-scale temperature anisotropies is an intriguing correlation introduced by magnetic fields . Scalar perturbations can only generate E-type polarization, while tensor perturbations generate both E- and B-type polarization. A magnetic field also generates both E- and B-type polarization, but in addition, it induces a correlation via the Faraday rotation coupling in the evolution equations for polarization. This means that the B-type polarization will be correlated with the temperature anisotropies. Such a correlation does not arise in the context of statistical isotropy, but a large-scale magnetic field breaks the isotropy and produces an novel signature, which may be more accessible to observation.
Another potentially important (although probably extremely small) effect on CMB temperature anisotropies arises from the general relativistic interaction between gravity and electromagnetism, whereby electromagnetic radiation may be induced from a magnetic field by gravitational waves .
## IV Magnetized structure formation
The effects of a weak cosmological magnetic field on structure formation in the linear regime are necessarily very small. The pioneering analysis was given in (see also ). In the matter era on sub-Hubble scales, a Newtonian approach is justified, based on the magnetized Euler equation
$$\frac{\stackrel{}{v}}{\eta }+aH\stackrel{}{v}=c_\mathrm{s}^2\stackrel{}{}\delta \stackrel{}{}\mathrm{\Phi }+\frac{1}{\rho }(\stackrel{}{}\times \stackrel{}{B})\times \stackrel{}{B},$$
(24)
where $`\mathrm{\Phi }`$ is the gravitational potential perturbation. The standard, non-magnetized adiabatic growing mode of density perturbations is slightly damped by magnetism :
$$\delta a^n,n=\frac{1}{4}\left[1+5\sqrt{1\alpha _{\mathrm{mag}}k^2}\right],$$
(25)
where $`\alpha _{\mathrm{mag}}`$ is a constant determined by $`c_\mathrm{a}^2`$, and $`k`$ is the wave number. New non-adiabatic constant and decaying modes are also introduced by the magnetic field. A magnetic field can induce density perturbations in a homogeneous medium, although it cannot on its own reproduce the features of the observed power spectrum . An analysis of the complex dynamics of magnetized damping during recombination shows that incompressible and slow magnetosonic modes may survive on scales well below the Silk scale, and this could lead to interesting variations on the non-magnetized scenario of structure formation. These small-scale modes that survive damping could seed early star or galaxy formation and could also precipitate fragmentation of early structures.
The magnetic field also acts as a source of incompressible rotational instabilities, which satisfy the wave equation
$`\ddot{W}_\mu +\left[{\displaystyle \frac{c_\mathrm{a}^2}{3(1+w)}}\right]\text{D}^2W_\mu `$ $`=`$ $`(43w)H\dot{W}_\mu `$ (26)
$`+\frac{1}{2}\rho \left[17w+3c_\mathrm{s}^2(1+w)\right]W_\mu ,`$ (27)
where $`D^\mu W_\mu =0`$. On small scales, these vortices may have some interesting effects on structure formation. Magnetic fields can generate not only vorticity, but also anisotropic distortion in the density distribution .
On super-Hubble scales, a fully general relativistic analysis is needed, and this is developed in (see for a dynamical-systems analysis of the equations). During the radiation era, the non-magnetized adiabatic growing mode is incorrectly predicted to suffer small magnetic damping via an analysis which does not incorporate all relativistic effects. In fact, there is a crucial magneto-curvature coupling, which arises from the non-commutation of the projected covariant derivatives of the magnetic field:
$$\text{D}_{[\mu }\text{D}_{\nu ]}B_\alpha =\frac{1}{2}_{\mu \nu \alpha \beta }B^\beta \epsilon _{\mu \nu \beta }\omega ^\beta \dot{B}_\alpha ,$$
(28)
where the projected curvature tensor is
$$_{\mu \nu \alpha \beta }=h_\mu {}_{}{}^{\sigma }h_{\nu }^{}{}_{}{}^{\chi }h_{\alpha }^{}{}_{}{}^{\gamma }h_{\beta }^{}{}_{}{}^{\delta }R_{\sigma \chi \gamma \delta }^{}V_{\mu \alpha }V_{\nu \beta }+V_{\mu \beta }V_{\nu \alpha },$$
(29)
with
$$V_{\mu \nu }=\frac{1}{3}\mathrm{\Theta }h_{\mu \nu }+\sigma _{\mu \nu }+\epsilon _{\mu \nu \alpha }\omega ^\alpha .$$
(30)
This coupling combines with the tension of the magnetic force-lines to reverse the damping effect and leads to a small enhancement of the growing mode, which satisfies the equation
$$a^2\frac{d^2\delta }{da^2}(2c_\mathrm{a}^2)\delta =c_\mathrm{a}^2\left(C+2a^2\right),$$
(31)
where $`C`$ is a constant and $`=h^{\mu \nu }h^{\alpha \beta }_{\mu \alpha \nu \beta }`$ is the projected curvature scalar.
The coupling between magnetism and curvature essentially injects the elastic properties of magnetic field lines into space itself, and can lead to rather unexpected dynamical and kinematical effects .
## V Conclusion
Cosmic magnetic fields provide a fascinating set of unsolved problems challenging theorists in cosmology. Not only do we need to resolve the key question as to whether these fields are primordial or protogalactic in origin, but we also need to develop a satisfactory theory of magnetogenesis and amplification. Furthermore, there are a number of open issues in calculating the magnetic effects on structure formation and on CMB anisotropies. The required theoretical developments will be driven by advances in observations, both directly of magnetic fields beyond the galactic scale, and indirectly via future advances in CMB observations and large-scale structure surveys.
Acknowledgements
I thank Christos Tsagas and Alejandra Kandus for many helpful discussions. I am grateful to the organisers of the conference for their wonderful hospitality and kindness, with special thanks to Sayan Kar and Naresh Dadhich. It was an honour to speak at Kharagpur, on the site of resistance to a former colonial prison, where many freedom fighters refused to surrender.
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# Gravitational Collapse of Cylindrical Shells Made of Counter-Rotating Dust Particles
## I Introduction
Gravitational collapse of a realistic body has been one of the most thorny and important problems in Einstein’s theory of General Relativity. Due to the complexity of the Einstein field equations, the problem even in simple cases, such as, spacetimes with spherical symmetry, is still not well understood , and new phenomena keep emerging . Particularly, in 1991 Shapiro and Teukolsky studied numerically the problem of a dust spheroid, and found that only the spheroid is compact enough, a black hole can be formed. Otherwise, the collapse most likely ends with a naked singularity. Later, Barrabés, Israel and Letelier constructed an analytical model of a collapsing convex thin shell and found that in certain cases no apparent horizons are formed, too . Their results were soon generalized to more general cases . However, since in all the cases considered by them, the external gravitational field of the collapsing shell is not known, one cannot exclude, similar to the ST case, the formation of an outer event horizon . Since then, the gravitational collapse with non-spherical symmetry has been attracting more and more attention. In particular, by studying the collapse of a cylindrical shell that is made of counter-rotating particles, Apostolatos and Thorne (AT) showed analytically that the centrifugal forces associated with an arbitrarily small amount of rotation, by themselves, without the aid of any pressure, can halt the collapse at some non-zero, minimum radius, and the shell will then oscillate until it settles down at some final, finite radius, whereby a spacetime singularity is prevented from forming on the symmetry axis . Soon after AT’s work, Shapiro and Teukolsky studied numerically the gravitational collapse of rotating spheroids, and found that the rotation indeed significantly modifies the evolution when it is sufficiently large. However, for small enough angular momentum, their simulations showed that spindle singularities appeared to arise without apparent horizons, too. Hence, it is possible that even spheroids with some angular momentum may still form naked singularities .
It should be noted that in the AT work it was considered only the case where the shell has zero total angular momentum and is momentarily static and radiation-free. In a realistic case, the spacetime has neither cylindrical symmetry nor zero angular momentum, and gravitational and particle radiations are always expected to occur. As a generalization of the AT work, in this paper we shall consider the case where cylindrical shell radiates gravitational waves and massless particles, as it is collapsing, while still keep the requirement that the total angular momentum of the shell be zero. Specifically, the paper is organized as follows: In Sec. $`II`$, the formulas for a general dynamic timelike thin shell that connects two arbitrary cylindrical regions are given, using Israel’s formula , while in Sec. $`III`$, a collapsing thin shell made of counter-rotating dust particles is studied. To model the particle radiation of the shell, we consider the case where the spacetime outside the shell is described by an out-going radiation fluid . The paper is ended with Sec. $`IV`$, where our main conclusions are presented.
## II Dynamics of Cylindrical Thin Shells Without Rotation
Both static and dynamic cylindrical thin shells with zero total angular momentum have been studied previously. However, in most of these studies a specific form of metric was usually assumed, which is valid only in some particular cases, such as, the spacetime is vacuum outside and inside the shell . In this section, we shall give a general treatment that is valid for any dynamic timelike cylindrical thin shell, connecting two arbitrary cylindrical regions.
To begin with, let us consider the cylindrical spacetimes described by the metric,
$$ds_{}^2=f^{}(t,r)dt^2g^{}(t,r)dr^2h^{}(t,r)dz^2l^{}(t,r)d\phi ^2,$$
(1)
where $`\{x^\mu \}\{t,r,z,\phi \},(\mu =0,1,2,3)`$, are the usual cylindrical coordinates. For the spacetimes to be cylindrical, several criteria have to be satisfied . When the symmetry axis is regular, those conditions are easily imposed. However, when it is singular, it is still not clear which kind of conditions should be imposed .
In general the spacetimes described by Eq.(1) have two Killing vectors, one is associated with the invariant translations along the symmetry axis, $`\xi _{(z)}=z`$, where $`z`$ is the Killing coordinate length with $`\mathrm{}<z<+\mathrm{}`$, and the other is associated with the invariant rotations about the axis, $`\xi _{(\phi )}=\phi `$ with $`0\phi 2\pi `$, where the hypersurface $`\phi =0`$ is identical with the one $`\phi =2\pi `$. Clearly, for the metric given above, the two Killing vectors are orthogonal. Consequently, the metric represents spacetimes without rotation, and the polarization of gravitational waves has only one degree of freedom .
Assume that a given spacetime is divided by a hypersurface $`\mathrm{\Sigma }`$ into two regions, say, $`V^\pm `$, where the region $`V^{}`$ is described by the metric (1), while the region $`V^+`$ is described by the metric
$$ds_+^2=f^+(T,R)dT^2g^+(T,R)dR^2h^+(T,R)dz^2l^+(T,R)d\phi ^2,$$
(2)
where $`\{x^{+\mu }\}\{T,R,z,\phi \},(\mu =0,1,2,3)`$, is another set of the cylindrical coordinates. The hypersurface $`\mathrm{\Sigma }`$ in the coordinates $`x^{\pm \mu }`$ is given, respectively, by
$$r=r_0(t),R=R_0(T).$$
(3)
On the surface, the metrics (1) and (2) reduce, respectively, to
$`ds_{}^2|_{r=r_0(t)}`$ $`=`$ $`\left[f^{}(t,r_0(t))g^{}(t,r_0(t))r_{}^{}{}_{0}{}^{2}(t)\right]dt^2h^{}(t,r_0(t))dz^2l^{}(t,r_0(t))d\phi ^2,`$ (4)
$`ds_+^2|_{R=R_0(T)}`$ $`=`$ $`\left[f^+(T,R_0(T))g^+(T,R_0(T))R_{}^{}{}_{0}{}^{2}(T)\right]dT^2h^+(T,R_0(T))dz^2l^+(T,R_0(T))d\phi ^2,`$ (5)
where a prime denotes the ordinary differentiation with respect to the indicated argument. In this paper, we shall consider only the case where $`\mathrm{\Sigma }`$ is timelike. Then, if we choose the intrinsic coordinates of the hypersurface as $`\{\xi ^a\}=\{\tau ,z,\phi \},(a=1,2,3)`$, where $`\tau `$ denotes the proper time of the surface, we find that the metric on the hypersurface can be written as,
$$ds^2|_\mathrm{\Sigma }=\gamma _{ab}d\xi ^ad\xi ^b=d\tau ^2h(\tau )dz^2l(\tau )d\phi ^2,$$
(6)
where
$`d\tau `$ $`=`$ $`\left[f^{}(t,r_0(t))g^{}(t,r_0(t))r_{}^{}{}_{0}{}^{2}(t)\right]^{1/2}dt=\left[f^+(T,R_0(T))g^+(T,R_0(T))R_{}^{}{}_{0}{}^{2}(T)\right]^{1/2}dT,`$ (7)
$`h(\tau )`$ $``$ $`h^{}(t,r_0(t))=h^+(T,R_0(T)),l(\tau )l^{}(t,r_0(t))=l^+(T,R_0(T)),`$ (8)
where the function dependence of $`\tau `$ on $`t`$ and $`T`$ is given by the first equation. Note that in writing the above expressions, we had chosen $`d\tau ,dT`$ and $`dt`$, without loss of generality, to have the same sign, and already applied the first junction conditions,
$$ds_{}^2|_{r=r_0(t)}=ds_+^2|_{R=R_0(T)}.$$
(9)
It can be shown that the unit spacelike normal vector to the hypersurface $`\mathrm{\Sigma }`$ in the coordinates $`x^{\pm \mu }`$ is given, respectively, by
$`n_\mu ^+`$ $`=`$ $`\left[{\displaystyle \frac{f^+g^+}{f^+g^+R_{}^{}{}_{0}{}^{2}(T)}}\right]^{1/2}\left\{R_{}^{}{}_{0}{}^{}(T)\delta _\mu ^T+\delta _\mu ^R\right\},`$ (10)
$`n_\mu ^{}`$ $`=`$ $`\left[{\displaystyle \frac{f^{}g^{}}{f^{}g^{}r_{}^{}{}_{0}{}^{2}(t)}}\right]^{1/2}\left\{r_{}^{}{}_{0}{}^{}(t)\delta _\mu ^t+\delta _\mu ^r\right\}.`$ (11)
Then, the non-vanishing components of the extrinsic curvature tensor $`K_{ab}^\pm `$, defined by<sup>*</sup><sup>*</sup>* Note that in this paper the definition for the extrinsic curvature tensor is different from that of Israel by a “$``$” sign .
$$K_{ab}=n_\alpha \left(\frac{^2x^\alpha }{\xi ^a\xi ^b}+\mathrm{\Gamma }_{\beta \delta }^\alpha \frac{x^\beta }{\xi ^a}\frac{x^\delta }{\xi ^b}\right),$$
(12)
are given by
$`K_{\tau \tau }^+`$ $`=`$ $`{\displaystyle \frac{\left(f^+g^+\right)^{1/2}}{2\left[f^+g^+R_{}^{}{}_{0}{}^{2}(T)\right]^{3/2}}}\{{\displaystyle \frac{f_{,R}^+}{g^+}}+({\displaystyle \frac{f_{,T}^+}{f^+}}2{\displaystyle \frac{g_{,T}^+}{g^+}})R_{}^{}{}_{0}{}^{}(T)`$ (14)
$`+(2{\displaystyle \frac{f_{,R}^+}{f^+}}{\displaystyle \frac{g_{,R}^+}{g^+}})R_{}^{}{}_{0}{}^{2}(T)+{\displaystyle \frac{g_{,R}^+}{f^+}}R_{}^{}{}_{0}{}^{3}(T)2R_{}^{\prime \prime }{}_{0}{}^{}(T)\},`$
$`K_{zz}^+`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{f^+g^+}{f^+g^+R_{}^{}{}_{0}{}^{2}(T)}}\right]^{1/2}\left\{{\displaystyle \frac{h_{,R}^+}{g^+}}+{\displaystyle \frac{h_{,T}^+}{f^+}}R_{}^{}{}_{0}{}^{}(T)\right\},`$ (15)
$`K_{\phi \phi }^+`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{f^+g^+}{f^+g^+R_{}^{}{}_{0}{}^{2}(T)}}\right]^{1/2}\left\{{\displaystyle \frac{l_{,R}^+}{g^+}}+{\displaystyle \frac{l_{,T}^+}{f^+}}R_{}^{}{}_{0}{}^{}(T)\right\},`$ (16)
where $`f_{,T}^+f^+(T,R)/T`$ etc., and $`K_{ab}^{}`$ can be obtained from the above expressions by the replacement
$$f^+,g^+,h^+,l^+,R_0(T),T,Rf^{},g^{},h^{},l^{},r_0(t),t,r.$$
(17)
In terms of $`K_{ab}^\pm `$ and $`\gamma _{ab}`$, the surface energy-momentum tensor, $`\tau _{ab}`$, is defined as ,
$$\tau _{ab}=\frac{1}{\kappa }\left\{\left[K_{ab}\right]^{}\gamma _{ab}[K]^{}\right\},$$
(18)
where $`\kappa [8\pi G/c^4]`$ is the Einstein constant, and
$$\left[K_{ab}\right]^{}K_{ab}^+K_{ab}^{},[K]^{}\gamma ^{ab}\left[K_{ab}\right]^{}.$$
(19)
Inserting Eq.(14) and the corresponding expressions for $`K_{ab}^{}`$ into Eq.(18), we find that $`\tau _{ab}`$ can be written in the form
$$\tau _{ab}=\rho w_aw_b+p_zz_az_b+p_\phi \phi _a\phi _b,(a,b=\tau ,z,\phi ),$$
(20)
where
$`\rho `$ $`=`$ $`{\displaystyle \frac{1}{\kappa }}\left\{{\displaystyle \frac{\left[K_{zz}\right]^{}}{h(\tau )}}+{\displaystyle \frac{\left[K_{\phi \phi }\right]^{}}{l(\tau )}}\right\},`$ (21)
$`p_z`$ $`=`$ $`{\displaystyle \frac{1}{\kappa }}\left\{\left[K_{\tau \tau }\right]^{}{\displaystyle \frac{\left[K_{\phi \phi }\right]^{}}{l(\tau )}}\right\},`$ (22)
$`p_\phi `$ $`=`$ $`{\displaystyle \frac{1}{\kappa }}\left\{\left[K_{\tau \tau }\right]^{}{\displaystyle \frac{\left[K_{zz}\right]^{}}{h(\tau )}}\right\},`$ (23)
and $`w_a,z_a`$ and $`\phi _a`$ are unit vectors, defined as,
$$w_a=\delta _a^\tau ,z_a=h^{1/2}(\tau )\delta _a^z,\phi _a=l^{1/2}(\tau )\delta _a^\phi .$$
(24)
Clearly, the surface energy-momentum tensor given by Eq.(20) can be interpreted as representing a massive thin shell with its velocity $`w_a`$, and principal pressures $`p_z`$ and $`p_\phi `$, respectively, in the direction, $`z_a`$ and $`\phi _a`$, provided that it satisfies some energy conditions .
Using Eq.(6) and Eqs.(20) - (24), one can show that the conservation law on the hypersurface $`\mathrm{\Sigma }`$ ,
$$\tau _{a|b}^b=\left[T_{\alpha \beta }^+n^{+\alpha }e_{(a)}^{+\beta }T_{\alpha \beta }^{}n^\alpha e_{(a)}^\beta \right],$$
(25)
has only one non-vanishing component, which can be written as
$$\frac{d\rho }{d\tau }+\frac{(\rho +p_z)}{2h(\tau )}\frac{dh(\tau )}{d\tau }+\frac{(\rho +p_\phi )}{2l(\tau )}\frac{dl(\tau )}{d\tau }=\left[T_{\alpha \beta }^+n^{+\alpha }e_{(\tau )}^{+\beta }T_{\alpha \beta }^{}n^\alpha e_{(\tau )}^\beta \right],$$
(26)
where “$`|_b`$” denotes the covariant differentiation with respect to the three-metric $`\gamma _{ab}`$, and $`T_{\alpha \beta }^\pm `$ are the energy-momentum tensors calculated, respectively, in $`V^+`$ and $`V^{}`$, and
$`e_{(\tau )}^{+\mu }`$ $``$ $`{\displaystyle \frac{x^{+\mu }}{\tau }}=\left(f^+g^+R_{}^{}{}_{0}{}^{2}(T)\right)^{1/2}\left\{\delta _T^\mu +R_{}^{}{}_{0}{}^{2}(T)\delta _R^\mu \right\},`$ (27)
$`e_{(z)}^{+\mu }`$ $``$ $`{\displaystyle \frac{x^{+\mu }}{z}}=\delta _z^\mu ,e_{(\phi )}^{+\mu }{\displaystyle \frac{x^{+\mu }}{\phi }}=\delta _\phi ^\mu ,`$ (28)
$`e_{(\tau )}^\mu `$ $``$ $`{\displaystyle \frac{x^\mu }{\tau }}=\left(f^{}g^{}r_{}^{}{}_{0}{}^{2}(t)\right)^{1/2}\left\{\delta _t^\mu +r_{}^{}{}_{0}{}^{2}(t)\delta _r^\mu \right\},`$ (29)
$`e_{(z)}^\mu `$ $``$ $`{\displaystyle \frac{x^\mu }{z}}=\delta _z^\mu ,e_{(\phi )}^\mu {\displaystyle \frac{x^\mu }{\phi }}=\delta _\phi ^\mu .`$ (30)
When no matter shell appears on the hypersurface $`\mathrm{\Sigma }`$, we have $`\tau _{ab}=0`$, and the hypersurface represents a boundary surface , with the junction conditions being given by Eq.(7) and Eq.(26). The latter can be written in the form
$$T_{\alpha \beta }^+n^{+\alpha }e_{(\tau )}^{+\beta }|_\mathrm{\Sigma }=T_{\alpha \beta }^{}n^\alpha e_{(\tau )}^\beta |_\mathrm{\Sigma },(\tau _{ab}=0).$$
(31)
Once we have the general formulas, let us turn to consider their applications to some specific cases.
## III Gravitational Collapse of Cylindrical Shells Made of Counter-Rotating Dust particles
In this section, we shall consider the gravitational collapse of a cylindrical shells made of counter-rotating dust particles. The shell emits gravitational and particle radiations, when it is collapsing. The metric inside the shell will be chosen as that of Minkowski,
$$ds_{}^2=dt^2dr^2dz^2r^2d\phi ^2,$$
(32)
so that the symmetry axis is well defined and the local-flatness condition is satisfied . The metric outside the shell will be chosen as that representing out-going radiation fluid, given by
$$ds_+^2=e^{b(\xi )}\left(dT^2dR^2\right)dz^2R^2d\phi ^2,$$
(33)
where $`b(\xi )`$ is an arbitrary function of $`\xi `$ with $`\xi TR`$. Corresponding to the metric (33), the energy-momentum tensor is given by
$$T_{\mu \nu }^+=\frac{b^{}(\xi )}{R}k_\mu k_\nu ,$$
(34)
where $`k_\mu `$ is a null vector, defined as
$$k_\mu =\frac{1}{\sqrt{2}}\left(\delta _\mu ^T\delta _\mu ^R\right),$$
(35)
which is the generator of the out-going radial null geodesic congruence . The presence of the out-going gravitational waves is indicated by the only non-vanishing component of the Weyl tensor, $`C_{\mu \nu \lambda \sigma }`$, given by ,
$$\mathrm{\Psi }_0C_{\mu \nu \lambda \sigma }L^\mu M^\nu L^\lambda M^\sigma =\frac{b^{}(\xi )}{2R}e^{b(\xi )},$$
(36)
where $`L^\mu `$ and $`M^\mu `$ are null vectors, the definitions of which are given by Eq.(8) in .
From Eqs.(32) and (33) we find that the first junction conditions (7) now reduce to
$$d\tau =\left[1r_{}^{}{}_{0}{}^{2}(t)\right]^{1/2}dt=e^{b(\xi _0)/2}\left[1R_{}^{}{}_{0}{}^{2}(T)\right]^{1/2}dT,r_0(t)=R_0(T),$$
(37)
where $`\xi _0`$ is defined as $`\xi _0=TR_0(T)`$. From the above expressions we find
$$\left(\frac{dT}{dt}\right)^2=\frac{1}{\mathrm{\Delta }}\left[R_{}^{}{}_{0}{}^{2}(T)+e^{b(\xi _0)}\left(1R_{}^{}{}_{0}{}^{2}(T)\right)\right]^1,$$
(38)
which results in
$`{\displaystyle \frac{d^2T}{dt^2}}`$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{\Delta }^2}}\{2R_{}^{}{}_{0}{}^{}R_{}^{\prime \prime }{}_{0}{}^{}`$ (40)
$`e^{b(\xi _0)}[b^{}(\xi _0)(1R_{}^{}{}_{0}{}^{})(1R_{}^{}{}_{0}{}^{2})+2R_{}^{}{}_{0}{}^{}R_{}^{\prime \prime }{}_{0}{}^{}]\},`$
$`r_{}^{\prime \prime }{}_{0}{}^{}(t)`$ $`=`$ $`{\displaystyle \frac{d^2T}{dt^2}}R_{}^{}{}_{0}{}^{}+\left({\displaystyle \frac{dT}{dt}}\right)^2R_{}^{\prime \prime }{}_{0}{}^{}`$ (41)
$`=`$ $`{\displaystyle \frac{e^{b(\xi _0)}}{2\mathrm{\Delta }^2}}\left\{2R_{}^{\prime \prime }{}_{0}{}^{}+b^{}(\xi _0)R_{}^{}{}_{0}{}^{}(1R_{}^{}{}_{0}{}^{})(1R_{}^{}{}_{0}{}^{2})\right\}.`$ (42)
Inserting Eqs.(14) and the corresponding expressions for $`K_{ab}^{}`$ into Eq.(21), and considering Eq.(40), we find
$`\rho `$ $`=`$ $`{\displaystyle \frac{e^{b(\xi _0)/2}}{\kappa R_0(1R_{}^{}{}_{0}{}^{2})^{1/2}}}(\mathrm{\Delta }1),`$ (43)
$`p_z`$ $`=`$ $`{\displaystyle \frac{e^{b(\xi _0)/2}}{\kappa \mathrm{\Delta }R_{}^{}{}_{0}{}^{}(1R_{}^{}{}_{0}{}^{2})^{3/2}}}\{\mathrm{\Delta }(1\mathrm{\Delta })(1R_{}^{}{}_{0}{}^{2})(1\mathrm{\Delta })R_{}^{}{}_{0}{}^{}R_{}^{\prime \prime }{}_{0}{}^{}`$ (45)
$`{\displaystyle \frac{1}{2}}b^{}(\xi _0)R_{}^{}{}_{0}{}^{}(R_{}^{}{}_{0}{}^{}\mathrm{\Delta })(1R_{}^{}{}_{0}{}^{})(1R_{}^{}{}_{0}{}^{2})\},`$
$`p_\phi `$ $`=`$ $`{\displaystyle \frac{e^{b(\xi _0)/2}}{\kappa \mathrm{\Delta }(1R_{}^{}{}_{0}{}^{2})^{3/2}}}\{(\mathrm{\Delta }1)R_{}^{\prime \prime }{}_{0}{}^{}`$ (47)
$`{\displaystyle \frac{1}{2}}b^{}(\xi _0)(R_{}^{}{}_{0}{}^{}\mathrm{\Delta })(1R_{}^{}{}_{0}{}^{})(1R_{}^{}{}_{0}{}^{2})\}.`$
When the cylindrical thin shell is made of counter-rotating dust particles, where half of the dust particles orbit around the symmetry axis in a right-handed direction with angular momentum per unit rest mass $`p`$, and the other half orbit in the opposite, left-handed direction with angular momentum per unit rest mass $`p`$, the surface energy-momentum tensor is given by Eq.(20) with $`p_z=0`$ . Thus, setting $`p_z=0`$ in Eq.(43), we find
$$R_{}^{\prime \prime }{}_{0}{}^{}=\frac{1R_{}^{}{}_{0}{}^{2}}{R_{}^{}{}_{0}{}^{}}\left\{\mathrm{\Delta }+\frac{1}{2}b^{}(\xi _0)R_{}^{}{}_{0}{}^{}(1R_{}^{}{}_{0}{}^{})\frac{R_{}^{}{}_{0}{}^{}\mathrm{\Delta }}{\mathrm{\Delta }1}\right\},(p_z=0).$$
(48)
Substituting the above expression into Eq.(43), we obtain
$$\rho =p_\phi =\frac{e^{b(\xi _0)/2}}{\kappa R_0(1R_{}^{}{}_{0}{}^{2})^{1/2}}(\mathrm{\Delta }1),(p_z=0).$$
(49)
To further study the dynamics of the thin shell with $`p_z=0`$, we need to solve Eq.(48), which is found too difficult to be done in the general case. Therefore, in the following we shall consider a particular case in which
$`R_{}^{\prime \prime }{}_{0}{}^{}`$ $`=`$ $`{\displaystyle \frac{1R_{}^{}{}_{0}{}^{2}}{R_{}^{}{}_{0}{}^{}}}\beta ,`$ (50)
$`\beta `$ $`=`$ $`\mathrm{\Delta }+{\displaystyle \frac{1}{2}}b^{}(\xi _0)R_{}^{}{}_{0}{}^{}(1R_{}^{}{}_{0}{}^{}){\displaystyle \frac{R_{}^{}{}_{0}{}^{}\mathrm{\Delta }}{\mathrm{\Delta }1}},`$ (51)
where $`\beta `$ is an arbitrary constant. Once $`R_{}^{}{}_{0}{}^{}(T)`$ is known, the function $`b^{}(\xi )`$ can be obtained from Eq.(51) by quadrature. Since we are mainly interested in the dynamics of the shell, in the following we shall concentrate ourselves on Eq.(50). Integrating it we find that
$$R_{}^{}{}_{0}{}^{}(T)=\pm \left(1e^{2\beta T}\right)^{1/2},$$
(52)
where the “+” sign corresponds to an expanding shell, while the “$``$” sign corresponds to a contracting shell. In the following let us consider the two cases separately.
### A Expanding Thin Shells
When the “+” sign in Eq.(52) is chosen, the integration of it yields,
$$R_0(T)=(T+R_{min})\frac{1}{\beta }\left\{\left(1e^{2\beta T}\right)^{1/2}\mathrm{ln}\left[1+\left(1e^{2\beta T}\right)^{1/2}\right]\right\},$$
(53)
where $`R_{min}`$ is an integration constant. When $`\beta >0`$, we find that
$`R_0(T)`$ $`=`$ $`\{\begin{array}{cc}R_{min},\hfill & T=0\text{,}\hfill \\ +\mathrm{},\hfill & T+\mathrm{}\text{,}\hfill \end{array}`$ (54)
$`R_{}^{}{}_{0}{}^{}(T)`$ $`=`$ $`\{\begin{array}{cc}0,\hfill & T=0\text{,}\hfill \\ 1,\hfill & T+\mathrm{}\text{,}\hfill \end{array}(\beta >0),`$ (55)
which shows that in this case the corresponding solution represents the expansion of a thin shell made of counter-rotating dust particles. The expansion starts from the moment $`T=0`$ with the radius of the shell given by $`R(0)=R_{min}`$. At this moment the shell has zero radial velocity but infinitely large acceleration pointing outwards, as one can see from Eq.(50). Thus, the shell will expand until $`T=+\mathrm{}`$, where it reaches to its maximal radius $`R_0(+\mathrm{})=+\mathrm{}`$, with its radial velocity $`R_{}^{}{}_{0}{}^{}(+\mathrm{})=+1`$ and acceleration $`R_{}^{\prime \prime }{}_{0}{}^{}(+\mathrm{})=0`$.
When $`\beta <0`$, we find
$`R_0(T)`$ $`=`$ $`\{\begin{array}{cc}R_{min},\hfill & T=0\text{,}\hfill \\ 0,\hfill & T=|T_1|\text{,}\hfill \end{array}`$ (56)
$`R_{}^{}{}_{0}{}^{}(T)`$ $`=`$ $`\{\begin{array}{cc}0,\hfill & T=0\text{,}\hfill \\ \text{finite},\hfill & T=|T_1|\text{,}\hfill \end{array}(\beta <0).`$ (57)
Thus, now the solution represents the expansion of a thin shell, starting from zero radius at the moment $`T=|T_1|`$. It expands until the moment $`T=0`$, where its radial velocity and acceleration are given, respectively, by $`R_{}^{}{}_{0}{}^{}(T=0)=0`$ and $`R_{}^{\prime \prime }{}_{0}{}^{}(T=0)=\mathrm{}`$. Because of its huge acceleration that now points inwards, the shell will collapse from this moment on, by following a process similar to that to be described below.
### B Collapsing Thin Shells
When the “$``$” sign in Eq.(52) is chosen, we find that
$$R_0(T)=(R_{min}T)+\frac{1}{\beta }\left\{\left(1e^{2\beta T}\right)^{1/2}\mathrm{ln}\left[1+\left(1e^{2\beta T}\right)^{1/2}\right]\right\},$$
(58)
where $`R_{min}`$ is another integration constant. Thus, when $`\beta >0`$, from Eq.(58) we find that
$`R_0(T)`$ $`=`$ $`\{\begin{array}{cc}R_{min},\hfill & T=0\text{,}\hfill \\ 0,\hfill & T=|T_1|\text{,}\hfill \end{array}`$ (59)
$`R_{}^{}{}_{0}{}^{}(T)`$ $`=`$ $`\{\begin{array}{cc}0,\hfill & T=0\text{,}\hfill \\ \left(1e^{2\beta |T_1|}\right)^2,\hfill & T=|T_1|\text{,}\hfill \end{array}(\beta >0),`$ (60)
which shows that now the shell starts to collapse at the moment $`T=0`$ with zero radial velocity. The collapse ends up at the moment $`T=|T_1|`$, where the whole shell contracts into a line-like spacetime singularity, as Eqs.(34) and (49) show. Therefore, unlike the case studied by AT , in the present case the centrifugal forces of the courter-rotating dust particles are not strong enough to prevent the collapse from forming a spacetime singularity.
When $`\beta <0`$, from Eqs.(52) and (58), we find
$`R_0(T)`$ $`=`$ $`\{\begin{array}{cc}+\mathrm{}\hfill & T\mathrm{}\text{,}\hfill \\ R_{min},\hfill & T=0\text{,}\hfill \end{array}`$ (61)
$`R_{}^{}{}_{0}{}^{}(T)`$ $`=`$ $`\{\begin{array}{cc}1,\hfill & T\mathrm{}\text{,}\hfill \\ 0,\hfill & T=0\text{,}\hfill \end{array}(\beta <0).`$ (62)
Thus, in the present case the shell starts to collapse at the moment $`T=\mathrm{}`$ with its radius $`R_0(\mathrm{})=+\mathrm{}`$ and its radial velocity $`R_{}^{}{}_{0}{}^{}(\mathrm{})=1`$. As it collapses, it is radiating massless particles and gravitational waves, as one can see from Eqs.(34) and (36). at the moment $`T=0`$, it collapses to its minimal radius $`R_0(0)=R_{min}`$, where its velocity becomes zero. As far as $`R_{min}0`$, in this case no spacetime singularity is formed, and the centrifugal forces of the dust particles now are strong enough to halt the collapse. On the other hand, from Eq.(50) we can see that at T = 0 the acceleration of the shell becomes infinitely large and points outwards. So, from this moment on, the shell will expand, by following a process similar to that described in the last subsection. When $`R_{min}=0`$, the centrifugal forces are still not strong enough to prevent the shell from collapsing into a zero radius, whereby a spacetime singularity is formed.
## IV Conclusions
In this paper, the general formulas of a non-rotating dynamic thin shell that connects two arbitrary cylindrical regions have been given in terms of the metric coefficients and their first derivatives, using Israel’s method. As an application of these formulas, the dynamics of a thin shell made of counter-rotating non-interacting particles, which emits both gravitational waves and massless particles, has been studied. It has been found that in some cases the models represent an expanding shell and others a collapsing shell. For the collapsing shell, two possible final states exist. In one case, after the shell collapses to a minimal non-zero radius, it starts to expands, that is, the angular momentum of the dust particles is strong enough to halt the collapse, so that a spacetime singularity is prevented from forming on the symmetry axis. However, in the other case, the rotation is not strong enough to halt the collapse at a finite non-zero radius, and as a result a spacetime singularity is formed finally. These results are different from the ones obtained by AT in the radiation-free case , but similar to the ones obtained by ST for rotating spheroids with radiation .
## Acknowledgment
The financial assistance from CNPq and FAPERJ (AW) is gratefully acknowledged.
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# Mermin’s n-particle Bell inequality and operators’ noncommutativity
(Published on 13 August 2001)
## Abstract
The relationship between the noncommutativity of operators and the violation of the Bell inequality is exhibited in the light of the $`n`$-particle Bell-type inequality discovered by Mermin \[Phys. Rev. Lett. 65 (1990) 1838\]. It is shown, in particular, that the maximal amount of violation of Mermin’s inequality predicted by quantum mechanics decreases exponentially by a factor of $`2^{m/2}`$ whenever any $`m`$ among the $`n`$ single-particle commutators happen to vanish.
PACS numbers: 03.65.Bz
The quantum-mechanical violation of the Clauser-Horne-Shimony-Holt (CHSH) form of Bell’s inequality is a direct consequence of the fact that quantum operators obey a noncommutative algebra. For two spin-$`\frac{1}{2}`$ particles, the CHSH inequality can be written as $`|\widehat{}_{\text{CHSH}}|2`$, where $`\widehat{}_{\text{CHSH}}`$ denotes the expectation value of the Bell operator
$$\widehat{}_{\text{CHSH}}=\sigma (\widehat{\text{n}}_1)\sigma (\widehat{\text{n}}_2)+\sigma (\widehat{\text{n}}_1)\sigma (\widehat{\text{n}}_2^{})+\sigma (\widehat{\text{n}}_1^{})\sigma (\widehat{\text{n}}_2)\sigma (\widehat{\text{n}}_1^{})\sigma (\widehat{\text{n}}_2^{}).$$
(1)
In Eq. (1), $`\sigma (\widehat{\text{n}}_j)`$ ($`\sigma (\widehat{\text{n}}_j^{})`$) denotes the spin operator for particle $`j`$ along the direction $`\widehat{\text{n}}_j`$ ($`\widehat{\text{n}}_j^{}`$). From the fact that $`\sigma ^2(\widehat{\text{n}}_1)=\sigma ^2(\widehat{\text{n}}_1^{})=\sigma ^2(\widehat{\text{n}}_2)=\sigma ^2(\widehat{\text{n}}_2^{})=\widehat{I}`$, it follows immediately that
$$\widehat{}_{\text{CHSH}}^2=4\widehat{I}[\sigma (\widehat{\text{n}}_1),\sigma (\widehat{\text{n}}_1^{})][\sigma (\widehat{\text{n}}_2),\sigma (\widehat{\text{n}}_2^{})].$$
(2)
From Eq. (2), we can see that the CHSH inequality may be violated by the quantum-mechanical predictions provided that $`[\sigma (\widehat{\text{n}}_1),\sigma (\widehat{\text{n}}_1^{})]0`$ and $`[\sigma (\widehat{\text{n}}_2),\sigma (\widehat{\text{n}}_2^{})]0`$. Whenever either commutator $`[\sigma (\widehat{\text{n}}_1),\sigma (\widehat{\text{n}}_1^{})]`$ or $`[\sigma (\widehat{\text{n}}_2),\sigma (\widehat{\text{n}}_2^{})]`$ vanishes then $`\widehat{}_{\text{CHSH}}=\pm 2\widehat{I}`$, and so $`|\widehat{}_{\text{CHSH}}|=2`$ for any joint state describing the spin of the particles. On the other hand, for three spin-$`\frac{1}{2}`$ particles the appropriate Bell inequality is of the form $`|\widehat{}_\text{H}|2`$ (see Eq. (14) of ), where now the representative Bell operator can be written as
$$\widehat{}_\text{H}=\sigma (\widehat{\text{n}}_1^{})\sigma (\widehat{\text{n}}_2)\sigma (\widehat{\text{n}}_3)+\sigma (\widehat{\text{n}}_1)\sigma (\widehat{\text{n}}_2^{})\sigma (\widehat{\text{n}}_3)+\sigma (\widehat{\text{n}}_1)\sigma (\widehat{\text{n}}_2)\sigma (\widehat{\text{n}}_3^{})\sigma (\widehat{\text{n}}_1^{})\sigma (\widehat{\text{n}}_2^{})\sigma (\widehat{\text{n}}_3^{}).$$
(3)
It is simple algebra to verify that
$$\begin{array}{c}\widehat{}_\text{H}^2=4\widehat{I}[\sigma (\widehat{\text{n}}_1),\sigma (\widehat{\text{n}}_1^{})][\sigma (\widehat{\text{n}}_2),\sigma (\widehat{\text{n}}_2^{})]\hfill \\ \hfill [\sigma (\widehat{\text{n}}_2),\sigma (\widehat{\text{n}}_2^{})][\sigma (\widehat{\text{n}}_3),\sigma (\widehat{\text{n}}_3^{})][\sigma (\widehat{\text{n}}_1),\sigma (\widehat{\text{n}}_1^{})][\sigma (\widehat{\text{n}}_3),\sigma (\widehat{\text{n}}_3^{})].\end{array}$$
(4)
From Eq. (4), we can see that in order for the inequality $`|\widehat{}_\text{H}|2`$ to be violated by the quantum predictions it is necessary that at least two of the commutators involved be nonzero. When one of the commutators, say $`[\sigma (\widehat{\text{n}}_3),\sigma (\widehat{\text{n}}_3^{})]`$, vanishes then the expression in Eq. (4) reduces to $`\widehat{}_\text{H}^2=4\widehat{I}[\sigma (\widehat{\text{n}}_1),\sigma (\widehat{\text{n}}_1^{})][\sigma (\widehat{\text{n}}_2),\sigma (\widehat{\text{n}}_2^{})]`$ which corresponds to the square of the two-particle Bell operator $`\widehat{}_{\text{CHSH}}`$, Eq. (2). This means that whenever the commutator $`[\sigma (\widehat{\text{n}}_j),\sigma (\widehat{\text{n}}_j^{})]`$ associated with any one of the particles $`j`$, $`j=1,2,3`$, is equal to zero then this particle plays no role in the violation of the corresponding three-particle Bell inequality, so that this Bell inequality of order 3 reduces to a Bell inequality of order 2, where by order of a Bell inequality we want to mean the number of “effective” particles involved in the associated Bell operator. The maximum eigenvalue of the operator $`\widehat{}_{\text{CHSH}}`$ is (in terms of the absolute value) $`2\sqrt{2}`$, whereas the maximum eigenvalue of $`\widehat{}_\text{H}`$ is 4. Consequently, the maximal quantum-mechanical violation of the inequality $`|\widehat{}_\text{H}|2`$ decreases by a factor of $`2\sqrt{2}/4=1/\sqrt{2}`$ when any one of the three commutators vanishes. (Of course, for two vanishing commutators we have $`\widehat{}_\text{H}=\pm 2\widehat{I}`$, and then the inequality at issue is saturated by the quantum-mechanical predictions.)
In this Letter, I would like to extend these results further to the general case in which $`n`$ spin-$`\frac{1}{2}`$ particles are considered. We will show that, under certain specific conditions which will be precisely stated below, a Bell inequality of order $`n`$ reduces to a Bell inequality of order $`nm`$ whenever any $`m`$ of the $`n`$ single-particle commutators $`[\sigma (\widehat{\text{n}}_j),\sigma (\widehat{\text{n}}_j^{})]`$ vanish. This is done in a way that explicitly shows the recognised relation between the presence of commutators in the Bell operator and the peculiar nonlocal properties exhibited by quantum mechanics. For this purpose we consider the Bell-type inequality derived by Mermin in 1990 for a system of $`n`$ spin-$`\frac{1}{2}`$ particles ($`n3`$) in the Greenberger-Horne-Zeilinger state
$$|\mathrm{\Phi }=(1/\sqrt{2})(|\mathrm{}+i|\mathrm{}),$$
(5)
where $``$ or $``$ in the $`j`$th position corresponds to the component of the spin of the $`j`$th particle along its own $`z`$ axis. Mermin’s Bell inequality can be summarised as
$`|\widehat{}_\text{M}|`$ $`2^{n/2},n\text{ even},`$ (6a)
$`|\widehat{}_\text{M}|`$ $`2^{(n1)/2},n\text{ odd},`$ (6b)
where $`\widehat{}_\text{M}`$ denotes the expectation value of the (Hermitian) Bell operator
$$\widehat{}_\text{M}=\frac{1}{2i}\left(\underset{j=1}{\overset{n}{}}(\sigma _x^j+i\sigma _y^j)\underset{j=1}{\overset{n}{}}(\sigma _x^ji\sigma _y^j)\right).$$
(7)
Mermin showed that the prediction that quantum mechanics makes for state (5) violates the inequalities in Eqs. (6) by an exponentially large factor of $`2^{(n2)/2}`$ for $`n`$ even or $`2^{(n1)/2}`$ for $`n`$ odd. This was the first spectacular demonstration of the fact that there is no limit to the amount by which the quantum-mechanical correlations can exceed the limits imposed by a Bell inequality.<sup>1</sup><sup>1</sup>1A noteworthy alternative derivation of the locally realistic bounds in Eqs. (6), along with the bound prescribed by quantum mechanics, can be found in Ref. (see, in particular, the equations (3.15), (3.17), and (3.18) of ). In order to achieve a greater generality, we will consider the most general form of the Bell operator in Eq. (7), namely,
$$\widehat{}_\text{M}=\frac{1}{2i}\left(\underset{j=1}{\overset{n}{}}(\sigma (\widehat{\text{n}}_j)+i\sigma (\widehat{\text{n}}_j^{}))\underset{j=1}{\overset{n}{}}(\sigma (\widehat{\text{n}}_j)i\sigma (\widehat{\text{n}}_j^{}))\right),$$
(8)
where $`\widehat{\text{n}}_j`$ and $`\widehat{\text{n}}_j^{}`$ are arbitrary directions. Of course, the Mermin’s Bell inequalities (6) still hold for the general Bell operator (8). The treatment that follows does not rely on any particular state like that of Eq. (5). Rather, it is based on the general properties exhibited by the square of the Bell operator (8) when expressed in terms of the commutators $`[\sigma (\widehat{\text{n}}_j),\sigma (\widehat{\text{n}}_j^{})]`$. In order to abbreviate the notation, we shall henceforth drop the $`\sigma `$’s of the commutators so that they will be written simply as $`[\widehat{\text{n}}_j,\widehat{\text{n}}_j^{}]`$. Now we must distinguish between the cases of $`n`$ odd and $`n`$ even. For an odd number of particles ($`n3`$) the square of operator (8) is given by
$$\begin{array}{c}\widehat{}_\text{M}^2(n\text{ odd})=2^{n1}\widehat{I}2^{n3}\stackrel{\left(\begin{array}{c}n\\ 2\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}]}}\hfill \\ \hfill +2^{n5}\stackrel{\left(\begin{array}{c}n\\ 4\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2<m_3<m_4}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}][\widehat{\text{n}}_{m_3},\widehat{\text{n}}_{m_3}^{}][\widehat{\text{n}}_{m_4},\widehat{\text{n}}_{m_4}^{}]}}\\ \hfill +\mathrm{}+(1)^k\mathrm{\hspace{0.17em}2}^{n2k1}\stackrel{\left(\begin{array}{c}n\\ 2k\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2<\mathrm{}<m_{2k}}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}]\mathrm{}[\widehat{\text{n}}_{m_{2k}},\widehat{\text{n}}_{m_{2k}}^{}]}}\\ \hfill +\mathrm{}+(1)^{(n1)/2}\stackrel{\left(\begin{array}{c}n\\ n1\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2<\mathrm{}<m_{n1}}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}]\mathrm{}[\widehat{\text{n}}_{m_{n1}},\widehat{\text{n}}_{m_{n1}}^{}]}},\end{array}$$
(9)
where {$`m_1,m_2,\mathrm{},m_{n1}\}`$ is a set of $`n1`$ indices each of which running from 1 to $`n`$. Analogously, for $`n`$ even ($`n4`$) we have the slightly more complicated expression
$$\begin{array}{c}\widehat{}_\text{M}^2(n\text{ even})=2^{n1}\widehat{I}2^{n3}\stackrel{\left(\begin{array}{c}n\\ 2\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}]}}\hfill \\ \hfill +2^{n5}\stackrel{\left(\begin{array}{c}n\\ 4\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2<m_3<m_4}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}][\widehat{\text{n}}_{m_3},\widehat{\text{n}}_{m_3}^{}][\widehat{\text{n}}_{m_4},\widehat{\text{n}}_{m_4}^{}]}}\\ \hfill +\mathrm{}+(1)^k\mathrm{\hspace{0.17em}2}^{n2k1}\stackrel{\left(\begin{array}{c}n\\ 2k\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2<\mathrm{}<m_{2k}}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}]\mathrm{}[\widehat{\text{n}}_{m_{2k}},\widehat{\text{n}}_{m_{2k}}^{}]}}\\ \hfill +\mathrm{}+(1)^{(n2)/2}\mathrm{\hspace{0.17em}2}\stackrel{\left(\begin{array}{c}n\\ n2\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2<\mathrm{}<m_{n2}}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}]\mathrm{}[\widehat{\text{n}}_{m_{n2}},\widehat{\text{n}}_{m_{n2}}^{}]}}\\ \hfill +(1)^{n/2}\frac{1}{2}\left([\widehat{\text{n}}_1,\widehat{\text{n}}_1^{}][\widehat{\text{n}}_2,\widehat{\text{n}}_2^{}]\mathrm{}[\widehat{\text{n}}_n,\widehat{\text{n}}_n^{}]\{\widehat{\text{n}}_1,\widehat{\text{n}}_1^{}\}\{\widehat{\text{n}}_2,\widehat{\text{n}}_2^{}\}\mathrm{}\{\widehat{\text{n}}_n,\widehat{\text{n}}_n^{}\}\right),\end{array}$$
(10)
where $`\{m_1,m_2,\mathrm{},m_{n2}\}`$ is a set of $`n2`$ indices each of which running from 1 to $`n`$, and $`\{\widehat{\text{n}}_j,\widehat{\text{n}}_j^{}\}`$ is the anticommutator of the operators $`\sigma (\widehat{\text{n}}_j)`$ and $`\sigma (\widehat{\text{n}}_j^{})`$.
Let us first examine what happens when any two of the commutators in Eq. (9) are equal to zero. In the first place, it is clear that all last $`\left(\genfrac{}{}{0pt}{}{n}{n1}\right)`$ terms in Eq. (9) vanish when any two commutators are equal to zero, since each of these $`n`$ terms is a product of $`n1`$ distinct commutators. Furthermore, it can be seen that, for each group of $`\left(\genfrac{}{}{0pt}{}{n}{2k}\right)`$ terms in Eq. (9) (with each of the terms in a group being a product of $`2k`$ commutators, $`2kn3`$), there is a total of $`\left(\genfrac{}{}{0pt}{}{n2}{2k}\right)`$ terms which remain “untouched” when any two of the commutators are made to vanish, the remaining $`\left(\genfrac{}{}{0pt}{}{n}{2k}\right)\left(\genfrac{}{}{0pt}{}{n2}{2k}\right)`$ terms in the group being equal to zero due to the vanishing of the two commutators. Therefore, whenever any two commutators are equal to zero, expression (9) reduces to
$$\begin{array}{c}\widehat{}_\text{M}^2(n\text{ odd})2^2(2^{n3}\widehat{I}2^{n5}\stackrel{\left(\begin{array}{c}n2\\ 2\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}]}}\hfill \\ \hfill +2^{n7}\stackrel{\left(\begin{array}{c}n2\\ 4\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2<m_3<m_4}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}][\widehat{\text{n}}_{m_3},\widehat{\text{n}}_{m_3}^{}][\widehat{\text{n}}_{m_4},\widehat{\text{n}}_{m_4}^{}]}}\\ \hfill +\mathrm{}+(1)^k\mathrm{\hspace{0.17em}2}^{n2k3}\stackrel{\left(\begin{array}{c}n2\\ 2k\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2<\mathrm{}<m_{2k}}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}]\mathrm{}[\widehat{\text{n}}_{m_{2k}},\widehat{\text{n}}_{m_{2k}}^{}]}}\\ \hfill +\mathrm{}+(1)^{(n3)/2}\stackrel{\left(\begin{array}{c}n2\\ n3\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2<\mathrm{}<m_{n3}}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}]\mathrm{}[\widehat{\text{n}}_{m_{n3}},\widehat{\text{n}}_{m_{n3}}^{}]}}).\end{array}$$
(11)
We may enumerate the $`n2`$ particles corresponding to the $`n2`$ nonvanishing commutators, so that each of the $`n3`$ indices $`\{m_1,m_2,\mathrm{},m_{n3}\}`$ in Eq. (11) runs from 1 to $`n2`$. We introduce the notation $`\widehat{}_\text{M}^2(n|m)`$ to denote the squared $`n`$-particle Bell operator that obtains when $`m`$ single-particle commutators vanish. Then we can put the whole expression (11) in a compact notation as $`\widehat{}_\text{M}^2(n\text{ odd}|2)=2^2\widehat{}_\text{M}^2(n2)`$. Suppose now that any two commutators appearing in Eq. (11) happen in turn to vanish. Then it is clear that an analysis similar to the one we have just performed for Eq. (9) would enable us to deduce that $`\widehat{}_\text{M}^2(n\text{ odd}|4)=2^4\widehat{}_\text{M}^2(n4)`$. Iterating this procedure successively would lead us to finally conclude that
$$\widehat{}_\text{M}^2(n\text{ odd}|2k)=2^{2k}\widehat{}_\text{M}^2(n2k),2k=0,2,4,\mathrm{},n3.$$
(12)
Likewise, we may determine the expression which results when two of the commutators in Eq. (10) are equal to zero. To do this, we need the auxiliary result according to which, for a spin-$`\frac{1}{2}`$ particle, if $`[\widehat{\text{n}}_j,\widehat{\text{n}}_j^{}]=0`$ then necessarily the anticommutator $`\{\widehat{\text{n}}_j,\widehat{\text{n}}_j^{}\}`$ is either $`2\widehat{I}_j`$ or $`2\widehat{I}_j`$, with $`\widehat{I}_j`$ being the identity operator acting on the Hilbert space pertaining to particle $`j`$. Without any loss of generality, we shall take the axes $`\widehat{\text{n}}_i`$, $`\widehat{\text{n}}_i^{}`$, $`\widehat{\text{n}}_j`$, and $`\widehat{\text{n}}_j^{}`$, $`ij`$, such that whenever we have $`[\widehat{\text{n}}_i,\widehat{\text{n}}_i^{}]=[\widehat{\text{n}}_j,\widehat{\text{n}}_j^{}]=0`$ then $`\{\widehat{\text{n}}_i,\widehat{\text{n}}_i^{}\}\{\widehat{\text{n}}_j,\widehat{\text{n}}_j^{}\}=4\widehat{I}_i\widehat{I}_j`$. Then it can be shown that, for any two vanishing commutators, the expression in Eq. (10) reduces to
$$\begin{array}{c}\widehat{}_\text{M}^2(n\text{ even})2^2(2^{n3}\widehat{I}2^{n5}\stackrel{\left(\begin{array}{c}n2\\ 2\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}]}}\hfill \\ \hfill +2^{n7}\stackrel{\left(\begin{array}{c}n2\\ 4\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2<m_3<m_4}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}][\widehat{\text{n}}_{m_3},\widehat{\text{n}}_{m_3}^{}][\widehat{\text{n}}_{m_4},\widehat{\text{n}}_{m_4}^{}]}}\\ \hfill +\mathrm{}+(1)^k\mathrm{\hspace{0.17em}2}^{n2k3}\stackrel{\left(\begin{array}{c}n2\\ 2k\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2<\mathrm{}<m_{2k}}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}]\mathrm{}[\widehat{\text{n}}_{m_{2k}},\widehat{\text{n}}_{m_{2k}}^{}]}}\\ \hfill +\mathrm{}+(1)^{(n4)/2}\mathrm{\hspace{0.17em}2}\stackrel{\left(\begin{array}{c}n2\\ n4\end{array}\right)\text{ terms}}{\stackrel{}{\underset{m_1<m_2<\mathrm{}<m_{n4}}{}[\widehat{\text{n}}_{m_1},\widehat{\text{n}}_{m_1}^{}][\widehat{\text{n}}_{m_2},\widehat{\text{n}}_{m_2}^{}]\mathrm{}[\widehat{\text{n}}_{m_{n4}},\widehat{\text{n}}_{m_{n4}}^{}]}}\\ \hfill +(1)^{(n2)/2}\frac{1}{2}([\widehat{\text{n}}_1,\widehat{\text{n}}_1^{}][\widehat{\text{n}}_2,\widehat{\text{n}}_2^{}]\mathrm{}[\widehat{\text{n}}_{n2},\widehat{\text{n}}_{n2}^{}]\{\widehat{\text{n}}_1,\widehat{\text{n}}_1^{}\}\{\widehat{\text{n}}_2,\widehat{\text{n}}_2^{}\}\mathrm{}\{\widehat{\text{n}}_{n2},\widehat{\text{n}}_{n2}^{}\})),\end{array}$$
(13)
where the particles corresponding to the $`n2`$ nonvanishing commutators have been enumerated so that each of the $`n4`$ indices $`\{m_1,m_2,\mathrm{},m_{n4}\}`$ in Eq. (13) runs from 1 to $`n2`$. Eq. (13) can be expressed as $`\widehat{}_\text{M}^2(n\text{ even}|2)=2^2\widehat{}_\text{M}^2(n2)`$. If an even number $`2k`$ of commutators are equal to zero then, by applying successively the reasoning leading to Eq. (13), we could equally conclude that (cf. Eq. (12))
$$\widehat{}_\text{M}^2(n\text{ even}|2k)=2^{2k}\widehat{}_\text{M}^2(n2k),2k=0,2,4,\mathrm{},n4.$$
(14)
A natural question that arises is whether the squared Bell operator $`\widehat{}_\text{M}^2(n)`$ does “collapse” into $`\widehat{}_\text{M}^2(n1)`$ when one of the commutators is equal to zero. A glance at expressions (9) and (10) tells us that $`\widehat{}_\text{M}^2(n)`$ cannot, in general, reduce to $`\widehat{}_\text{M}^2(n1)`$ because of the presence of the product of anticommutators in the last term of Eq. (10). If this product were equal to zero then there would be a “continuous” transition between $`\widehat{}_\text{M}^2(n)`$ and $`\widehat{}_\text{M}^2(n1)`$ when one commutator vanishes or, more generally, between $`\widehat{}_\text{M}^2(n)`$ and $`\widehat{}_\text{M}^2(nm)`$ when $`m`$ commutators vanish. Obviously, in order for the product $`\{\widehat{\text{n}}_1,\widehat{\text{n}}_1^{}\}\{\widehat{\text{n}}_2,\widehat{\text{n}}_2^{}\}\mathrm{}\{\widehat{\text{n}}_n,\widehat{\text{n}}_n^{}\}`$ to be equal to zero, it suffices that any given one of the factors vanishes. So we shall assume that one of the anticommutators, $`\{\widehat{\text{n}}_i,\widehat{\text{n}}_i^{}\}`$ say, is zero, which amounts to making the directions $`\widehat{\text{n}}_i`$ and $`\widehat{\text{n}}_i^{}`$ perpendicular between themselves. Regarding the quantum-mechanical violation of the Mermin’s Bell inequalities (6), the restriction $`\{\widehat{\text{n}}_i,\widehat{\text{n}}_i^{}\}=0`$ does not entail any real limitation. In fact, as we shall see below, in order to achieve the maximum quantum violation of Mermin’s inequalities, it is necessary that $`\widehat{\text{n}}_j`$ be perpendicular to $`\widehat{\text{n}}_j^{}`$ for all $`j=1,2,\mathrm{},n`$. At any event, for the case in which one of the anticommutators vanishes, an analysis similar to that used to derive the Eqs. (12) and (14) would enable us to conclude that
$$\widehat{}_\text{M}^2(n\text{ odd}|2k+1)=2^{2k+1}\widehat{}_\text{M}^2(n2k1),2k+1=1,3,\mathrm{},n4,$$
(15)
and
$$\widehat{}_\text{M}^2(n\text{ even}|2k+1)=2^{2k+1}\widehat{}_\text{M}^2(n2k1),2k+1=1,3,\mathrm{},n3.$$
(16)
We stress that the condition $`\{\widehat{\text{n}}_i,\widehat{\text{n}}_i^{}\}=0`$ upon which relations (15) and (16) are based, restricts in no way their “practical” validity since, as can be seen from Eqs. (9) and (10), it is absolutely necessary that at least two of the commutators be nonzero, if we want that the quantum-mechanical predictions can violate inequalities (6). Hence, as the condition $`\{\widehat{\text{n}}_i,\widehat{\text{n}}_i^{}\}=0`$ implies that $`[\widehat{\text{n}}_i,\widehat{\text{n}}_i^{}]0`$, we can always choose one of the nonvanishing commutators to be $`[\widehat{\text{n}}_i,\widehat{\text{n}}_i^{}]`$.
Eqs. (12) and (14)-(16) can be unified in the single relation
$$\widehat{}_\text{M}^2(n|m)=2^m\widehat{}_\text{M}^2(nm),m=0,1,2,\mathrm{},n3,$$
(17)
which applies for any $`n`$ ($`n3`$).<sup>2</sup><sup>2</sup>2The case $`m=n2`$ is excluded because the Mermin’s Bell operator $`\widehat{}_\text{M}(2)=\sigma (\widehat{\text{n}}_1)\sigma (\widehat{\text{n}}_2^{})+\sigma (\widehat{\text{n}}_1^{})\sigma (\widehat{\text{n}}_2)`$ does not entail any meaningful Bell-type inequality. The transition from the three-particle Bell operator $`\widehat{}_\text{H}`$ (which can be identified with $`\widehat{}_\text{M}(3)`$) to the two-particle Bell operator $`\widehat{}_{\text{CHSH}}`$ has been treated separately at the opening part of this Letter. So, as the operator $`\widehat{}_\text{M}^2(n|m)`$ is proportional to $`\widehat{}_\text{M}^2(nm)`$, we have shown that a Bell inequality of order $`n`$ reduces to a Bell inequality of order $`nm`$ whenever any $`m`$ of the $`n`$ single-particle commutators $`[\widehat{\text{n}}_j,\widehat{\text{n}}_j^{}]`$ are equal to zero. The maximum eigenvalue of the Bell operator $`\widehat{}_\text{M}(n)`$ is $`2^{n1}`$. From Eq. (17), we can see that the maximum eigenvalue of the operator $`\widehat{}_\text{M}(n|m)`$ is $`2^{m/2}`$ times the maximum eigenvalue of $`\widehat{}_\text{M}(nm)`$, namely $`2^{m/2}2^{n1}`$. It thus follows that for the case in which $`m`$ single-particle commutators vanish, the maximum amount of violation predicted by quantum mechanics of the $`n`$-particle Bell-type inequalities (6) diminishes by a factor of $`2^{m/2}`$ with respect to the maximal violation obtained when all $`n`$ anticommutators vanish. In particular, for one vanishing commutator, the decrease factor is $`1/\sqrt{2}`$.
We conclude by noting an important feature of the general expressions in Eqs. (9) and (10), namely, that all products in such expressions involve an even number of commutators (except the last product in Eq. (10) which involves an even number of anticommutators). This fact implies that all the eigenvectors of the squared Bell operator $`\widehat{}_\text{M}^2(n)`$ are two-fold degenerate. To see this, let us consider the case of $`n`$ even, the treatment and conclusions for the case of $`n`$ odd being the same as for $`n`$ even. Without loss of generality we may take the axes $`\widehat{\text{n}}_j`$ and $`\widehat{\text{n}}_j^{}`$ as defining the $`x`$$`y`$ plane associated with particle $`j`$, so that such directions $`\widehat{\text{n}}_j`$ and $`\widehat{\text{n}}_j^{}`$ are specified by the azimuthal angles $`\varphi _j`$ and $`\varphi _j^{}`$, respectively. Thus, replacing in Eq. (10) each commutator $`[\widehat{\text{n}}_j,\widehat{\text{n}}_j^{}]`$ (anticommutator $`\{\widehat{\text{n}}_j,\widehat{\text{n}}_j^{}\}`$) by its value $`2i\mathrm{sin}\theta _j\widehat{\text{z}}_j`$ ($`2\mathrm{cos}\theta _j\widehat{I}_j`$), we obtain
$$\begin{array}{c}\widehat{}_\text{M}^2(n\text{ even})=2^{n1}(\widehat{I}+\underset{m_1<m_2}{}\mathrm{sin}\theta _{m_1}\mathrm{sin}\theta _{m_2}\widehat{\text{z}}_{m_1}\widehat{\text{z}}_{m_2}\hfill \\ \hfill +\underset{m_1<m_2<m_3<m_4}{}\mathrm{sin}\theta _{m_1}\mathrm{sin}\theta _{m_2}\mathrm{sin}\theta _{m_3}\mathrm{sin}\theta _{m_4}\widehat{\text{z}}_{m_1}\widehat{\text{z}}_{m_2}\widehat{\text{z}}_{m_3}\widehat{\text{z}}_{m_4}\\ \hfill +\mathrm{}+\underset{m_1<m_2<\mathrm{}<m_{2k}}{}\mathrm{sin}\theta _{m_1}\mathrm{sin}\theta _{m_2}\mathrm{}\mathrm{sin}\theta _{m_{2k}}\widehat{\text{z}}_{m_1}\widehat{\text{z}}_{m_2}\mathrm{}\widehat{\text{z}}_{m_{2k}}\\ \hfill +\mathrm{}+\underset{m_1<m_2<\mathrm{}<m_{n2}}{}\mathrm{sin}\theta _{m_1}\mathrm{sin}\theta _{m_2}\mathrm{}\mathrm{sin}\theta _{m_{n2}}\widehat{\text{z}}_{m_1}\widehat{\text{z}}_{m_2}\mathrm{}\widehat{\text{z}}_{m_{n2}}\\ \hfill +\mathrm{sin}\theta _1\mathrm{sin}\theta _2\mathrm{}\mathrm{sin}\theta _n\widehat{\text{z}}_1\widehat{\text{z}}_2\mathrm{}\widehat{\text{z}}_n(1)^{n/2}\mathrm{cos}\theta _1\mathrm{cos}\theta _2\mathrm{}\mathrm{cos}\theta _n\widehat{I}),\end{array}$$
(18)
where $`\theta _j=\varphi _j^{}\varphi _j`$ is the angle included between $`\widehat{\text{n}}_j`$ and $`\widehat{\text{n}}_j^{}`$, and $`\widehat{\text{z}}_j`$ is the spin operator for particle $`j`$ along its own $`z`$-axis. From Eq. (18), it is apparent that every one of the $`2^n`$ vectors $`|z_1,z_2,\mathrm{},z_n`$ is an eigenvector of $`\widehat{}_\text{M}^2(n\text{ even})`$, where $`|z_j`$ designates the eigenvector of $`\widehat{\text{z}}_j`$ with eigenvalue $`z_j=+1\text{ or }1`$. Also, from Eq. (18), it is clear that if $`|z_1,z_2,\mathrm{},z_n`$ is an eigenvector of $`\widehat{}_\text{M}^2(n\text{ even})`$ with associated eigenvalue $`\mu `$, then the same holds true for the eigenvector $`|z_1,z_2,\mathrm{},z_n`$. For the special case where $`\theta _j=\pi /2`$ for each $`j=1,2,\mathrm{},n`$, Eq. (18) becomes
$$\begin{array}{c}\widehat{}_\text{M}^2(n\text{ even})=2^{n1}(\widehat{I}+\underset{m_1<m_2}{}\widehat{\text{z}}_{m_1}\widehat{\text{z}}_{m_2}+\underset{m_1<m_2<m_3<m_4}{}\widehat{\text{z}}_{m_1}\widehat{\text{z}}_{m_2}\widehat{\text{z}}_{m_3}\widehat{\text{z}}_{m_4}\hfill \\ \hfill +\mathrm{}+\underset{m_1<m_2<\mathrm{}<m_{2k}}{}\widehat{\text{z}}_{m_1}\widehat{\text{z}}_{m_2}\mathrm{}\widehat{\text{z}}_{m_{2k}}\\ \hfill +\mathrm{}+\underset{m_1<m_2<\mathrm{}<m_{n2}}{}\widehat{\text{z}}_{m_1}\widehat{\text{z}}_{m_2}\mathrm{}\widehat{\text{z}}_{m_{n2}}+\widehat{\text{z}}_1\widehat{\text{z}}_2\mathrm{}\widehat{\text{z}}_n).\end{array}$$
(19)
Since the total number of terms in (19) is $`_{k=0}^{n/2}\left(\genfrac{}{}{0pt}{}{n}{2k}\right)=2^{n1}`$, we can immediately conclude that, for the considered case in which $`\theta _j=\pi /2`$, $`j=1,2,\mathrm{},n`$, both $`|\mathrm{}`$ and $`|\mathrm{}`$ are eigenvectors of $`\widehat{}_\text{M}^2(n\text{ even})`$ with eigenvalue $`2^{2(n1)}`$.<sup>3</sup><sup>3</sup>3It will be noted that each single-particle operator $`\widehat{\text{z}}_j`$ appears $`2^{n1}/2`$ times in Eq. (19). As can readily be seen, this very fact makes the eigenvalues corresponding to the $`2n`$ eigenvectors, $`|\mathrm{}`$, $`|\mathrm{},\mathrm{},`$ $`|\mathrm{}`$, $`|\mathrm{}`$, $`|\mathrm{},\mathrm{},\text{ and }|\mathrm{}`$ to be zero. Indeed, it can be shown that all $`2^n2`$ eigenvalues of the operator in Eq. (19) corresponding to the eigenvectors $`|z_1,z_2,\mathrm{},z_n`$ (with $`|z_1,z_2,\mathrm{},z_n|\mathrm{}\text{ or }|\mathrm{}`$) are equal to zero. Obviously, the latter is the maximum eigenvalue of $`\widehat{}_\text{M}^2(n\text{ even})`$. (Of course, as noted earlier, the same conclusions hold for the operator $`\widehat{}_\text{M}^2(n\text{ odd})`$.) It turns out that the eigenvector $`|\mathrm{\Phi }^\pm `$ of $`\widehat{}_\text{M}(n\text{ even})`$ (or $`\widehat{}_\text{M}(n\text{ odd})`$) with maximum eigenvalue $`\lambda _{\text{max}}=\pm 2^{n1}`$ does consist of an equally weighted superposition of the two eigenvectors of $`\widehat{}_\text{M}^2(n\text{ even})`$ (or $`\widehat{}_\text{M}^2(n\text{ odd})`$) with eigenvalue $`\mu _{\text{max}}=2^{2(n1)}`$,
$$|\mathrm{\Phi }^\pm =(1/\sqrt{2})(|\mathrm{}\pm e^{i\varphi }|\mathrm{}),$$
(20)
where the phase $`\varphi `$ is given by $`\varphi =\varphi _1+\varphi _2+\mathrm{}+\varphi _n+\pi /2`$. (Recall that $`\varphi _j^{}=\varphi _j+\pi /2`$ if the Bell inequalities (6) are to be maximally violated by the state $`|\mathrm{\Phi }^\pm `$.) For the Bell operator (7) discussed by Mermin we have $`\varphi _1=\varphi _2=\mathrm{}=\varphi _n=0`$, and then $`|\mathrm{\Phi }^\pm =(1/\sqrt{2})(|\mathrm{}\pm i|\mathrm{})`$ is the (nondegenerate) eigenvector of operator (7) with eigenvalue $`\pm 2^{n1}`$.
In conclusion, by expressing the square of the Mermin’s Bell operator in terms of the $`n`$ single-particle commutators $`[\widehat{\text{n}}_j,\widehat{\text{n}}_j^{}]`$, we have made explicit the relationship between the operators’ noncommutativity and the quantum-mechanical violation of the Bell inequality for the general case in which $`n`$ spin-$`\frac{1}{2}`$ particles are considered. We have seen that in order for the quantum-mechanical predictions to maximally violate Mermin’s inequality, it is necessary that all $`n`$ anticommutators vanish. Furthermore, we have shown that the maximal violation of Mermin’s inequality predicted by quantum mechanics decreases exponentially with the number of vanishing commutators. Last, but not least, it is the case that the diagonalisation of the operator $`\widehat{}_\text{M}^2(n)`$ (and hence the diagonalisation of $`\widehat{}_\text{M}(n)`$ itself) can readily be performed when such Bell operator squared is expressed in terms of the commutators.
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# Identification of the physical parameters of the paramagnetic phase of the one-dimensional Kondo lattice model done by introducting a nonmagnetic quantum state with rotating order parameters
## I INTRODUCTION
Strong correlations between electronic degrees of freedom in several materials, like high-temperature superconductors, heavy-fermion materials, etc., are responsible, or are at least thought to be responsible, for the occurrence of various quantum phenomena such as quantum magnetism, superconductivity, non-conventional metallic phases, etc. Heavy-fermion materials may be modeled using the canonical Kondo-lattice model for which intensive studies have been reported in the last years about its one-dimensional version. In this case, the zero-temperature phase diagram in terms of the conduction-electrons density $`n`$ below half-filling and the Kondo exchange coupling constant $`J_K`$ between the localized and itinerant spins has been established mainly using numerical methods.<sup>-</sup> For strong exchange couplings, the ground state is an unsaturated ferromagnet, whereas for weak couplings, it is a paramagnetic phase. At exactly half-filling, the phase is a spin-liquid insulator for any Kondo coupling, Fig. 1 (Ref. 1 contains detailed discussion of this phase diagram). The ferromagnetic and the spin-liquid phases are now well understood because these are characterized by order parameters which are the magnetization for the ferromagnetic phase, and the spin and charge energy gaps for the spin-liquid phase. For the paramagnetic phase however, no physical parameter has been proposed so far in order to describe it. One of the purposes of this work is to identify such a parameter (or parameters). Moreover, the question concerning the crossover from the high-temperarture regime to the low-temperature region needs to be addressed. In the high-temperature regime, the conduction electrons and localized spins behave as being almost independent due to the fact that thermal fluctuations wash out any characteristic energy related to the Kondo exchange coupling. In the low-temperature regime, correlation effects become strong enough to dominate the physical properties of the Kondo lattice model. The paramagnetic phase was proposed to be a Luttinger liquid at zero temperature. If we were to assume that this true, then a crossover would take place between the high-temperature metallic and the low temperature Luttinger-liquid phases.
Our findings can be summarized as follows. The paramagnetic phase is found to be governed by the quantum dynamics (spin-flip processes) of both the localized and itinerant spins. We pave the way to a new quantum state in which the local magnetization is finite in the xy-plane if measured in a rotating reference frame, but vanishes once averaged over the angle of rotation, hence ensuring that Mermin-Wagner theorem is not violated; the magnetization along the z-axis being equal to zero. This is an example of a rotating order prameter with a finite length but a phase angle assuming any value between 0 and $`2\pi `$. We should mention that interpreting our findings in terms of the Luttinger liquid state is not a simple matter. We are still investigating this question. Note however that quantitative comparison between our results and the Luttinger-liquid description results is found to be fairly good (see below). What remains to do is understanding: why should the Fermi liquid picture give rise to the Luttinger liquid one at low temperature?
In section II, a full description of the present approach is provided. The physical foundations upon which our theoretical calculations are based, are explained. We identify the physical parameters as being the averages of spin flip operators of the localized and itinerant electrons. Section III is devoted to showing some results and their discussions. The temperature dependence of the parameters, crossover temperature, phase-diagram boundary between the ferromagnetic and paramagnetic phases, as well as some of thermodynamic functions are reported. Comparison with numerical data turns out to be very satisfactory. In section IV, conclusions are drawn about the validity of the present approach.
## II Description of the approach
### A Physical foundations
In this work, the numerical data available in the literature are considered as numerical experiments upon which the present approach is founded.
The Kondo lattice Hamiltonian in one dimension is given by:
$$H=t\underset{i,\sigma }{}(c_{i,\sigma }^{}c_{i+1,\sigma }+\mathrm{H}.\mathrm{c})+J_K\underset{i}{}𝐒_i𝐬_i,$$
(1)
where $`𝐒_i`$ and $`𝐬_i`$ are respectively the localized and itinerant spin operators. $`c_{i,\sigma }^{}`$ and $`c_{i,\sigma }`$ are the creation and annihilation operators at site $`i`$ of an electron with the z-component of the spin being $`\sigma =\pm 1/2`$. $`t`$ is the hoping energy of conduction electrons. Because neither the Kondo singlets formation nor magnetic order occurs in the paramagnetic phase for densities below half-filling, a result that is strongly pointed out by numerical calculations, it is justified to choose the following canal for decoupling the interacting term of (1) where only the averages $`s_i^{}`$ and $`S_i^{}`$ and their complex conjugates are taken into account. These parameters represent spin-flip processes, and are therefore a good probe of the quantum dynamics of the system. We suppose that these processes dominate over all other processes which are the Kondo screening and magnetic ordering. We will see that our results turn out to be consistent with this hypothesis. It is worth noting here that in order to avoid breaking SU(2) symmetry, the phase angles of these parameters will not be calculated in mean-field theory. We, rather, perform a summation over all possible values of these angles. This procedure ensures that the physical phase we obtain is nonmagnetic because it does not break SU(2) symmetry.
### B Effective Hamiltonian
Using this new canal, we get the following approximate expression for the interacting term of (1):
$`𝐒_i𝐬_i`$ $`{\displaystyle \frac{1}{2}}[S_i^{}A_i^{}+s_i^{}Q_i^{}+\mathrm{H}.\mathrm{c}.2\mathrm{R}\mathrm{e}(Q_iA_i^{})],`$ (2)
where $`Q_i=S_i^{}|Q_i|e^{i\varphi _i}`$ and $`A_i=s_i^{}|A_i|e^{i\psi _i}`$. Then, the moduli $`|Q_i|`$ and $`|A_i|`$ are calculated in mean-field approximation, but summation over the phase angles $`\psi _i`$ and $`\varphi _i`$ are performed to guarantee that the continuous SU(2) symmetry remains unbroken. The total Hamiltonian is averaged over $`\varphi _i`$ and $`\psi _i`$, and the summation over the phases $`\varphi _i`$ and $`\psi _i`$ of the lowering spin operators lead to $`S_i^{}_{\varphi _i}=_0^{2\pi }d\varphi _i|Q_i|e^{i\varphi _i}=0`$ and $`s_i^{}_{\psi _i}=_0^{2\pi }d\psi _i|A_i|e^{i\psi _i}=0`$ where $`|Q_i|`$ and $`\varphi _i`$ on one hand, and $`|A_i|`$ and $`\psi _i`$ on the other hand are considered to be independent variables. The minimization of the average of the magnetic energy $`J_K𝐒_i𝐬_i=J_K|A_i||Q_i|\mathrm{cos}(\varphi _i\psi _i)`$ imposes the constraint $`\varphi _i\psi _i=\pi `$ on the angles. For this reason, for example, the sum over these phases in the last term of (2), which involves $`A_i^{}Q_i`$, is $`d(\varphi _i\psi _i\pi )|A_iQ_i|e^{i(\varphi _i\psi _i)}\delta (\varphi _i\psi _i\pi )=|A_iQ_i|`$. The delta function implements the constraint $`\varphi _i\psi _i=\pi `$.
Using the second quantization form for the itinerant spins, the Hamiltonian (1) leads to the following effective Hamiltonian where averaging over the phase angles $`\psi _i`$ and $`\varphi _i`$ is done:
$`=`$ $`{\displaystyle \frac{J_K}{2}}{\displaystyle \underset{i}{}}{\displaystyle _0^{2\pi }}{\displaystyle \frac{\mathrm{d}\psi _i}{2\pi }}S_i^+|A_i|e^{i\psi _i}+\mathrm{H}.\mathrm{c}.`$ (7)
$`+{\displaystyle \frac{J_K}{2}}{\displaystyle \underset{i}{}}{\displaystyle _0^{2\pi }}{\displaystyle \frac{\mathrm{d}\varphi _i}{2\pi }}\{Q_ic_{i,}^{}c_{i,}+Q_i^{}c_{i,}^{}c_{i,}\}`$
$`t{\displaystyle \underset{i,\sigma =,}{}}c_{i,\sigma }^{}c_{i+1,\sigma }+\mathrm{H}.\mathrm{c}`$
$`{\displaystyle \underset{i}{}}{\displaystyle }d(\psi _i\varphi _i)\times `$
$`{\displaystyle \frac{J_K}{2}}\{|A_iQ_i|e^{i[\varphi _i\psi _i]}\delta (\varphi _i\psi _i\pi )+\mathrm{C}.\mathrm{c}.\}.`$
### C Unifrom rotating configuration
In the rest of this work, we focus our attention on the uniform configuration, which is obtained for $`Q_i=|Q|e^{i\varphi }`$ and $`A_i=|A|e^{i\psi }`$. $`Q`$, $`A`$, $`\varphi `$ and $`\psi `$ are considered to be site independent but angles vary between $`0`$ and $`2\pi `$ while satisfying the constraint $`\varphi \psi =\pi `$. This state realized in this way is called the uniform rotating configuration (URC). In this state, the vector parameters $`Q_i=S_i^{}`$ and $`A_i=s_i^{}`$ point in opposite directions while rotating at the same rate. The local and total magnetizations are equal to zero.
In addition, we treat the up and down fermions as being different, and use the following canonical transformation:
$`c_k=e^{i\varphi /2}(\rho _k+\sigma _k)/\sqrt{2}`$ (8)
$`c_k=e^{i\varphi /2}(\rho _k\sigma _k)/\sqrt{2}`$ (9)
in order to diagonalize the effective Hamiltonian (7) in the case of the URC. Here, $`c_k^{()}=_ic_i^{()}e^{ir_ik}/\sqrt{N}`$ with $`N`$ being the number of lattice sites. Next, we perform the following transformation to absorb the phase terms in the transformations (9):
$`c_{k,}e^{i\varphi /2}c_{k,},c_{k,}e^{i\varphi /2}c_{k,}`$ (10)
which is equivalent to making a rotation by angle $`\varphi `$ about the z-axis for the x- and y-components of the itinerant spin operator, with the matrix of rotation given by
$`\left(\begin{array}{ccc}& \mathrm{cos}\varphi \hfill & \mathrm{sin}\varphi \hfill \\ \hfill & \mathrm{sin}\varphi \hfill & \mathrm{cos}\varphi \hfill \end{array}\right).`$
Due to the constraint $`\varphi \psi =\pi `$, rotating the itinerant-spins x and y components causes the rotation of the x and y components of the localized spins by angle $`\psi =\varphi \pi `$ about th z axis. The matrix of rotation in this case is given by:
$`\left(\begin{array}{ccc}& \mathrm{cos}\psi \hfill & \mathrm{sin}\psi \hfill \\ \hfill & \mathrm{sin}\psi \hfill & \mathrm{cos}\psi \hfill \end{array}\right).`$
Because we sum over the angles $`\varphi `$ and $`\psi `$ between $`0`$ and $`2\pi `$, the spin components are continuously rotating, and the Hamiltonian may be written in the rotating reference frame for which the x axis and y axis coincide with the rotating x and y components of the itinerant spins.
### D Mean field Hamiltonian in the rotating reference frame
The fact that the effective Hamiltonian is invariant under the above rotations by angles $`\varphi `$ and $`\psi `$ is consistent with the absence of SU(2) symmetry breaking. Here we perform such rotations before diagonalizing the simplified Hamiltonian obtained in the URC approximation. This is equivalent to using the transformations
$`c_{i,}e^{i\varphi /2}c_{i,},`$ (11)
$`c_{i,}e^{i\varphi /2}c_{i,}`$ (12)
$`S_i^+e^{i\psi }S_i^+`$ (13)
on both the itinerant and localized spins. The result is a much simpler expression for the Hamiltonian:
$`\mu N={\displaystyle \underset{k}{}}`$ $`\{E_\rho (k)\rho _k^{}\rho _k+E_\sigma (k)\sigma _k^{}\sigma _k\}`$ (15)
$`+AJ_K{\displaystyle \underset{i}{}}S_i^x+J_K{\displaystyle \underset{i}{}}AQ`$
where $`A`$ and $`Q`$ stand now for the magnitudes $`|A|`$ and $`|Q|`$ respectively. Here, $`E_{\rho ,\sigma }=ϵ(𝐤)\mu \pm |Q|J_K/2`$; $`ϵ(𝐤)=2t\mathrm{cos}k`$ is the tight-binding spectrum, and $`\mu `$ is the chemical potential. Note that we should keep in mind that (15) is obtained in a rotating frame as $`\varphi `$ and $`\psi `$ take values in the interval $`[0,2\pi [`$ subject to the constraint $`\varphi \psi =\pi `$. Therefore, it is inappropriate to interpret this Hamiltonian as that of conduction electrons and localized spins coupled to magnetic fields along the x-direction in a reference frame at rest. Now, we understand that the up and down spins could be treated differently as a consequence of the effective rotating magnetic fields $`AJ_K`$ for the localised spins, and $`QJ_K`$ for the itinerant spins. Note also that the energy spectrum of the itinerant electrons splits into two bands under the effect of the Kondo exchange interaction.
### E Phase fluctuations
Allowing for the summation over the phase angles means that we do not minimize free energy with respect to these angles. However, to satisfy the requirement that the magnetic energy is minimized we constrained their difference to be $`\pi `$. This leads us to questioning whether this minimum is stable againt fluctuations about the value $`\varphi \psi =\pi `$. To study the effect of these fluctuations, the following treatment is done. We replace the delta function in the last term of $``$ in Eq. (7) by the broader Gaussian distribution
$$\frac{1}{ϵ\sqrt{\pi }}e^{(\varphi \psi \pi )^2/ϵ^2}$$
to allow for other angles difference to contribute. If $`ϵ0`$ the Gaussian distribution reduces to the Dirac distribution. The last term in (7) takes the form:
$`{\displaystyle \mathrm{d}(\varphi \psi )}`$ (16)
$`{\displaystyle \frac{J_K}{2}}{\displaystyle \underset{i}{}}\{|AQ|e^{i[\varphi \psi ]}{\displaystyle \frac{1}{ϵ\sqrt{\pi }}}e^{(\varphi \psi \pi )^2/ϵ^2}+\mathrm{C}.\mathrm{c}.\}.`$ (17)
Integration over $`ϵ`$ from $`\mathrm{}`$ to $`+\mathrm{}`$ is undertaken, once the integration over $`\varphi \psi `$ is done, to guarantee that all possible fluctuations of the phase angles are embodied in the present approach. The result in the URC is
$$\alpha J_K\underset{i}{}|AQ|$$
which differs from the result $`J_K_i|AQ|`$, obtained without fluctuations, by the factor $`\alpha `$. $`\alpha `$ s a number larger than 1 but smaller than $`2\sqrt{\pi }`$ which is the value obtained when the integrations are carried out to infinity in the integral on $`\varphi \psi `$. This reduces the values of the mean-field parameters (as one would expect) by a factor $`\alpha 1`$ without destroying the mean-field picture. Thus this means that the mean-field solution $`\varphi \psi =\pi `$ is stable against phase fluctuations.
## III Results
### A Parameters of the KLM model
In the rest of this paper, we set $`\alpha =1`$, and seek some quantitative understanding of the present approach. To calculate the parameters $`Q`$ and $`A`$, we use the self-consistent equations obtained by minimizating the free energy:
$`F=`$ $`{\displaystyle \frac{1}{N\beta }}{\displaystyle \underset{k,\nu =\rho ,\sigma }{}}\mathrm{ln}\{1+e^{\beta E_\nu (k))}\}`$ (19)
$`{\displaystyle \frac{1}{\beta }}\mathrm{ln}\{2\mathrm{cosh}[\beta AJ_K/2]\}+J_KQA.`$
with respect to $`A`$ and $`Q`$. This yields:
$`Q=`$ $`{\displaystyle \frac{1}{2}}\mathrm{tanh}(\beta J_KA/2)`$ (20)
$`A=`$ $`{\displaystyle \frac{1}{2N}}{\displaystyle \underset{k}{}}\{f[E_\sigma (k)]f[E_\rho (k)]\},`$ (21)
where the summation $`_k`$ runs over the Brillouin zone of the conduction electrons, and $`f(x)=1/(1+e^{\beta (x\mu )})`$ is the Fermi distribution factor. $`\beta =1/k_BT`$ with $`k_B`$ is the Boltzmann constant. $`\mu ϵ_F`$ with $`ϵ_F`$ being the Fermi energy of the conduction electrons.
As the paramagnetic phase occurs for $`J_K<4t`$, we expand $`f(E_\rho )f(E_\sigma )`$ to first order in $`QJ_K/2`$. We obtain:
$`A={\displaystyle \frac{1}{2}}QJ_K{\displaystyle \underset{k}{}}{\displaystyle \frac{f(ϵ(k))}{ϵ(k)}}={\displaystyle \frac{1}{2}}Q\chi J_K,`$ (22)
which leads to
$`A=`$ $`{\displaystyle \frac{1}{4}}\chi J_K\mathrm{tanh}(\beta AJ_K/2)`$ (23)
when (21) is used. Here $`\chi =_k\frac{f(ϵ(k))}{ϵ(k)}`$ is the uniform susceptibility of the free conduction electrons in one dimension.
At zero temperature, the non-zero solution is given by $`A=\chi J_K/4`$ and $`Q=1/2`$; $`\chi 𝒟(ϵ_F)=1/2\pi v_F`$ is the density of states of the free conduction electrons; $`v_F=2t\mathrm{sin}k_F`$ being their Fermi velocity. Slightly above zero temperature, we get for $`k_BT𝒟(ϵ_F)J_K^2/4`$
$`A{\displaystyle \frac{1}{4}}\chi J_K(12e^{\beta \chi J_K^2/4})`$ (24)
$`Q{\displaystyle \frac{1}{2}}e^{\beta \chi J_K^2/4}`$ (25)
where again $`\chi 𝒟(ϵ_F)`$ is evaluated at $`T=0`$ because $`\chi (T)`$ does not deviate much from its zero-temperature value for temperature well below the Fermi temperature of the free conduction electrons. In the rest of this paper we will consider this value for $`\chi `$.
### B Crossover temperature
The parameters $`A`$ and $`Q`$ decrease as temperature increases, and vanish at a temperature $`T_{cf}`$. Indeed, expanding the tangent hyperbolic to third order in the vicinity of this critical temperature in Eq. (23) and using Eq. (23), we find:
$`AA_0\{1T/T_{cf}\}^{1/2},QQ_0\{1T/T_{cf}\}^{1/2},`$ (26)
with $`A_0=2\sqrt{3}k_BT_{cf}/J_K`$, and $`Q_0=\sqrt{3}/2`$, and the temperature
$$T_{cf}=\chi J_K^2/8k_B.$$
(27)
Below $`T_{cf}`$, the Kondo lattice system deviates from the system of independent localized spins and conduction electrons. As might be expected $`T_{cf}`$ is proportional to $`\chi J_K^2`$, a result that is reminiscent of the RKKY interaction which dominates over Kondo screening, although not leading to magnetic order. Above $`T_{cf}`$, the magnitude of the magnitization is zero. Below, $`T_{cf}`$, the magnitude becomes finite, but averaging about its phase angles gives zero magnitization. It is for this reason that this change in behavior is not a true phase transition. The present appoach leads to a sharp change of regime at $`T_{cf}`$. We do not exclude that if the change in the regime is rather a smooth crossover, the crossover temperature will be given by $`T_xT_{cf}`$. $`T_{cf}`$ is called the crossover temperature. Shibata and Tsuntsugu reported the existence of a crossover temperature, but said that it is determined by the velocity of the excitations of the itinerant spins at zero temperarture.
### C Ground-state energy and the phase diagram
The idea behind calculating the ground-state energy is to find how the correction (due to $`J_K`$) to the ground-state energy of the conduction electrons behaves. This will indicate the type of correlations that dominant. This correction is found to be given by the equation:
$`\mathrm{\Delta }E_{GS}{\displaystyle \frac{J_K^2}{8}}𝒟(ϵ_F),(|\mathrm{\Delta }E_{GS}|/k_BT_{cf}1).`$ (28)
for $`J_K4t`$. Interestingly, this correction is of order $`J_K^2/t`$, and is again reminiscent of the RKKY interaction between localized spins. As $`J_K/t`$ increases (but keeping $`J_K/t<4`$), we found that the correction to the ground-state energy deviates from a $`J_K^2/t`$ law. Estimating the points where the change takes place for different values of $`J_K`$, we determined the line boundary, $`J_{K,c}(n)/t`$ versus the conduction electrons density $`n`$, separating the paramagnetic phase from the ferromagnetic phase at strong coupling $`J_K/t`$. This line is found to be very well fitted by:
$$J_{K,c}(n)/t2\pi n\mathrm{sin}(\pi n).$$
(29)
This is consistent with the fact that the weak coupling regime is equivalent to $`J_K/t<ϵ_F/t`$. Linearizing the energy spectrum of the conduction electrons around $`k_F`$, one finds the linear spectrum $`ϵ(k)ϵ_F=v_F(kk_F)`$, where $`v_F`$ is the Fermi velocity. For a particle with momentum $`k`$ ($`\mathrm{}=1`$) and velocity $`v_F`$, the energy is given by $`v_Fk_F`$. Thus, within the linear approximation, it is natural to consider an effective Fermi energy given by $`ϵ_F=v_Fk_F=2tk_F\mathrm{sin}k_F`$. This yields $`J_K<v_Fk_F=2t\pi n\mathrm{sin}(\pi n)`$, which leads to (29).
Eq. (29) is in very good agreement with the available exact numerical data, Fig. 1. For $`n=0.35`$, Eq. (29) yields $`J_{K,c}1.96`$, which is very close to the density-matrix-renormalization-group exact result, $`2`$. The agreement is however less accurate for $`n=1/6`$ where we obtained $`J_{K,c}0.5`$ compared to the result of Ref., namely $`J_{K,c}1`$, obtained using quantum Monte Carlo simulation. The reason for this discordance may however be attributed to the fact that Monte Carlo simulations are difficult to implement at very low temperatures, and to the fact that the present theory is mean-field like. Using the bosonization technique Honner and Gulacsi, reported $`J_{K,c}(n)2.5\mathrm{sin}(\pi n)`$. But as we notice on figure 1, this does not capture the general trend, while Eq. (29) does it fairly well. Honner and Gulacsi noted that a dependence as that I report here is also possible.
### D Excitation velocities
For densities below the half-filled band, and in the presence of the Kondo coupling, we have found that the conduction electrons as charge and spin entities still form a gapless phase because the term $`\pm |Q|J_K/2`$ in $`E_{\rho ,\sigma }`$ does not open a gap in the energy spectrum. The elementary excitations are described by the operators $`\rho _k`$ and $`\sigma _k`$. Their excitation velocities are given by $`v_{\rho ,\sigma }=2t\mathrm{sin}k_{\rho ,\sigma }`$ where $`k_{\rho ,\sigma }=n\pi J_K/8t\mathrm{sin}(n\pi )`$ for small $`J_K`$. $`k_\rho `$ and $`k_\sigma `$ are determined by the conditions $`E_\rho =ϵ_F`$ and $`E_\sigma =ϵ_F`$ with $`ϵ_F=2t\mathrm{cos}(k_F)`$. Thus, the $`\rho `$\- and $`\sigma `$-velocities are given by $`v_\rho v_FJ_K\mathrm{cot}(n\pi )/4`$ and $`v_\sigma v_F+J_K\mathrm{cot}(n\pi )/4`$ for $`J_K<J_{K,c}(n)`$. This leads to
$`\mathrm{\Delta }v=v_\sigma v_\rho {\displaystyle \frac{J_K}{2}}\mathrm{cot}(n\pi ).`$ (30)
It is not clear how we could relate these velocities to the spin and charge velocities of the conduction electrons obtained within the bosonization approach. In Fig. 2, we draw $`\mathrm{\Delta }v`$ in terms of $`J_K`$ for $`n=1/3`$, and compare it with Shibata et al.’s density-matrix-renormalization-group results. Note that this comparison should be considered with caution as our results for $`v_\rho `$ and $`v_\sigma `$ disagree with theirs. Only the difference compares well. We should also emphasize here that we do not pretend by any means that the present mean-field theory is equivalent to the bosonization technique which is used to obtain the spin and charge velocities when spin-charge separation takes place in a system of interacting fermions.
### E Thermodynamic functions
Finally, we would like to briefly analyze the entropy $`S=F/T`$ and heat capacity $`C=TS/T`$. In Fig. 3, we display $`C`$, for $`n=0.35`$ and $`J_K/t=1.6`$ where $`T_{cf}=0.057t/k_B`$ agrees very well with the crossover temperature predicted by Shibata and Tsunetsugu in their numerical work . The entropy $`S`$ is also reported on the same figure. It behaves as predicted in Ref.. For $`T/T_{cf}0`$, the entropy $`S1.24T`$ yields a slope in very good agreement with the Tomonaga-Luttinger result $`S=\pi T(v_\sigma ^1+v_\rho ^1)/31.20T`$. A sharp maximum is found in $`C`$ at $`T_{cf}`$ as a consequnece of the change in the regime from hight temperatute to temperatures below $`T_{cf}`$. The linear behavior of entropy in terms of temperature is consistent with the gapless quasi-particle excitations, and is an indication of the metallic character of the paramagnetic state.
## IV conclusions
In summary, the paramagnetic phase of the one-dimensional Kondo lattice model is investigated using a new mean-field approach. This approach is valid away from half-filling where Kondo screening is negligible. We predict that this phase is described by a new nonmagnetic-quantum state characterized by two rotating order parameters with nonzero magnitudes below a crossover temperature. We calculated the ground-state line boundary in the density-coupling phase diagram, the elementary-excitations velocities of the conduction electrons, heat capacity, and entropy. The crossover temperature separating the low-temperature paramagnetic quantum phase from the normal high-temperature phase is evaluated as a function of the Kondo coupling constant and the density of states of the conduction electrons. Overall, our results agree very well with many of the available numerical data well below half-filling, and are compatible with the Luttinger-liquid picture as put forward by several authors, (see Ref. and references therein). Note that it is not yet clear how we could interpret the results of our approach in terms of the Luttinger-liquid description. For densities close to half-filling, residual effects due to Kondo screening have to be taken into account. It is then natural that, quantitatively, our approach is less accurate in this limit. This is what we come to face when we compare Shibata and Tsunetsugu’s results with ours for the entropy $`S`$ for example. The results start to agree well only for conduction-electrons densities below 0.35. Finally, the RKKY oscillations constitute an issue that needs to be addressed within the present approach. One possible avenue is to allow for the parameters $`A`$ $`Q`$ to be $`k`$-dependent. However, even with these two limitations, the present approach is very promising because it is very simple to use, and the calculations can, to a very large extent, be done analytically. And most importantly, it leads to many very satisfactory physical results.
The author would like to thank Prof. P. Fulde for his comments on the manuscript.
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# Universal trend of the information entropy of a fermion in a mean field
## I Introduction
The information entropy for a continuous probability distribution $`p(x)`$ in one dimension is defined by the expression:
$$S=p(x)\mathrm{ln}(p(x))𝑑x$$
(1)
where $`p(x)𝑑x=1`$.
$`S`$ is measured in bits if the base of the logarithm is 2 and nats (natural units of information) if the logarithm is natural. It represents the information content of a probability distribution as well as a measure of uncertainty of the corresponding state. We note that the information and thermodynamic entropy are different concepts but can be connected employing some assumptions.
An important step in the past was the discovery in of an entropic uncertainty relation (EUR), which for a three-dimensional system has the form:
$$S_r+S_k3(1+\mathrm{ln}\pi )6.434,(\mathrm{}=1)$$
(2)
where
$$S_r=\rho (\stackrel{}{r})\mathrm{ln}\rho (\stackrel{}{r})𝑑\stackrel{}{r}$$
(3)
is the information entropy in position-space
$$S_k=n(\stackrel{}{k})\mathrm{ln}n(\stackrel{}{k})𝑑\stackrel{}{k}$$
(4)
is the information entropy in momentum-space and $`\rho (\stackrel{}{r})`$, $`n(\stackrel{}{k})`$ are the density distributions in position- and momentum- space respectively, normalized to unity.
The lower bound in (2) is attained for gaussian density distributions. The physical meaning of the above inequality is the following: an increase of $`S_k`$ corresponds to a decrease of $`S_r`$ and vice versa, which indicates that a diffuse density distribution $`n(k)`$ in momentum space is associated with a localised density distribution $`\rho (r)`$ in configuration space and vice versa.
Relation (2) represents a strengthened version of Heisenberg’s uncertainty principle for two reasons: first EUR leads to Heisenberg’s uncertainty relation but the inverse is not true. Second the right-hand-side of EUR does not depend on the state of the system, while in Heisenberg’s relation does depend. It is obvious from (3) and (4) that $`S_r`$, $`S_k`$ depend on the unit of length in measuring $`\rho (\stackrel{}{r})`$ and $`n(\stackrel{}{k})`$. However, the important quantity is the entropy sum $`S_r+S_k`$ (net information content of the state) which is invariant to uniform scaling of coordinates.
Information entropy was employed in the past for the study of quantum mechanical systems \[1-9\]. Recently we studied the position- and momentum- space information entropies $`S_r`$, $`S_k`$ respectively for various systems: the nuclear density distribution of nuclei, the electron density distribution of atoms and the valence electron density distribution of atomic clusters. We showed that a similar functional form $`S=a+b\mathrm{ln}N`$ for the entropy as function of the number of particles $`N`$ holds approximately for the above systems. We conjectured that this is a universal property of a many-fermion system in a mean field.
The concept of information entropy proved also to be fruitful in a different context . We used the formalism of Ghosh, Berkowitz and Parr within the ground state density functional framework, to define the concept of an information entropy associated with the density distribution of a nuclear system. It turned out that $`S`$ increases with the quality of the wave function and can serve as a criterion of the quality of a nuclear model.
Another interesting result is the fact that the entropy of an $`N`$-photon state subjected to Gaussian noise increases linearly with the logarithm of $`N`$.
Encouraged by previous work we attempt in the present paper to calculate $`S_r`$, $`S_k`$ for the wave functions of single-particle states (instead of the total densities as in ) for various systems i.e. a nucleon in a nucleus, an electron in an atomic cluster and a $`\mathrm{\Lambda }`$ particle in a hypernucleus. We employ for these systems models existing in the literature.
Our aim is to investigate the dependence of $`S_r`$, $`S_k`$ on the excitation of a fermion in a quantum-mechanical system as well as its dependence on the system under consideration and the number of the particles $`N`$. We also attempt to connect the information entropy with the energy of the single-particle state. The study of the dependence of $`S`$ on the quantum state of a system is also interesting (as stated in ) for two reasons: (i) The information-theoretical and physical entropy are connected via Boltzmann’s constant $`k_B`$ by the Jayne’s relation $`S_{phys}=k_BS_{inf}`$. Thus one can ascribe to any quantum object a certain value of its physical entropy $`S_{phys}`$ if one calculates $`S_{inf}`$. (ii) It is interesting to know the value of the information entropy which is a measure of the spatial ”spreading out” of the wave function for various states of various systems.
The present paper is organized as follows: In sec. 2 we calculate $`S_r`$, $`S_k`$ for single-particle states of a nucleon in a nucleus as function of the number of nucleons $`N`$ using the simplest model available i.e. the harmonic oscillator potential and a more realistic one (Skyrme). In sec. 3 we calculate $`S_r`$, $`S_k`$ for a $`\mathrm{\Lambda }`$ particle in a hypernucleus employing a simple and (semi-) analytical relativistic model. In sec 4 we determine $`S_r`$, $`S_k`$ for the single-particle states of an electron in atomic (metallic) clusters using the Woods-Saxon potential. In sec. 5 we present a relationship of $`S_r+S_k`$ with the energy. Finally, sec. 6 contains a discussion of our results (comparison of sec. 2, 3 and 4) and our conclusions.
## II Information entropy for a nucleon in a nucleus
The information entropy $`S_r`$ in position-space for a single-particle wave function $`\psi (\stackrel{}{r})`$ is defined as
$$S_r=|\psi (\stackrel{}{r})|^2\mathrm{ln}|\psi (\stackrel{}{r})|^2d\stackrel{}{r}$$
(5)
while the entropy $`S_k`$ in momentum-space is
$$S_k=|\varphi (\stackrel{}{k})|^2\mathrm{ln}|\varphi (\stackrel{}{k})|^2d\stackrel{}{k}$$
(6)
where $`\varphi (\stackrel{}{k})`$ is the Fourier transform of $`\psi (\stackrel{}{r})`$.
In this section we calculate $`S_r`$ and $`S_k`$ for the single-particle states $`1s`$, $`1p`$, $`1d,\mathrm{}`$ of a nucleon in a nucleus in the framework of the harmonic oscillator (HO) model. We use for the HO parameter the well-known expression $`\mathrm{}\omega =41A^{1/3}`$MeV.
We find that the value of $`\mathrm{}\omega `$ is important only for $`S_r`$, $`S_k`$, while the net information content $`S=S_r+S_k`$ is independent of $`\mathrm{}\omega `$ and consequently of $`A`$. It depends only on the state under consideration and characterises it. These values for the states $`1s`$, $`1p`$, $`1d`$ and $`2s`$ are 6.4341, 7.8388, 8.6651 and 8.3015 respectively.
However, the HO model is a simplification. Thus, we employed a more realistic parametrization of the nuclear mean field i.e. the Skyrme (Sk III) interaction . In this model protons and neutrons move in different potentials. We choose to work with protons. However, similar results can be obtained for neutrons. We found that the values for $`S_r`$, $`S_k`$ obtained from the wave functions of single-particle states calculated according to Sk III are represented well by the expression
$$S_r(\mathrm{or}S_k)=a+bN^{1/3}$$
(7)
while $`S_r+S_k`$ is a slowly varying function of $`N`$ of the same form as (7). The values of the parameters are shown in Table 1.
In Fig. 1a we plot our fitted expressions (Sk III) $`S_r`$ $`(\mathrm{or}S_k)=a+bN^{1/3}`$ for the entropies $`S_r`$, $`S_k`$, $`S_r+S_k`$ of $`1s`$-states as functions of $`N^{1/3}`$. The lines correspond to our fitted expressions, while the corresponding values of our numerical calculations are denoted by squares for $`S_r`$, circles for $`S_k`$ and triangles for $`S_r+S_k`$. Similar graphs can be plotted for the higher states $`1p,1d,2s,\mathrm{}`$. From Fig. 1a we see that the values of the entropies are represented well by our fitted expressions. In Fig. 1b we compare the sum $`S_r+S_k=a+bN^{1/3}`$ for various single-particle states. We observe that the entropy sum $`S_r+S_k`$ enhances with the excitation of the single-particle states. We see that $`S_r+S_k`$ is a slowly varying function of $`N`$. We also note that the spin-orbit splitting is reproduced correctly i.e. the state $`1p_{3/2}`$ is lower than $`1p_{1/2}`$ e.t.c. (as for the energy) although their difference is small and cannot be shown in the figure.
## III Information entropy for a $`\mathrm{\Lambda }`$ in a hypernucleus
We employ a simple and (semi-) analytical relativistic model of a hypernucleus of ref. , where the Dirac equation with a scalar potential $`U_S(r)`$ and the fourth component of a vector potential $`U_V(r)`$ was considered in the case of rectangular shapes of these potentials with the same radius:
$$R=r_0A_{core}^{1/3}$$
In the Dirac equation was solved and gave the wave functions $`G(r)`$ and $`F(r)`$ for the large and small components for a $`\mathrm{\Lambda }`$ particle in a hypernucleus. These components can be found in relations (10) and (11) of .
The Dirac spinors in terms of large $`(G)`$ and small $`(F)`$ components can be expressed:
$$\psi =\left(\begin{array}{c}iG(r)/r\\ F(r)/r\end{array}\right)$$
(8)
The density distribution of a $`\mathrm{\Lambda }`$ in position space is:
$$\rho (r)=\frac{1}{4\pi }[G^2(r)/r^2+F^2(r)/r^2]$$
(9)
and the normalization is
$$4\pi _0^{\mathrm{}}\rho (r)r^2𝑑r=1$$
In momentum-space we have:
$$\varphi (k)=\left(\begin{array}{c}iX(k)\\ Y(k)\end{array}\right)$$
(10)
where $`X(k)`$ and $`Y(k)`$ are the Fourier transforms of $`G(r)/r`$ and $`F(r)/r`$ respectively. Thus the density distribution in momentum-space is given by:
$$n(k)=\frac{1}{4\pi }[X^2(k)+Y^2(k)]$$
(11)
and the normalization is:
$$4\pi _0^{\mathrm{}}n(k)k^2𝑑k=1$$
The information entropies of the $`\mathrm{\Lambda }`$ particle are calculated according to the relations:
$$S_r=4\pi \rho (r)\mathrm{ln}\rho (r)r^2𝑑r$$
(12)
$$S_k=4\pi n(k)\mathrm{ln}n(k)k^2𝑑k$$
(13)
where $`\rho (r)`$ and $`n(k)`$ are given by (9) and (11) respectively.
For the depths of the potential we used the values : $`D_+=30.55`$ MeV, $`D_{}=300`$ MeV, $`r_0=1.01`$ fm and the radius parameter $`R=r_0A_{core}^{1/3}`$ obtained by fitting the experimental binding energies of the ground state of the $`\mathrm{\Lambda }`$ particle. In the following we put $`A_{core}=N=`$ number of particles.
Next we fitted the expressions $`S_r`$ $`(\mathrm{or}S_k)=a+bN^{1/3}`$ to the values of $`S_r`$, $`S_k`$ calculated from (12) and (13) and found that these values are represented well. The values of the parameters $`a`$ and $`b`$ for various states are shown in Table 1.
In Fig. 2a we plot our fitted expressions for $`S_r,S_k,S_r+S_k`$ as functions of $`N^{1/3}`$ for the $`1s`$ state. This is done for a $`\mathrm{\Lambda }`$ in a hypernucleus in a similar way as for a nucleon in a nucleus (Fig 1a). Similar graphs can be plotted for the higher states. In fig. 2b we compare the sum $`S_r+S_k`$ for various single-particle states of a $`\mathrm{\Lambda }`$ (similar with Fig 1b). The spin orbit splitting is reproduced correctly as in nuclei (Sec. 2).
## IV Information entropy for an electron in an atomic cluster.
We consider atomic (metallic) clusters composed of neutral sodium atoms, where the electrons move in an effective radial electronic potential parametrized by Woods-Saxon potential of the form:
$$V_{WS}(r)=\frac{V_0}{1+\mathrm{exp}[(rR)/a]}$$
(14)
with $`V_0=6`$ eV, $`R=r_0N^{1/3}`$, $`r_0=2.25\AA `$ and $`a=0.74\AA `$. For a detailed study regarding the parametrization of Ekardt’s potentials see ref .
We found the wave functions of the single-particle states in configuration space and by Fourier transform the corresponding ones in momentum space by solving numerically the Schröndinger equation for atomic clusters for various values of the number of valence electrons $`N`$. Using the above wave functions, we calculated the information entropies $`S_r`$, $`S_k`$ (relations (5) and (6) ) for the single-particle states instead of the total density distributions as in ref. . Then we fitted the form $`S_r`$ $`(\mathrm{or}S_k)=a+bN^{1/3}`$ to these values and found that these expressions represent well the values of $`S_r`$,$`S_k`$. In Fig. 3a we plot $`S_r`$, $`S_k`$ and $`S_r+S_k`$ as functions of $`N^{1/3}`$ (similar as Fig. 1a, 2a) and in Fig. 3b we compare $`S_r+S_k`$ for various states (similar as in Fig. 1b, 2b). In Table 1 we present the values of the parameters $`a`$ and $`b`$ which were obtained from the fitting.
## V Relationship of the information entropy with the energy of single-particle states
In Fig. 4 we plot $`S_r+S_k`$, obtained with the HO model of the nucleus, versus the energy of the single-particle states. We use $`\mathrm{}\omega =41A^{1/3}`$ with $`A=208`$ (Pb) and keep the quantum number $`n`$ equal to 1. A fitting procedure gives for $`n=1`$, the relation:
$$S=k\mathrm{ln}(\mu E+\nu )$$
(15)
where $`k=2.0206`$, $`\mu =3.5373\mathrm{MeV}^1`$ and $`\nu =12.5320`$. Similar relations hold for $`n>1`$.
Next we plot the sum $`S_r+S_k`$ as function of the energy $`E`$ of single-particle states for a proton in a nucleus according to Sk III interaction for <sup>208</sup>Pb (Fig. 5) and an electron in atomic cluster with $`N=198`$ (Fig. 6) for $`n=1`$. Similar curves hold for higher values of $`n>1`$. In both cases the dependence of $`S_r+S_k`$ on $`E`$ can be represented well by the functional form (15). The values of the constants are the following
$`k=1.5262,\mu =17.3043\mathrm{MeV}^1,\nu =793.109\mathrm{for}\mathrm{a}\mathrm{proton}\mathrm{in}\mathrm{a}\mathrm{nucleus}`$
$`k=1.2386,\mu =1481.48\mathrm{eV}^1,\nu =8730.52\mathrm{for}\mathrm{an}\mathrm{electron}\mathrm{in}\mathrm{a}\mathrm{cluster}`$
A similar relation may be obtained for a $`\mathrm{\Lambda }`$ in a hypernucleus but the number of values of $`S_r+S_k`$ available is small. It is the first time in our research on information entropy that we observe such a relationship of $`S_r+S_k`$ with a fundamental quantity as the energy. We note that in Fig. 5 there are pairs of points with almost the same $`S_r+S_k`$ which reproduce the spin-orbit splitting.
## VI Discussion and Conclusions
Comparing our results in sections 2, 3, and 4, we see that a similar functional form $`S_r`$ $`(\mathrm{or}S_k)=a+bN^{1/3}`$ describes well the information entropies $`S_r`$, $`S_k`$ of the single-particle states for a nucleon in a nucleus, a $`\mathrm{\Lambda }`$ in a hypernucleus and a valence electron in an atomic cluster, although the single-particle potentials are different. We conjecture that this is a universal trend of the information entropies $`S_r`$, $`S_k`$ for a fermion in a mean field, while the net information content $`S_r+S_k`$ of the single-particle states of a fermion in a mean field is a slowly varying function of $`N`$ of the form $`S=a+bN^{1/3}`$ for the systems considered above. For nuclei and the simple HO potential $`S_r+S_k`$ is exactly a constant independent of $`N`$ i.e. $`b=0`$.
We note that in we found the universal property $`S=a+b\mathrm{ln}N`$ for the total density distributions of various systems.
In both cases it is not clear why $`S`$ depends linearly on $`\mathrm{ln}N`$ (total densities) or linearly on $`N^{1/3}`$ (single-particle states) but we note that in atomic physics there is already a connection of the information entropy with experiment i.e. with fundamental and/or experimental quantities e.g. the kinetic energy or the magnetic susceptibility. Both characteristics have been used in the study of the dynamics of atomic and molecular systems . This connection established the information entropy as an interesting entity for atomic physics. In the present paper we obtained a relationship of $`S_r+S_k`$ with a fundamental quantity as the energy of the single-particle states, i.e. $`S=k\mathrm{ln}(\mu E+\nu )`$. It is remarkable that the same functional form holds for various systems.
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# Brill-Noether Loci and the Gonality Stratification of ℳ_𝑔
## 1 Introduction
For an irreducible smooth projective complex curve $`C`$ of genus $`g`$, the gonality defined as $`\text{gon}(C)=\text{min}\{d_1:\text{there exists a }𝔤_d^1\text{ on }C\}`$ is perhaps the second most natural invariant: it gives an indication of how far $`C`$ is from being rational, in a way different from what the genus does. For $`g3`$ we consider the stratification of the moduli space $`_g`$ of smooth curves of genus $`g`$ given by gonality:
$$_{g,2}^1_{g,3}^1\mathrm{}_{g,k}^1\mathrm{}_g,$$
where $`_{g,k}^1:=\{[C]_g:C\text{ has a }𝔤_k^1\}.`$ It is well-known that the $`k`$-gonal locus $`_{g,k}^1`$ is an irreducible variety of dimension $`2g+2k5`$ when $`k(g+2)/2`$; when $`k[(g+3)/2]`$ one has that $`_{g,k}^1=_g`$ (see for instance \[AC\]). The number $`[(g+3)/2]`$ is thus the generic gonality for curves of genus $`g`$.
For positive integers $`g,r`$ and $`d`$, we introduce the Brill-Noether locus
$$_{g,d}^r=\{[C]_g:C\text{ carries a }𝔤_d^r\}.$$
The Brill-Noether Theorem (cf. \[ACGH\]) asserts that when the Brill-Noether number $`\rho (g,r,d)=g(r+1)(gd+r)`$ is negative, the general curve of genus $`g`$ has no $`𝔤_d^r`$’s, hence in this case the locus $`_{g,d}^r`$ is a proper subvariety of $`_g`$. We study the relative position of the loci $`_{g,d}^r`$ when $`r3`$ and $`\rho (g,r,d)<0`$ with respect to the gonality stratification of $`_g`$. Typically, we would like to know the gonality of a ‘general’ point $`[C]_{g,d}^r`$, or equivalently the gonality of a ‘general’ smooth curve $`C^r`$ of genus $`g`$ and degree $`d`$. Since the geometry of the loci $`_{g,d}^r`$ is very messy (existence of many components, some nonreduced and/or not of expected dimension), we will content ourselves with computing $`\text{gon}(C)`$ when $`[C]`$ is a general point of a ‘nice’ component of $`_{g,d}^r`$ (i.e. a component which is generically smooth, of the expected dimension and with general point corresponding to a curve with a very ample $`𝔤_d^r`$).
Our main result is the following:
###### Theorem 1
Let $`g15`$ and $`d14`$ be integers with $`g`$ odd and $`d`$ even, such that $`d^2>8g`$, $`4d<3g+12`$, $`d^28g+8`$ is not a square and either $`d18`$ or $`g<4d31`$. If
$$(d^{},g^{})\{(d,g),(d+1,g+1),(d+1,g+2),(d+2,g+3)\},$$
then there exists a regular component of the Hilbert scheme $`\mathrm{Hilb}`$$`_{d^{},g^{},3}`$ whose general point $`[C^{}]`$ is a smooth curve such that $`\mathrm{gon}`$$`(C^{})=`$$`\mathrm{min}`$$`(d^{}4,[(g^{}+3)/2])`$.
Here by $`\text{Hilb}_{d,g,r}`$ we denote the Hilbert scheme of curves $`C^r`$ with $`p_a(C)=g`$ and $`\text{deg}(C)=d`$. A component of $`\text{Hilb}_{d,g,r}`$ is said to be regular if its general point corresponds to a smooth irreducible curve $`C^r`$ such that the normal bundle $`N_{C/^r}`$ satisfies $`H^1(C,N_{C/^r})=0`$. By standard deformation theory (cf. \[Mod\] or \[Se\]), a regular component of $`\text{Hilb}_{d,g,r}`$ is generically smooth of the expected dimension $`\chi (C,N_{C/^r})=(r+1)d(r3)(g1)`$. Note that for $`r=3`$ the expected dimension of the Hilbert scheme is just $`4d`$. We refer to Section 4 for a natural extension of Theorem 1 for curves in higher dimensional projective spaces.
As for the numerical conditions entering Theorem 1, we note that the inequality $`d^2>8g`$ ensures the existence of smooth curves $`C^3`$ with $`g(C)=g`$ and $`\text{deg}(C)=d`$ (see Section 2), $`4d<3g+12\rho (g,3,d)<0`$ is just the condition that $`_{g,d}^3`$ is a proper subvariety of $`_g`$, while the remaining requirements are mild technical conditions.
A remarkable application of Theorem 1 is a new proof of our result (cf. \[Fa\]):
###### Theorem 2
The Kodaira dimension of the moduli space of curves of genus $`23`$ is $`2`$.
We recall that for $`g24`$ Harris, Mumford and Eisenbud proved (cf. \[HM\],\[EH\]) that $`_g`$ is of general type whereas for $`g16,g14`$ we have that $`\kappa (_g)=\mathrm{}.`$ The famous Slope Conjecture of Harris and Morrison predicts that $`_g`$ is uniruled for all $`g22`$ (see \[Mod\]). Therefore the moduli space $`_{23}`$ appears as an intriguing transition case between two extremes: uniruledness and being of general type.To put our Theorem 1 into perspective, let us note that for $`r=2`$ we have the following result of M. Coppens (cf. \[Co\]): let $`\nu :C\mathrm{\Gamma }`$ be the normalization of a general, irreducible plane curve of degree $`d`$ with $`\delta =g\left(\genfrac{}{}{0pt}{}{d1}{2}\right)`$ nodes. Assume that $`0<\delta <(d^27d+18)/2`$. Then $`\mathrm{gon}`$$`(C)=d2.`$
This theorem says that there are no $`𝔤_{d3}^1`$’s on $`C`$. On the other hand a $`𝔤_{d2}^1`$ is given by the lines through a node of $`\mathrm{\Gamma }`$. The condition $`\delta <(d^27d+18)/2`$ from the statement is equivalent with $`\rho (g,1,d3)<0`$. This is the range in which the problem is non-trivial: if $`\rho (g,1,d3)0`$, the Brill-Noether Theorem provides $`𝔤_{d3}^1`$’s on $`C`$. For $`r3`$ we might hope for a similar result. Let $`C^r`$ be a suitably general smooth curve of genus $`g`$ and degree $`d`$, with $`\rho (g,r,d)<0`$. We can always assume that $`dg1`$ (by duality $`𝔤_d^r|K_C𝔤_d^r|`$ we can always land in this range). One can expect that a $`𝔤_k^1`$ computing $`\text{gon}(C)`$ is of the form $`𝔤_d^r(D)=\{ED:E𝔤_d^r,ED\}`$ for some effective divisor $`D`$ on $`C`$. Since the expected dimension of the variety of $`e`$-secant $`(r2)`$-plane divisors
$$V_e^{r1}(𝔤_d^r):=\{DC_e:\text{ dim }𝔤_d^r(D)1\}$$
is $`2r2e`$ (cf. \[ACGH\]), we may ask whether $`C`$ has finitely many $`(2r2)`$-secant $`(r2)`$-planes (and no $`(2r1)`$-secant $`(r2)`$-planes at all). This is known to be true for curves with general moduli, that is, when $`\rho (g,r,d)0`$ (cf. \[Hir\]): for instance a smooth curve $`C^3`$ with general moduli has only finitely many $`4`$-secant lines and no $`5`$-secant lines. No such principle appears to be known for curves with special moduli.Definition: We call the number $`\text{min}(d2r+2,[(g+3)/2])`$ the expected gonality of a smooth nondegenerate curve $`C^r`$ of degree $`d`$ and genus $`g`$.
One can approach such problems from a different angle: find recipes to compute the gonality of various classes of curves $`C^r`$. Our knowledge in this respect is very scant: we know how to compute the gonality of extremal curves $`C^r`$ (that is, curves attaining the Castelnuovo bound, see \[ACGH\]) and the gonality of complete intersections in $`^3`$ (cf. \[Ba\]): If $`C^3`$ is a smooth complete intersection of type $`(a,b)`$ then $`\text{gon}(C)=abl`$, where $`l`$ is the degree of a maximal linear divisor on $`C`$. Hence an effective divisor $`DC`$ computing $`\text{gon}(C)`$ is residual to a linear divisor of degree $`l`$ in a plane section of $`C`$.
Acknowledgments: This paper is part of my thesis written at the Universiteit van Amsterdam. The help of my advisor Gerard van der Geer, and of Joe Harris, is gratefully acknowledged.
## 2 Linear systems on $`K3`$ surfaces in $`^r`$
We will construct smooth curves $`C^r`$ having the expected gonality starting with sections of smooth $`K3`$ surfaces. We recall a few basic facts about linear systems on $`K3`$ surfaces (cf. \[SD\]).
Let $`S`$ be a smooth $`K3`$ surface. For an effective divisor $`DS`$, we have $`h^1(S,D)=h^0(D,𝒪_D)1`$. If $`CS`$ is an irreducible curve then $`H^1(S,C)=0`$, and by Riemann-Roch we have that $`\text{dim}|C|=1+C^2/2=p_a(C).`$ In particular $`C^22`$ for every irreducible curve $`C`$. Moreover we have equivalences
$`C^2=2\text{dim}|C|=0C\text{ is a smooth rational curve }\text{ and }`$$`C^2=0\text{dim}|C|=1p_a(C)=1.`$For a $`K3`$ surface one also has a ‘strong Bertini’ Theorem (cf. \[SD\]):
###### Proposition 2.1
Let $``$ be a line bundle on a $`K3`$ surface $`S`$ such that $`||\mathrm{}`$. Then $`||`$ has no base points outside its fixed components. Moreover, if $`\mathrm{bs}||=\mathrm{}`$ then either
* $`^2>0,\text{ }h^1(S,)=0`$ and the general member of $`||`$ is a smooth, irreducible curve of genus $`^2/2+1`$, or
* $`^2=0`$ and $`=𝒪_S(kE)`$, where $`k_1`$, $`ES`$ is an irreducible curve with $`p_a(E)=1`$. We have that $`h^0(S,)=k+1,\text{ }h^1(S,)=k1`$ and all divisors in $`||`$ are of the form $`E_1+\mathrm{}+E_k`$ with $`E_iE`$.
We are interested in space curves sitting on $`K3`$ surfaces and the starting point is Mori’s Theorem (cf. \[Mo\]): if $`d>0`$, $`g0`$, there is a smooth curve $`C^3`$ of degree $`d`$ and genus $`g`$, lying on a smooth quartic surface $`S`$, if and only if (1) $`g=d^2/8+1`$, or (2) $`g<d^2/8`$ and $`(d,g)(5,3)`$. Moreover, we can choose $`S`$ such that $`\text{Pic}(S)=H=(4/d)C`$ in case (1) and such that $`\text{Pic}(S)=HC`$, with $`H^2=4,C^2=2g2`$ and $`HC=d`$, in case (2). In each case $`H`$ denotes a plane section of $`S`$. Note that from the Hodge Index Theorem one has the necessary condition $`(CH)^2H^2C^2=d^28(g1)0.`$
Mori’s result has been extended by Rathmann to curves in higher dimensional projective spaces (cf. \[Ra\], see also \[Kn\]): For integers $`d>0,g>0`$ and $`r3`$ such that $`d^24g(r1)+(r1)^2`$, there exists a smooth $`K3`$ surface $`S^r`$ of degree $`2r2`$ and a smooth curve $`CS`$ of genus $`g`$ and degree $`d`$ such that $`\text{Pic}(S)=HC`$, where $`H`$ is a hyperplane section of $`S`$.
We will repeatedly use the following simple observation:
###### Proposition 2.2
Let $`S^r`$ be a smooth $`K3`$ surface of degree $`2r2`$ with a smooth curve $`CS`$ such that $`\mathrm{Pic}`$$`(S)=HC`$ and assume that $`S`$ has no $`(2)`$ curves. A divisor class $`D`$ on $`S`$ is effective if and only if $`D^20`$ and $`DH>2`$.
Remark: If $`S^r`$ is a smooth $`K3`$ surface of degree $`2r2`$ with Picard number $`2`$ as above, $`S`$ has no $`(2)`$ curves when the equation
$$(r1)m^2+mnd+(g1)n^2=1$$
(1)
has no solutions $`m,n`$. This is the case for instance when $`d`$ is even and $`g`$ and $`r`$ are odd. Furthermore, a necessary condition for $`S`$ to have genus $`1`$ curves is that $`d^24(g1)(r1)`$ is a square.
## 3 Brill-Noether special linear series on curves on $`K3`$ surfaces
The first important result in the study of special linear series on curves lying on $`K3`$ surfaces was Lazarsfeld’s proof of the Brill-Noether-Petri Theorem (cf. \[Laz\]). He noticed that there is no Brill-Noether type obstruction to embed a curve in a $`K3`$ surface: if $`C_0S`$ is a smooth curve of genus $`g2`$ on a $`K3`$ surface such that $`\text{Pic}(S)=C_0`$, then the general curve $`C|C_0|`$ satisfies the Brill-Noether-Petri Theorem, that is, for any line bundle $`A`$ on $`C`$, the Petri map $`\mu _0(C,A):H^0(C,A)H^0(C,K_CA^{})H^0(C,K_C)`$ is injective. We mention that Petri’s Theorem implies (trivially) the Brill-Noether Theorem.
The general philosophy when studying linear series on a $`K3`$-section $`CS`$ of genus $`g2`$, is that the type of a Brill-Noether special $`𝔤_d^r`$ often does not depend on $`C`$ but only on its linear equivalence class in $`S`$, i.e. a $`𝔤_d^r`$ on $`C`$ with $`\rho (g,r,d)<0`$ is expected to propagate to all smooth curves $`C^{}|C|`$. This expectation, in such generality, is perhaps a bit too optimistic, but it was proved to be true for the Clifford index of a curve (see \[GL\]): for $`CS`$ a smooth $`K3`$-section of genus $`g2`$, one has that $`\text{Cliff}(C^{})=\text{Cliff}(C)`$ for every smooth curve $`C^{}|C|`$. Furthermore, if $`\text{Cliff}(C)<[(g1)/2]`$ (the generic value of the Clifford index), then there exists a line bundle $``$ on $`S`$ such that for all smooth $`C^{}|C|`$ the restriction $`_{|C^{}}`$ computes $`\text{Cliff}(C^{})`$. Recall that the Clifford index of a curve $`C`$ of genus $`g`$ is defined as
$$\text{Cliff}(C):=\text{min}\{\text{Cliff}(D):D\text{Div}(C),h^0(D)2,h^1(D)2\},$$
where for an effective divisor $`D`$ on $`C`$, we have $`\text{Cliff}(D)=\text{deg}(D)2(h^0(D)1)`$. Note that in the definition of $`\text{Cliff}(C)`$ the condition $`h^1(D)2`$ can be replaced with $`\text{deg}(D)g1`$. Another invariant of a curve is the Clifford dimension of $`C`$ defined as
$$\text{Cliff-dim}(C):=\text{min}\{r1:𝔤_d^r\text{ on }C\text{ with }dg1,\text{ such that }d2r=\text{Cliff}(C)\}.$$
Curves with Clifford dimension $`2`$ are rare: smooth plane curves are precisely the curves of Clifford dimension $`2`$, while curves of Clifford dimension $`3`$ occur only in genus $`10`$ as complete intersections of two cubic surfaces in $`^3`$.
Harris and Mumford during their work in \[HM\] conjectured that the gonality of a $`K3`$-section should stay constant in a linear system: if $`CS`$ carries an exceptional $`𝔤_d^1`$ then every smooth $`C^{}|C|`$ carries an equally exceptional $`𝔤_d^1`$. This conjecture was later disproved by Donagi and Morrison (cf. \[DMo\]). They came up with the following counterexample: let $`\pi :S^2`$ be a $`K3`$ surface, double cover of $`^2`$ branched along a smooth sextic and let $`=\pi ^{}𝒪_^2(3)`$. The genus of a smooth $`C||`$ is $`10`$. The general $`C||`$ carries a very ample $`𝔤_6^2`$, hence $`\text{gon}(C)=5`$. On the other hand, any curve in the codimension $`1`$ linear system $`|\pi ^{}H^0(^2,𝒪_^2(3))|`$ is bielliptic, therefore has gonality $`4`$. Under reasonable assumptions this turns out to be the only counterexample to the Harris-Mumford conjecture. Ciliberto and Pareschi proved that if $`CS`$ is such that $`|C|`$ is base-point-free and ample, then either $`\text{gon}(C^{})=\text{gon}(C)`$ for all smooth $`C^{}|C|`$, or $`(S,C)`$ are as in the previous counterexample (cf. \[CilP\]).
Although $`\text{gon}(C)`$ can drop as $`C`$ varies in a linear system, base-point-free $`𝔤_d^1`$’s on $`K3`$-sections do propagate:
###### Proposition 3.1 (Donagi-Morrison)
Let $`S`$ be a $`K3`$ surface, $`CS`$ a smooth, nonhyperelliptic curve and $`|Z|`$ a complete, base-point-free $`𝔤_d^1`$ on $`C`$ such that $`\rho (g,1,d)<0`$. Then there is an effective divisor $`DS`$ such that:
* $`h^0(S,D)2,\text{ h}^0(S,CD)2,\mathrm{deg}_\mathrm{C}(\mathrm{D}_{|\mathrm{C}})\mathrm{g}1.`$
* $`\mathrm{Cliff}`$$`(C^{},D_{|C^{}})\mathrm{Cliff}`$$`(C,Z)`$, for any smooth $`C^{}|C|`$.
* There is $`Z_0|Z|`$, consisting of distinct points such that $`Z_0DC`$.
Throughout this paper, for a smooth curve $`C`$ we denote, as usual, by $`W_d^r(C)`$ the scheme whose points are line bundles $`A\text{Pic}^d(C)`$ with $`h^0(C,A)r+1`$, and by $`G_d^r(C)`$ the scheme parametrizing $`𝔤_d^r`$’s on $`C`$.
## 4 The gonality of curves in $`^r`$
For a wide range of $`d,g`$ and $`r`$ we construct curves $`C^r`$ of degree $`d`$ and genus $`g`$ having the expected gonality. We start with a case when we can realize our curves as sections of $`K3`$ surfaces.
###### Theorem 3
Let $`r3,dr^2+r`$ and $`g0`$ be integers such that $`\rho (g,r,d)<0`$ and with $`d^2>4(r1)(g+r2)`$ when $`r4`$ while $`d^2>8g`$ when $`r=3`$. Let us assume moreover that $`0`$ and $`1`$ are not represented by the quadratic form
$$Q(m,n)=(r1)m^2+mnd+(g1)n^2,\text{ }\text{ }m,n.$$
Then there exists a smooth curve $`C^r`$ of degree $`d`$ and genus $`g`$ such that $`\mathrm{gon}`$$`(C)=\mathrm{min}`$$`(d2r+2,[(g+3)/2])`$. If $`\mathrm{gon}`$$`(C)=d2r+2<[(g+3)/2]`$ then $`\mathrm{dim}`$$`\text{ }W_{d2r+2}^1(C)=0`$ and every $`𝔤_{d2r+2}^1`$ is given by the hyperplanes through a $`(2r2)`$-secant $`(r2)`$-plane.
Proof: By Rathmann’s Theorem there exists a smooth $`K3`$ surface $`S^r`$ with $`\text{deg}(S)=2r2`$ and $`CS`$ a smooth curve of degree $`d`$ and genus $`g`$ such that $`\text{Pic}(S)=HC`$, where $`H`$ is a hyperplane section. The conditions $`d,g`$ and $`r`$ are subject to, ensure that $`S`$ does not contain $`(2)`$ curves or genus $`1`$ curves.
We prove first that $`\text{Cliff-dim}(C)=1`$. It suffices to show that $`CS`$ is an ample divisor, because then by using Prop.3.3 from \[CilP\] we obtain that either $`\text{Cliff-dim}(C)=1`$ or $`C`$ is a smooth plane sextic, $`g=10`$ and $`(S,C)`$ are as in Donagi-Morrison’s example (then $`\text{Cliff-dim}(C)=2`$). The latter case obviously does not happen.
We prove that $`CD>0`$ for any effective divisor $`DS`$. Let $`DmH+nC`$, with $`m,n`$, such a divisor. Then $`D^2=(2r2)m^2+2mnd+n^2(2g2)0`$ and $`DH=(2r2)m+dn>2.`$ The case $`m0,n0`$ is impossible, while the case $`m0,n0`$ is trivial. Let us assume $`m>0,n<0`$. Then $`DC=md+n(2g2)>n\left(d^2/(2r2)2g+2\right)+d/(r1)>0,`$ because $`d^2/(2r2)>2g`$. In the remaining case $`m<0,n>0`$ we have that $`nDCmDH>0`$, so $`C`$ is ample by Nakai-Moishezon.
Our assumptions imply that $`dg1`$, so $`𝒪_C(1)`$ is among the line bundles from which $`\text{Cliff}(C)`$ is computed. We get thus the following estimate on the gonality of $`C`$:
$$\text{gon}(C)=\text{Cliff}(C)+2\text{Cliff}(C,H_{|C})+2=d2r+2,$$
which yields $`\text{gon}(C)\text{min}(d2r+2,[(g+3)/2]).`$
For the rest of the proof let us assume that $`\text{gon}(C)<[(g+3)/2]`$. We will then show that $`\text{gon}(C)=d2r+2.`$ Let $`|Z|`$ be a complete, base point free pencil computing $`\text{gon}(C).`$ By applying Prop.3.1, there exists an effective divisor $`DS`$ satisfying
$$h^0(S,D)2,h^0(S,CD)2,\text{deg}(D_{|C})g1,\text{ gon}(C)=\text{Cliff}(D_{|C})+2\text{ and }ZDC.$$
We consider the exact cohomology sequence:
$$0H^0(S,DC)H^0(S,D)H^0(C,D_{|C})H^1(S,DC).$$
Since $`CD`$ is effective and $`0`$, one sees that $`DC`$ cannot be effective, so $`H^0(S,DC)=0`$. The surface $`S`$ does not contain $`(2)`$ curves, so $`|CD|`$ has no fixed components; the equation $`(CD)^2=0`$ has no solutions, therefore $`(CD)^2>0`$ and the general element of $`|CD|`$ is smooth and irreducible. Then it follows that $`H^1(S,DC)=H^1(S,CD)^{}=0.`$ Thus $`H^0(S,D)=H^0(C,D_{|C})`$ and
$$\text{gon}(C)=2+\text{Cliff}(D_{|C})=2+DC2\text{ dim}|D|=DCD^2.$$
We consider the following family of effective divisors
$$𝒜:=\{D\text{Div}(S):h^0(S,D)2,h^0(S,CD)2,\text{ }CDg1\}.$$
Since we already know that $`d2r+2\text{gon}(C)\alpha `$, where $`\alpha =\text{min}\{DCD^2:D𝒜\}`$, we are done if we prove that $`\alpha d2r+2`$. Take $`D𝒜`$ such that $`DmH+nC`$, $`m,n`$. The conditions $`D^2>0,DCg1`$ and $`2<DH<d2`$ (use Prop.2.2 for the last inequality) can be rewritten as
$$(r1)m^2+mnd+n^2(g1)>0\text{ (i), }2<(2r2)m+nd<d2\text{ (ii), }md+(2n1)(g1)0\text{ (iii)}.$$
We have to prove that for any $`D𝒜`$ the following inequality holds
$$f(m,n)=DCD^2=(2r2)m^2+m(d2nd)+(nn^2)(2g2)f(1,0)=d2r+2.$$
We solve this standard calculus problem. Denote by
$$a:=\frac{d+\sqrt{d^24(r1)(g1)}}{2r2}\text{ }\text{ and }\text{ }b:=\frac{d\sqrt{d^24(r1)(g1)}}{2r2}\text{ }.$$
We dispose first of the case $`n<0`$. Assuming $`n<0`$, from (i) we have that either $`m<bn`$ or $`m>an`$. If $`m<bn`$ from (ii) we obtain that $`2<n(d(2r2)b)<0`$, because $`n<0`$ and $`d(2r2)b=\sqrt{d^24(r1)(g1)}>0`$, so we have reached a contradiction.
We assume now that $`n<0`$ and $`m>an`$. From (iii) we get that $`m(g1)(12n)/d`$. If $`an>(g1)(12n)/d`$ we are done because there are no $`m,n`$ satisfying (i), (ii) and (iii), while in the other case for any $`D𝒜`$ with $`DmH+nC`$, one has the inequality
$$f(m,n)>f(an,n)=\frac{(d^24(r1)(g1))+d\sqrt{d^24(r1)(g1)}}{2r2}(n).$$
When $`r4`$ since we assume that $`\sqrt{d^24(r1)(g1)}2r2`$, it immediately follows that $`f(m,n)d>d2r+2.`$ In the case $`r=3`$ when we only have the weaker assumption $`d^2>8g`$, we still get that $`f(an,n)>d4`$ unless $`n=1`$ and $`d^28g<8`$. In this last situation we obtain $`m(d+4)/4`$ so $`f(m,1)f((d+4)/4,1)>d4`$.
The case $`n>0`$ can be treated rather similarly. From (i) we get that either $`m<an`$ or $`m>bn`$. The first case can be dismissed immediately. When $`m>bn`$ we use that for any $`D𝒜`$ with $`DmH+nC`$,
$$f(m,n)\text{min}\{f((g1)(2n1)/d,n),\text{max}\{f(bn,n),f((2nd)/(2r2),n)\}\}.$$
Elementary manipulations give that
$$f((g1)(2n1)/d,n)=(g1)/2[(2n1)^2(d^24(r1)(g1))/d^2+1]d2r+2$$
(use only that $`dg1`$ and $`d^2>4(r1)g`$, so we cover both cases $`r=3`$ and $`r4`$ at once). Note that in the case $`n>0`$ we have equality if and only if $`n=1,m=1`$ and $`d=g1`$.
Moreover $`f(bn,n)=n(2g2bd)2g2bd`$ and $`2g2bd>d2r+22r2<\sqrt{d^24(r1)(g1)}<d2r+2`$. When this does not happen we proceed as follows: if $`\sqrt{d^24(r1)(g1)}d2r+2`$ then if $`n=1`$ we have that $`m>b1`$, that is $`m0`$, but this contradicts (ii). When $`n2`$, we have $`f((2nd)/(2r2),n)=[(d^24(r1)(g1))(n^2n)+(2d4)]/(2r2)>d2r+2.`$ Finally, the remaining possibility $`2r2\sqrt{d^24(r1)(g1)}`$ does not occur when $`r4`$ while in the case $`r=3`$ we either have $`f(bn,n)>d4`$ or else $`n=1`$ and then $`m>(d+4)/4`$ hence $`f(m,1)>f((d+4)/4,1)=d4.`$
All this leaves us with the case $`n=0`$, when $`f(m,0)=(2r2)m^2+md`$. Clearly $`f(m,0)f(1,0)`$ for all $`m`$ complying with (i),(ii) and (iii).
Thus we proved that $`\text{gon}(C)=d2r+2`$. We have equality $`DCD^2=d2r+2`$ where $`D𝒜`$, if and only if $`D=H`$ or in the case $`d=g1`$ also when $`D=CH`$. The latter possibility can be ruled out since $`d=g1`$ is not compatible with the assumptions $`dr^2+r`$ and $`d2r+2<[(g+3)/2]`$. Therefore we can always assume that the divisor on $`S`$ cutting a $`𝔤_{d2r+2}^1`$ on $`C`$ is the hyperplane section of $`S`$. Since $`ZHC`$, if we denote by $`\mathrm{\Delta }`$ the residual divisor of $`Z`$ in $`HC`$, we have that $`h^0(C,H_{|C}\mathrm{\Delta })=2`$, so $`\mathrm{\Delta }`$ spans a $`^{r2}`$ hence $`|Z|`$ is given by the hyperplanes through the $`(2r2)`$-secant $`(r2)`$-plane $`\mathrm{\Delta }`$. This shows that every pencil computing $`\text{gon}(C)`$ is given by the hyperplanes through a $`(2r2)`$-secant $`(r2)`$-plane.
There are a few ways to see that $`C`$ has only finitely many $`(2r2)`$-secant $`(r2)`$-planes. The shortest is to invoke Theorem 3.1 from \[CilP\]: since $`\text{gon}(C^{})=d2r+2`$ is constant as $`C^{}`$ varies in $`|C|`$, for the general smooth curve $`C^{}|C|`$ one has $`\text{dim }W_{d2r+2}^1(C^{})=0`$. $`\mathrm{}`$ Remarks: 1. Keeping the assumptions and the notations of Theorem 3 we note that when $`d2r+2<[(g+3)/2]`$ the linear system $`|C|`$ is $`(d2r1)`$-very ample, i.e. for any $`0`$-dimensional subscheme $`ZS`$ of length $`d2r`$ the map $`H^0(S,C)H^0(S,C𝒪_Z)`$ is surjective. Indeed, by applying Theorem 2.1 from \[BS\] if $`|C|`$ is not $`(d2r1)`$-very ample, there exists an effective divisor $`D`$ on $`S`$ such that $`C2D`$ is $``$-effective and
$$CD(d2r)D^2CD/2<d2r,$$
hence $`CDD^2d2r`$. On the other hand clearly $`D𝒜`$, thus $`CDD^2d2r+2`$, a contradiction.
2. One can find quartic surfaces $`S^3`$ containing a smooth curve $`C`$ of degree $`d`$ and genus $`g`$ in the case $`g=d^2/8+1`$ (which is outside the range Theorem 3 deals with). Then $`d=4m,g=2m^2+1`$ with $`m1`$ and $`C`$ is a complete intersection of type $`(4,m)`$. For such a curve, $`\text{gon}(C)=dl`$, where $`l`$ is the degree of a maximal linear divisor on $`C`$ (cf. \[Ba\]). If $`S`$ is sufficiently general so that it contains no lines, by Bezout, $`C`$ cannot have $`5`$-secant lines so $`\text{gon}(C)=d4`$ in this case too.
When $`r=3`$ we want to find out when the curves constructed in Theorem 3 correspond to smooth points of $`\text{Hilb}_{d,g,3}`$. We have the following:
###### Proposition 4.1
Let $`CS^3`$ be a smooth curve sitting on a quartic surface such that $`\mathrm{Pic}`$$`(S)=HC`$ with $`H`$ being a plane section and assume furthermore that $`S`$ contains no $`(2)`$ curves. Then $`H^1(C,N_{C/^3})=0`$ if and only if $`d18\text{ or }g<4d31`$.
Proof: We use the exact sequence
$$0N_{C/S}N_{C/^3}N_{S/^3}𝒪_C0,$$
(2)
where $`N_{S/^3}𝒪_C=𝒪_C(4)`$ and $`N_{C/S}=K_C`$. We claim that there is an isomorphism $`H^1(C,N_{C/^3})=H^1(C,𝒪_C(4))`$. Suppose this is not the case. Then the injective map $`H^1(C,K_C)H^1(C,N_{C/^3})`$ provides a section $`\sigma H^0(N_{C/^3}^{}K_C)`$ which yields a splitting of the dual of the exact sequence (2), hence (2) is split as well. Using a result from \[GH, p.252\] we obtain that $`C`$ is a complete intersection with $`S`$. This is clearly a contradiction. Therefore one has $`H^1(C,N_{C/^3})=H^1(C,𝒪_C(4))`$.
We have isomorphisms $`H^1(C,4H_{|C})=H^2(S,4HC)=H^0(S,C4H)^{}`$. According to Prop.2.2 the divisor $`C4H`$ is effective if and only if $`(C4H)^20`$ and $`(C4H)H>2`$, from which the conclusion follows. $`\mathrm{}`$ We need to determine the gonality of nodal curves not of compact type and which consist of two components meeting at a number of points. We have the following result:
###### Proposition 4.2
Let $`C=C_1_\mathrm{\Delta }C_2`$ be a quasi-transversal union of two smooth curves $`C_1`$ and $`C_2`$ meeting at a finite set $`\mathrm{\Delta }`$. Denote by $`g_1=g(C_1),g_2=g(C_2),\delta =\mathrm{card}(\mathrm{\Delta })`$. Let us assume that $`C_1`$ has only finitely many pencils $`𝔤_d^1`$, where $`\delta d`$ and that the points of $`\mathrm{\Delta }`$ do not occur in the same fibre of one of these pencils. Then $`\mathrm{gon}`$$`(C)d+1`$. Moreover if $`\mathrm{gon}`$$`(C)=d+1`$ then either (1) $`C_2`$ is rational and there is a degree $`d`$ map $`f_1:C_1^1`$ and a degree $`1`$ map $`f_2:C_2^1`$ such that $`f_{1|\mathrm{\Delta }}=f_{2|\mathrm{\Delta }}`$, or (2) there is a $`𝔤_{d+1}^1`$ on $`C_1`$ containing $`\mathrm{\Delta }`$ in a fibre.
Proof: Let us assume that $`C`$ is $`k`$-gonal, that is, a limit of smooth $`k`$-gonal curves. If $`g=g_1+g_2+\delta 1`$, we consider the space $`\overline{}_{g,k}`$ of Harris-Mumford admissible coverings of degree $`k`$ and we denote by $`\pi :\overline{}_{g,k}\overline{}_g`$ the proper map sending a covering to the stable model of its domain (see \[HM\]). Since $`[C]\overline{}_{g,k}^1=\text{Im}(\pi )`$, it follows that there exists a semistable curve $`C^{}`$ whose stable model is $`C`$ and a degree $`k`$ admissible covering $`f:C^{}Y`$, where $`Y`$ is a semistable curve of arithmetic genus $`0`$. We thus have that $`f^1(Y_{sing})=C_{sing}^{}`$ and if $`pC_1^{}C_2^{}`$ with $`C_1^{}`$ and $`C_2^{}`$ components of $`C^{}`$, then $`f(C_1^{})`$ and $`f(C_2^{})`$ are distinct components of $`Y`$ and the ramification indices at the point $`p`$ of the restrictions $`f_{|C_1^{}}`$ and $`f_{|C_2^{}}`$ are the same.
We have that $`C^{}=C_1C_2R_1\mathrm{}R_\delta `$, where for $`1i\delta `$ the curve $`R_i`$ is a (possibly empty) destabilizing chain of $`^1`$’s inserted at the nodes of $`C`$. Let us denote $`\{p_i\}=C_1R_i`$ and $`\{q_i\}=C_2R_i`$; if $`R_i=\mathrm{}`$ then we take $`p_i=q_i\mathrm{\Delta }C`$.
We first show that $`kd+1`$. Suppose $`kd`$. Since $`C_1`$ has no $`𝔤_{d1}^1`$’s it follows that $`k=d`$ and that $`f^1f(C_1)=C_1`$. If there were distinct points $`p_i`$ and $`p_j`$ such that $`f(p_i)f(p_j)`$, then $`f(R_i)f(R_j)`$ and the image curve $`Y`$ would no longer have genus $`0`$. Therefore $`f(p_i)=f(p_j)`$ for all $`i,j\{1,\mathrm{},\delta \}`$, that is $`\mathrm{\Delta }`$ appears in the fibre of a $`𝔤_d^1`$ on $`C_1`$, a contradiction.
Assume now that $`k=d+1`$. Then either $`\text{deg}(f_{|C_1})=d`$ or $`\text{deg}(f_{|C_1})=d+1.`$ If $`\text{deg}(f_{|C_1})=d+1`$, then again $`f^1f(C_1)=C_1`$ and by the same reasoning $`f`$ maps all the $`p_i`$’s to the same point and this yields case (2) from the statement of the Proposition. If $`\text{deg}(f_{|C_1})=d`$ then $`f^1f(C_1)=C_1D`$, where $`D`$ is a smooth rational curve mapped isomorphically to its image via $`f`$. If $`D=C_2`$ then the condition that the dual graph of $`Y`$ is a tree implies that $`f(p_i)=f(q_i)`$ for all $`i`$ and this yields case (1) from the statement. Finally, if $`DC_2`$ then $`f(C_1)f(C_2)`$. We know that there are $`1i<j\delta `$ such that $`f(p_i)f(p_j)`$. The image $`f(C_2)`$ belongs to a chain $`R`$ of $`^1`$’s such that either $`Rf(C_1)=\{f(p_i)\}`$ or $`Rf(C_1)=\{f(p_j)\}`$. In the former case $`f(p)=f(p_i)`$ for all $`p\mathrm{\Delta }\{p_j\}`$ while in the latter case $`f(p)=f(p_j)`$ for all $`p\mathrm{\Delta }\{p_i\}.`$ In each case by adding a base point we obtain a $`𝔤_{d+1}^1`$ on $`C_1`$ containing $`\mathrm{\Delta }`$ in a fibre. $`\mathrm{}`$
Theorem 3 provides curves $`C^3`$ of expected gonality when $`d`$ is even and $`g`$ is odd (equation (1) has no solutions in this case). Naturally, we would like to have such curves when $`d`$ and $`g`$ have other parities as well. We will achieve this by attaching to a ‘good’ curve of expected gonality either a $`2`$ or $`3`$-secant line or a $`4`$-secant conic.Theorem 1 Let $`g15`$ and $`d14`$ be integers with $`g`$ odd and $`d`$ even, such that $`d^2>8g,4d<3g+12`$, $`d^28g+8`$ is not a square and either $`d18`$ or $`g<4d31`$. If
$$(d^{},g^{})\{(d,g),(d+1,g+1),(d+1,g+2),(d+2,g+3)\},$$
then there exists a regular component of $`\mathrm{Hilb}_{\mathrm{d}^{},\mathrm{g}^{},3}`$ with general point $`[C^{}]`$ a smooth curve such that $`\mathrm{gon}`$$`(C^{})=\mathrm{min}(\mathrm{d}^{}4,[(\mathrm{g}^{}+3)/2])`$. Proof: For $`d`$ and $`g`$ as in the statement we know by Theorem 3 and by Prop.4.1 that there exists a smooth nondegenerate curve $`C^3`$ of degree $`d`$ and genus $`g`$, with $`\text{gon}(C)=\text{min}(d4,[(g+3)/2])`$ and $`H^1(C,N_{C/^3})=0`$. We can also assume that $`C`$ sits on a smooth quartic surface $`S`$ and $`\text{Pic}(S)=HC`$. Moreover, in the case $`d4<[(g+3)/2]`$ the curve $`C`$ has only finitely many $`𝔤_{d4}^1`$’s, all given by planes through a $`4`$-secant line. i) Let us settle first the case $`(d^{},g^{})=(d+1,g+1)`$. Take $`p,qC`$ general points, $`L=\overline{pq}^3`$ and $`X:=CL`$. By applying Lemma 1.2 from \[BE\], we know that $`H^1(X,N_X)=0`$ and the curve $`X`$ is smoothable in $`^3`$, that is, there exists a flat family of curves $`\{X_t\}`$ in $`^3`$ over a smooth and irreducible base, with the general fibre $`X_t`$ smooth while the special fibre $`X_0`$ is $`X`$. If $`d4<[(g+3)/2]`$, then since $`C`$ has only finitely many $`𝔤_{d4}^1`$’s, by applying Prop.4.2 we get that $`\text{gon}(X)=d3`$. In the case $`d4[(g+3)/2]`$ we just notice that $`\text{gon}(X)\text{gon}(C)=[(g^{}+3)/2]`$. ii) Next we tackle the case $`(d^{},g^{})=(d+1,g+2)`$. Assume first that $`d4<[(g+3)/2]d^{}4<[(g^{}+3)/2]`$. We apply Lemma 1.2 from \[BE\] to a curve $`X:=CL`$, where $`L`$ is a suitable trisecant line to $`C`$. In order to conclude that $`X`$ is smoothable in $`^3`$ and that $`H^1(X,N_X)=0`$, we have to make sure that the trisecant line $`L=\overline{pqq^{}}`$ with $`p,q,q^{}C`$ can be chosen in such a way that
$$L,T_p(C),T_q(C)\text{ and }T_q^{}(C)\text{ do not all lie in the same plane. }$$
(3)
We claim that when $`C|C|`$ is general in its linear system, at least one of its trisecants satisfies (3). Suppose not. Then for every smooth curve $`C|C|`$ and for every trisecant line $`L`$ to $`C`$ condition (3) fails.
We consider a $`0`$-dimensional subscheme $`ZS`$ where $`Z=p+q+q^{}+u+u^{}`$, with $`p,q,q^{}S`$ being collinear points while $`u`$ and $`u^{}`$ are general infinitely near points to $`q`$ and $`q^{}`$ respectively. The linear system $`|C|`$ is at least $`5`$-very ample (cf. Remark 1), hence a general curve $`C|CZ|`$ is smooth and possesses a trisecant line for which (3) holds, a contradiction.
Since the scheme of trisecants to a space curve is of pure dimension $`1`$, it follows that for a general curve $`C|C|`$, through a general point $`pC`$ there passes a trisecant line $`L`$ for which (3) holds. We have that $`X:=CL`$ is smoothable in $`^3`$ and $`H^1(X,N_X)=0`$. We conclude that $`\text{gon}(X)=d3`$ by proving that there is no $`𝔤_{d4}^1`$ on $`C`$ containing $`LC`$ in a fibre.
If $`C|C|`$ is general, any line in $`^3`$ (hence also a $`4`$-secant line to $`C`$) can meet only finitely many trisecants. Indeed, assuming that $`m^3`$ is a line meeting infinitely many trisecants, we consider the correspondence
$$T=\{(p,t)C\times m:\overline{pt}\text{ is a trisecant to }C\}$$
and the projections $`\pi _1:TC`$ and $`\pi _2:Tm`$. If $`\pi _2`$ is surjective, then $`\text{Nm}_{\pi _1}(\pi _2)`$ yields a $`𝔤_3^1`$ on $`C`$, a contradiction. If $`\pi _2`$ is not surjective then there exists a point $`t^3`$ such that $`\overline{pt}`$ is a trisecant to $`C`$ for each $`pC`$. This possibility cannot occur for a general $`C|C|`$: Otherwise we take general points $`t^3`$ and $`p,p^{}S`$ and if we denote
$$:=\{C|C|:p,p^{}C\text{ and }\overline{tx}\text{ is a trisecant to }C\text{ for each }xC\},$$
we have that $`\text{dim }g5`$. On the other hand since $`\overline{tp}`$ and $`\overline{tp^{}}`$ are trisecants for all curves $`C`$, there must be a $`0`$-dimensional subscheme $`Z(\overline{tp}\overline{tp^{}})S`$ of length $`6`$ such that $`|CZ|`$, hence $`\text{dim}\text{dim}|CZ|=g6`$ (use again that $`|C|`$ is $`5`$-very ample), a contradiction. In this way the case $`d4<[(g+3)/2]`$ is settled.
When $`d4[(g+3)/2]`$ we apply Theorem 3 to obtain a smooth curve $`C_1^3`$ of degree $`d`$ and genus $`g+2`$ such that $`\text{gon}(C_1)=(g+5)/2`$ and $`H^1(C_1,N_{C_1})=0`$. We take $`X_1:=C_1L_1`$ with $`L_1`$ being a general $`1`$-secant line to $`C_1`$. Then $`X_1`$ is smoothable and $`\text{gon}(X_1)=\text{gon}(C_1)=(g+5)/2`$. iii) Finally, we turn to the case $`(d^{},g^{})=(d+2,g+3)`$. Take $`H^3`$ a general plane meeting $`C`$ in $`d`$ distinct points in general linear position and pick 4 of them: $`p_1,p_2,p_3,p_4CH`$. Choose $`QH`$ a general conic such that $`QC=\{p_1,p_2,p_3,p_4\}`$. Theorem 5.2 from \[Se\] ensures that $`X:=CQ`$ is smoothable in $`^3`$ and $`H^1(X,N_X)=0`$.
Assume first that $`d^{}4[(g^{}+3)/2]`$. We claim that $`\text{gon}(X)\text{gon}(C)+2`$. According to Prop.4.2 the opposite could happen only in 2 cases: a) There exists a $`𝔤_{d3}^1`$ on $`C`$, say $`|Z|`$, such that $`|Z|(p_1p_2p_3p_4)\mathrm{}.`$ b) There exists a degree $`d4`$ map $`f:C^1`$ and a degree $`1`$ map $`f^{}:Q^1`$ such that $`f(p_i)=f^{}(p_i)`$, for $`i=1,\mathrm{},4`$.
Assume that a) does happen. We denote by $`U=\{DC_4:|𝒪_C(1)|(D)\mathrm{}\}`$ the irreducible $`3`$-fold of divisors of degree $`4`$ spanning a plane and also consider the correspondence
$$\mathrm{\Sigma }=\{(L,D)W_{d3}^1(C)\times U:|L|(D)\mathrm{}\},$$
with the projections $`\pi _1:\mathrm{\Sigma }W_{d3}^1(C)`$ and $`\pi _2:\mathrm{\Sigma }U`$. We know that $`\pi _2`$ is dominant, hence $`\text{dim }\mathrm{\Sigma }3`$ and therefore $`\text{dim }W_{d3}^1(C)2`$.
If $`\rho (g,1,d3)<0`$ by Prop.3.1 we get that every base-point-free $`𝔤_{d3}^1`$ on $`C`$ is cut out by a divisor $`D`$ on $`S`$ such that $`D𝒜`$ (see the proof of Theorem 3 for this notation) and $`CDD^2=\text{Cliff}(C,D_{|C})+2d3`$, hence $`CDD^2d4`$ for parity reasons. As pointed out at the end of the proof of Theorem 3 this forces $`DH`$, that is, all base-point-free $`𝔤_{d3}^1`$’s on $`C`$ are given by planes through a trisecant line. Thus $`C`$ has $`\mathrm{}^2`$ trisecants, a contradiction.
If $`\rho (g,1,d3)0`$, then $`g=2d9`$ and we can assume that there is $`L\pi _1(\mathrm{\Sigma })`$ such that $`|𝒪_C(1)L|\mathrm{}`$. The map $`\pi _1`$ is either generically finite hence $`\text{dim }W_{d4}^1(C)\text{dim }W_{d3}^1(C)21`$ (cf. \[FHL\]), a contradiction, or otherwise $`\pi _1`$ has fibre dimension $`1`$. This is possible only when there is a component $`A`$ of $`W_{d3}^1(C)`$ with $`\text{dim}(A)2`$ and such that the general $`LA`$ satisfies $`|𝒪_C(1)L|\mathrm{}`$ and every $`LA`$ has non-ordinary ramification so that the monodromy of each $`𝔤_{d3}^1`$ is not the full symmetric group. Applying again \[FHL\] there is $`LW_{d4}^1(C)`$ such that $`\{L\}+W_1^0(C)A`$, in particular $`L`$ has non-ordinary ramification too. It is easy to see that this contradicts the $`(d7)`$-very ampleness of $`|C|`$ asserted by Remark 1.
We now rule out case b). Suppose that b) does happen and denote by $`L^3`$ the $`4`$-secant line corresponding to $`f`$. Let $`\{p\}=LH`$, and pick $`lH`$ a general line. As $`Q`$ was a general conic through $`p_1,\mathrm{},p_4`$ we may assume that $`pQ`$. The map $`f^{}:Ql`$ is (up to a projective isomorphism of $`l`$) the projection from a point $`qQ`$, while $`f(p_i)=\overline{p_ip}l`$, for $`i=1,\mathrm{},4`$. By Steiner’s Theorem from classical projective geometry, the condition $`(f(p_1)f(p_2)f(p_3)f(p_4))=(f^{}(p_1)f^{}(p_2)f^{}(p_3)f^{}(p_4))`$ is equivalent with $`p_1,p_2,p_3,p_4,p`$ and $`q`$ being on a conic, a contradiction since $`pQ`$.
Finally, when $`d^{}4>[(g^{}+3)/2]`$, we have to show that $`\text{gon}(X)\text{gon}(C)+1`$. We note that $`\text{dim }G_{(g+3)/2}^1(C)=1`$ (for any curve one has the inequality $`\text{dim }G_{\text{gon}}^11`$). By taking $`H(^3)^{}`$ general enough, we obtain that $`p_1,\mathrm{},p_4`$ do not occur in the same fibre of a $`𝔤_{(g+3)/2}^1`$. $`\mathrm{}`$ Remark: Theorem 1 can be viewed as a non-containment relation $`_{g^{},d^{}}^3_{g^{},d^{}5}^1`$ between different Brill-Noether loci when $`d^{}`$ and $`g^{}`$ are as in Theorem 1 and moreover $`d^{}4[(g^{}+3)/2]`$. We can turn this problem on its head and ask the following question: given $`g`$ and $`k`$ such that $`k<(g+2)/2`$, when is it true that the general $`k`$-gonal curve of genus $`g`$ has no other linear series $`𝔤_d^r`$ with $`\rho (g,r,d)<0`$, that is, the pencil computing the gonality is the only Brill-Noether exceptional linear series?
In \[Fa2\] we prove using limit linear series the following result: fix $`g`$ and $`k`$ positive integers such that $`3\rho (g,1,k)<0`$. If $`\rho (g,1,k)=3`$ assume furthermore that $`k6`$. Then the general $`k`$-gonal curve $`C`$ of genus $`g`$ has no $`𝔤_d^r`$’s with $`\rho (g,r,d)<0`$ except $`𝔤_k^1`$ and $`|K_C𝔤_k^1|`$. In other words the $`k`$-gonal locus $`_{g,k}^1`$ is not contained in any other proper Brill-Noether locus $`_{g,d}^r`$ with $`r2,dg1`$ and $`\rho (g,r,d)<0`$.
In seems that other methods are needed to extend this result for more negative values of $`\rho (g,1,k)`$.
## 5 The Kodaira dimension of $`_{23}`$
In this section we explain how Theorem 1 gives a new proof of our result $`\kappa (_{23})2`$ (cf. \[Fa\]). We refer to \[Fa\] for a detailed analysis of the geometry of $`_{23}`$ ; in that paper we also conjecture that $`\kappa (_{23})=2`$ and we present evidence for such a possibility.
Let us denote by $`\overline{}_g`$ the moduli space of Deligne-Mumford stable curves of genus $`g`$. We study the multicanonical linear systems on $`\overline{}_{23}`$ by exhibiting three explicit multicanonical divisors on $`\overline{}_{23}`$ which are (modulo a positive combination of boundary classes coming from $`\overline{}_{23}_{23}`$) of Brill-Noether type, that is, loci of curves having a $`𝔤_d^r`$ when $`\rho (23,r,d)=1`$.
On $`_{23}`$ there are three Brill-Noether divisors corresponding to the solutions of the equation $`\rho (23,r,d)=1`$: the $`12`$-gonal divisor $`_{23,12}^1`$, the divisor $`_{23,17}^2`$ of curves having a $`𝔤_{17}^2`$ and finally the divisor $`_{23,20}^3`$ of curves possessing a $`𝔤_{20}^3`$. If we denote by $`\overline{}_{g,d}^r`$ the closure of $`_{g,d}^r`$ inside $`\overline{}_g`$, the classes $`[\overline{}_{g,d}^r]\text{Pic}_{}(\overline{}_g)`$ when $`\rho (g,r,d)=1`$ have been computed (see \[EH\],\[Fa\]). It is quite remarkable that for fixed $`g`$ all classes $`[\overline{}_{g,d}^r]`$ are proportional. One also knows the canonical divisor class (cf. \[HM\]):
$$K_{\overline{}_g}=13\lambda 2\delta _03\delta _12\delta _2\mathrm{}2\delta _{[g/2]},$$
and by comparing for $`g=23`$ this formula with the expression of the classes $`[\overline{}_{23,d}^r]`$, we find that there are constants $`m,m_1,m_2,m_3_{>0}`$ such that
$$mK_{\overline{}_{23}}=m_1[\overline{}_{23,12}^1]+E=m_2[\overline{}_{23,17}^2]+E=m_3[\overline{}_{23,20}^3]+E,$$
where $`E`$ is the same positive combination of the boundary classes $`\delta _1,\mathrm{},\delta _{11}`$.
As explained in \[Fa\], since $`\overline{}_{23,12}^1`$, $`\overline{}_{23,17}^2`$ and $`\overline{}_{23,20}^3`$ are mutually distinct irreducible divisors, we can show that the multicanonical image of $`\overline{}_{23}`$ cannot be a curve once we construct a smooth curve of genus $`23`$ lying in the support of exactly two of the divisors $`_{23,12}^1`$,$`_{23,17}^2`$ and $`_{23,20}^3`$. In this way we rule out the possibility of all three intersections of two Brill-Noether divisors being equal to base-locus$`(|mK_{\overline{}_{23}}|)_{23}`$.
In \[Fa\] we found such genus $`23`$ curves using an intricate construction involving limit linear series (cf. Proposition 5.4 in \[Fa\]). Here we can construct such curves in a much simpler way by applying Theorem 1 when $`(d,g)=(18,23)`$: there exists a smooth curve $`C^3`$ of genus $`23`$ and degree $`18`$ such that $`\text{gon}(C)=13`$; moreover $`C`$ sits on a smooth quartic surface $`S^3`$ such that $`\text{Pic}(S)=CH.`$
Since $`C`$ has a very ample $`𝔤_{18}^3`$, by adding $`2`$ base points it will also have plenty of $`𝔤_{20}^3`$’s and also $`𝔤_{17}^2`$’s of the form $`𝔤_{18}^3(p)=\{D𝔤_{18}^3:Dp\}`$, for any $`pC`$. Therefore $`[C](_{23,20}^3_{23,17}^2)_{23,12}^1`$, and Theorem 2 now follows.
Department of Mathematics, University of Michigan
525 East University, Ann Arbor, MI 48109
e-mail: gfarkas@math.lsa.umich.edu
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# Interstellar Seeing. I. Superresolution Techniques Using Radio Scintillations
## 1 Introduction
Diffractive interstellar scintillation (DISS) is caused by multipath scattering of radio waves from small-scale density irregularities in the ionized interstellar medium. It is sensitive to intrinsic sizes of radiation sources in much the same way that optical scintillation from atmospheric turbulence is quenched for planets while strong for stars. However, interstellar scintillation differs from the atmospheric case in that it can resolve sources at angular resolutions much smaller than those achievable with available apertures, including the longest baselines used in very long baseline interferometry (VLBI), those using space antennas. Optical techniques such as intensity interferometry, speckle interferometry and adaptive optics typically only restore the telescope resolution to what it would be in the absence of any atmospheric turbulence.
We define the superresolution regime where the source is unresolved by terrestrial interferometers but is sufficiently extended to modify the DISS. Let $`\theta _s,\theta _{ij}`$ and $`\theta _{iso}`$ be the source size, interferometer fringe spacing, and isoplanatic DISS patch, respectively. By definition, two point sources separated by much less than the isoplanatic angle will show identical DISS. The isoplanatic angle $`\theta _{iso}\lambda /d\theta _d`$, where $`\lambda `$ is the wavelength, $`\theta _d`$ is the size of the scattering (“seeing”) disk and $`d`$ is the source-Earth distance. Using typical numbers ($`d=1`$ kpc, $`\theta _d=1`$ mas at an observing frequency of 1 GHz), $`\theta _{iso}0.4\mu `$arc sec. For pulsars, whose light-cylinder radii, $`r_{\mathrm{LC}}=cP/2\pi `$ ($`P`$ = spin period) are smaller than 1 $`\mu `$arc sec at typical distances, we have $`\theta _{iso}\theta _s\theta _{ij}`$ and $`\theta _{iso}\theta _d`$. In this case, speckle methods can achieve far better resolution than the interferometer. By speckle methods, we mean observations that analyze differences in the DISS between source components, which are sensitive to the spatial separations of those components. Cornwell & Narayan (1993) have discussed particular superresolution techniques in the radio context. In optical astronomy, superresolution is not achievable because $`\theta _d\theta _{iso}`$. However, the superresolution regime has been identified in optical laboratory applications (Charnotskii, Myakinin & Zavorotnyy 1990).
The purpose of this paper is to provide a rigorous and general formulation of DISS that can be used in superresolution applications. Previous work has relied on DISS theories that are restricted to particular spatial geometries (e.g. thin screens) or to particular wavenumber spectra (e.g. a power law with the “Kolmogorov” slope.). In subsequent papers, we will apply the methods of this paper to scintillation observations and derive constraints on pulsar source sizes.
Our treatment builds upon work published in both the optical and radio propagation literature and is unique in the following ways: (1) it interfaces empirical astrophysical constraints on source radiation fields with the scattering geometry and scattering strength appropriate for radio observations; (2) it calculates fluctuation statistics of realistically measureable quantities in single-aperture and interferometric observations, taking into account all sources of fluctuation; (3) we present exact calculations for the intensity probability density function (PDF) that take into account arbitrary source brightness distributions and arbitrary amounts of time-bandwidth averaging; and (4) our results can be applied to a wide range of media with arbitrary spatial extent and wavenumber spectrum. Gwinn et al. (1998) discuss issues that are very similar to those contained in this paper. Our treatment is more general than Gwinn et al.’s because it is not limited to scattering media contained in thin screens. Also, our treatment of the probability distribution is not restricted to small source sizes (compared to the isoplanatic scale). Finally, our treatment includes the effects of intrinsic source fluctuations, which are assumed negligible by Gwinn et al.
We restrict the analysis to strong (saturated) DISS. Another component — refractive interstellar scintillations (RISS) — also modulates source intensitities, but on time scales much longer than for DISS. RISS can be seen from larger angular diameter sources, by a factor of 1000, than can DISS. We are interested in modeling data spans much shorter than the characteristic RISS time scale, so we ignore RISS in our discussion. The cases we consider are what would be called the “single speckle” regime in the optical literature. This corresponds to the case where the aperture size (either a single dish diameter or an interferometer baseline) is smaller than the diffraction length scale (e.g. the Fried scale). However, our formalism can be extended easily to include aperture averaging and multiple speckle cases.
In §2 we briefly summarize previous work on use of DISS to resolve radio sources. In §3 we present a general signal model and in §4 give expressions for the modulation index of the interferometric visibility and single aperture intensity. We take into account time-bandwidth averaging and source extent. We calculate an approximate PDF of the visibility and intensity using the number of degrees of freedom in the scintillations. We also outline the exact calculation of the PDF. §5 presents a Bayesian inference method for the source size. The paper is summarized in §6. The details of our definitions and calculations are given in three Appendices. In Appendix A we derive the statistics and scintillations of the scintillating amplitude modulated noise model. In Appendix B we derive second moments for the intensity and visiblity. In Appendix C we derive the PDFs for the scintillations, total intensity, and visiblity.
Notation required in discussions of wave propagation through random media is necessarily copious. Apart from standard definitions for wavelength, frequency, speed of light, and wavenumber ($`\lambda ,\nu ,c`$ and $`k=2\pi /\lambda `$, we list in Table 1 those symbols that are used throughout the paper.
## 2 Previous Applications of DISS Superresolution
DISS superresolution techniques have been used in several ways to probe source structure. DISS has been sought from various kinds of AGNs (Condon & Backer 1975; Armstrong, Spangler & Hardee 1977; Condon & Dennison 1978; Dennison & Condon 1981), leading to bounds on the brightness temperature and limits on the Lorentz factors of bulk relativistic flow. Lovelace (1970) first suggested that DISS might resolve pulsar magnetospheres. Backer (1975) placed coarse limits on the sizes of pulsar magnetospheres by considering the fractional modulation of DISS for a few pulsars. Cordes, Weisberg & Boriakoff (1983; hereafter CWB83) placed upper bounds on the source extent, and emission altitude, for two pulsars, finding upper bounds on emission altitudes of $`0.5r_{\mathrm{LC}}`$ and 0.1$`r_{\mathrm{LC}}`$. Wolszczan & Cordes (1987) exploited a remarkable episode of multiple imaging of a pulsar by a refracting interstellar structure; they found that (a) different pulse components from a long-period pulsar showed nonidentical DISS, signifying source resolution; and (b) the implied emission altitude is comparable to $`r_{\mathrm{LC}}`$, in contrast to other estimates, based on pulse widths and polarization, which suggest emission altitudes of only 1-10% of $`r_{\mathrm{LC}}`$ (Blaskiewicz et al. 1991). Kuzmin (1992) and Smirnova, Shishov & Malofeev (1996) reported similar results on four additional long-period pulsars, again with implied emission radii $`r_{\mathrm{LC}}`$.
Recently, Gwinn et al. (1997,2000) have analyzed VLBI observations of the Vela pulsar using time and frequency resolutions that exploit the same spatial resolving power of DISS as do the single-dish observations used by others. Through estimation of the probability density function (PDF) of the visibility function magnitude, they infer that DISS shows less modulation than expected from a point source and therefore conclude that the source must be extended. They estimate a transverse size $`500`$ km for the region responsible for the pulsed flux in a narrow range of pulse phase.
In a second paper, we reassess the conclusions of Gwinn et al. (1997) by considering how time-frequency averaging in the signal processing affects inferences on source size from visibility fluctuations. Gwinn et al. (1999) have also considered time-frequency averaging and the effects of source fluctuations. The approaches differ and yield different conclusions about the importance of averaging and source noise as well as differing on estimates of the source size.
Also recently, radio observations of gamma-ray burst (GRB) afterglows suggest that DISS occurs in the early stages and then is quenched as expanding synchrotron sources are first smaller than, and then exceed, the isoplanatic scale (Goodman 1997; Frail et al. 1997). If GRB synchrotron sources are incoherent sources, the angular size requirements for DISS to appear are severe. We defer to another paper a discussion of GRB scintillation.
## 3 Signal Model for Scintillating Sources
To account for all contributions to measureable quantities, we need a comprehensive statistical model for the received signal. The model presented here includes a source that is temporally incoherent but has arbitrary spatial coherence; diffractive interstellar scintillation from an arbitrary distribution of scattering material along the line of sight; and additive radiometer noise.
We define $`\epsilon (𝐫,t)`$ to be the complex, narrowband, baseband scalar electric field that is explicitly or implicitly manipulated in radio astronomy systems. (For definitions, see Appendix A.) It is determined by the source emission mechanism and by propagation effects through intervening media, as well as by receiver and background sky noise. We consider source emission that has underlying Gaussian statistics and propagation effects from the turbulent, ionized interstellar medium (ISM). The field measured at position $`𝐫`$ is the superposition of scintillating source components and radiometer noise, $`n(𝐫,t)`$,
$`\epsilon (𝐫,t)={\displaystyle 𝑑𝐫_𝐬\epsilon _s(𝐫_𝐬,t)g(𝐫,t,\nu ,𝐫_𝐬)}+n(𝐫,t).`$ (1)
All vectors are two-dimensional and transverse to the line of sight. The corresponding intensity is
$`I(𝐫,t)=|\epsilon (𝐫,t)|^2.`$ (2)
The quantity $`\epsilon _s(𝐫,t)`$ is the field emitted per unit area at the source but whose amplitude includes implicitly an inverse distance dependence. Propagation is described by the quantity $`g(𝐫,t,\nu ,𝐫_𝐬)`$, which is the propagator for a point source at location $`𝐫_𝐬`$. It includes a phase factor for free-space propagation as well as phase and amplitude factors associated with DISS.<sup>1</sup><sup>1</sup>1 The form of Eq. 1 is approximate. Factoring the integrand relies on $`\epsilon _s`$ being narrowband with bandwidth $`\mathrm{\Delta }\nu \nu `$, that it varies much faster than the DISS propagator, $`g`$, and that the bandwidth is much smaller than the characteristic scintillation bandwidth, i.e. the characteristic bandwidth on which $`g`$ changes. This last constraint is identical to requiring that differential propagation times be much less than the shortest characteristic time scale of the signal. Later, we also consider variations of $`g`$ with frequency. The signal model presented still applies if we consider the total frequency range to comprise many separate intervals in each of which $`g`$ is piecewise constant. Here $`\nu `$ is the center frequency of the passband with bandwidth $`\mathrm{\Delta }\nu `$ that is selected by the receiver and mixed to baseband. The notation for $`\epsilon (𝐫,t)`$ and $`n(𝐫,t)`$ leaves this center frequency implicit; but each of these quantities is a time series of a narrowband process and is the baseband equivalent of the narrowband radiation field selected by the receiver. Further justification for this model is given in Appendix A.
The field emitted by the source is taken to be amplitude modulated noise (Rickett 1975; Cordes 1976b), $`\epsilon _s(𝐫_𝐬,t)=a(𝐫_𝐬,t)m(𝐫_𝐬,t)`$, which is a physically motivated and empirically confirmed model for most astrophysical sources. The amplitude modulated noise model includes nonstationary modulations $`a(𝐫_𝐬,t)`$ of stationary noise $`m(𝐫_𝐬,t)`$. The noise correlation function is $`m(𝐫_{𝐬}^{}{}_{1}{}^{},t_1)m^{}(𝐫_{𝐬}^{}{}_{2}{}^{},t_2)=\delta (𝐫_{𝐬}^{}{}_{2}{}^{}𝐫_{𝐬}^{}{}_{1}{}^{})\mathrm{\Delta }(t_2t_1)`$, where the asterisk denotes complex conjugation, angular brackets denote ensemble average, and $`\mathrm{\Delta }(\tau )`$ is a continuous delta-function-like quantity with $`\mathrm{\Delta }(0)=1`$ and width equal to the reciprocal bandwidth of the receiving system. The additive noise, $`n(𝐫,t)`$, also is ‘$`\mathrm{\Delta }`$’ correlated in time and is assumed spatially uncorrelated across nonzero baselines, $`𝐛=𝐫_2𝐫_1`$.
### 3.1 Interferometer Visibility Function & Phase Structure Function
It is well known that the mean visibility function of a scattered point source (in strong scattering) is the product of the true source visibility $`\mathrm{\Gamma }_s(𝐛)`$ and the second moment of the DISS modulation, $`\gamma _\mathrm{g}`$ (e.g. Rickett 1990 and references therein):
$`\mathrm{\Gamma }_\epsilon (𝐛,\tau _i)=\epsilon (𝐫,t)\epsilon ^{}(𝐫+𝐛,t+\tau _i)=\mathrm{\Delta }(\tau _i)\gamma _\mathrm{g}(𝐛,0,0,0)\mathrm{\Gamma }_s(𝐛).`$ (3)
The source visibility is the usual Fourier transform of the brightness distribution, $`I_s(𝐫_𝐬)=A(𝐫_𝐬)=a^2(𝐫_𝐬)`$,
$`\mathrm{\Gamma }_s(𝐛)={\displaystyle 𝑑𝐫_𝐬e^{+ikd^1𝐫_𝐬𝐛}I_s(𝐫_𝐬)},`$ (4)
where $`d`$ is the source-observer distance and $`k=2\pi c^1\nu `$. We use a spatial vector, $`𝐫_𝐬`$, to define the source brightness distribution rather than using an angular variable, as is common practice. We assume that the propagation delay between the pair of sites has already been removed, so the visibility function maximizes at $`\tau _i(\mathrm{\Delta }\nu )^1`$.
The factor $`\gamma _\mathrm{g}`$ in Eq. 3 is proportional to the second cross moment of the propagator, $`g(𝐫,t,\nu ,𝐫_𝐬)`$, at two observation positions, times and frequencies separated by $`𝐛`$, $`\tau `$, $`\delta \nu `$, respectively, and for two point sources separated by $`\delta 𝐫_𝐬`$:
$`\mathrm{\Gamma }_g(𝐛,\tau ,\delta \nu ,\delta 𝐫_𝐬)`$ $`=`$ $`g(𝐫,t,\nu ,𝐫_𝐬)g^{}(𝐫+𝐛,t+\tau ,\nu +\delta \nu ,𝐫_𝐬+\delta 𝐫_𝐬)=e^{i\psi }\gamma _\mathrm{g}(𝐛,\tau ,\delta \nu ,\delta 𝐫_𝐬).`$ (5)
In this equation, the phase $`\psi `$ is determined by free-space propagation and drops out in much of what appears below, but is responsible for the Fourier relation in Eq. 4. Appendix A shows the details. The form for $`\gamma _\mathrm{g}`$ with zero frequency lag ($`\delta \nu =0`$) is simply expressed in the Gaussian limit using the phase structure function, $`D_\varphi (𝐛,\tau ,\delta 𝐫_𝐬)`$,
$`\gamma _\mathrm{g}(𝐛,\tau ,\delta \nu =0,\delta 𝐫_𝐬)`$ $`=`$ $`e^{\frac{1}{2}D_\varphi (𝐛,\tau ,\delta 𝐫_𝐬)}.`$ (6)
For nonzero frequency lags, a closed-form expression for $`\gamma _\mathrm{g}`$ is not usually available. The DISS “gain” $`G=|g|^2`$, has unit mean, $`G=1`$, and has a normalized autocovariance in the strong scattering (Rayleigh) limit,
$`\gamma _\mathrm{G}(𝐛,\tau ,\delta \nu ,\delta 𝐫_𝐬)=G(𝐫,t,\nu ,𝐫_𝐬)G(𝐫+𝐛,t+\tau ,\nu +\delta \nu ,𝐫_𝐬+\delta 𝐫_𝐬)1=|\gamma _\mathrm{g}(𝐛,\tau ,\delta \nu ,\delta 𝐫_𝐬)|^2.`$ (7)
For a medium in which scattering occurs with variable strength all along the line of sight, the phase structure function is (Lotova & Chashei 1981; Cordes & Rickett 1998)
$`D_\varphi (𝐛,\tau ,\delta 𝐫_𝐬)`$ $``$ $`\lambda ^2{\displaystyle _0^d}𝑑sC_n^2(s)|𝐛_{\mathrm{eff}}(s)|^\alpha `$ (8)
$`𝐛_{\mathrm{eff}}(s)`$ $`=`$ $`(s/d)𝐛+𝐕_{\mathrm{eff}}(s)\tau +(1s/d)\delta 𝐫_𝐬`$ (9)
$`𝐕_{\mathrm{eff}}(s)`$ $`=`$ $`(s/d)𝐕_{\mathrm{obs}}+(1s/d)𝐕_\mathrm{p}𝐕_\mathrm{m}(s).`$ (10)
$`C_n^2`$ is the coefficient in the wavenumber spectrum for electron density variations and $`\alpha `$ is the exponent of the structure function. For a square-law structure function $`\alpha =2`$, while for a Kolmogorov medium in strong, but not superstrong scattering (Cordes & Lazio 1991), $`\alpha =5/3`$. $`𝐕_\mathrm{p}`$ is the pulsar velocity, $`𝐕_{\mathrm{obs}}`$ is the observer’s velocity and $`𝐕_\mathrm{m}`$ is the velocity of the scattering material in the ISM. For the case $`C_n^2(s)\delta (sD_s)`$, where $`D_s`$ is the distance of a scattering screen from a pulsar, we retrieve the result applicable for a thin-screen.
Note that the phase structure function $`\lambda ^2`$ for radio propagation through tenuous plasmas. For optical and infrared (IR) propagation through the atmosphere, $`D_\varphi \lambda ^2`$.
Our expressions Eq. 7-10 are quite general, being based on Gaussian statistics for the wavefield and on the saturated (Rayleigh) regime of scattering. As such, they can be used for observations of Galactic and extragalactic radio sources, including pulsars, masers, microquasars, active-galactic nuclei, and gamma-ray burst sources.
Our form for $`D_\varphi `$ applies to isotropic scattering irregularities. Evidence exists for anisotropies in heavily scattered sources (Frail et al. 1994; Wilkinson, Narayan & Spencer 1994; Yusef-Zadeh 1994; Molnar et al. 1995; Desai & Gwinn 1998; Spangler & Cordes 1998; Trotter, Moran & Rodriguez 1998). Though, for simplicity, we consider only the isotropic case in this paper, it is a simple matter to extend our results to the anisotropic case, which we will do elsewhere.
### 3.2 Isoplanatic Scales
DISS is correlated over spatial and temporal scales at the observer’s location that are determined by contours of constant $`D_\varphi `$. For a thin screen at $`s=D_s`$ from the source, we write $`D_\varphi =(|𝐛_{\mathrm{eff}}|/b_e)^\alpha `$, with $`𝐛_{\mathrm{eff}}`$ given by Eq. 9 and where $`b_e`$ is the 1/e scale of the structure function. We define the isoplanatic length scale and time scale through $`D_\varphi (\mathrm{b}_{\mathrm{iso}},0,0)=1`$ and $`D_\varphi (0,\mathrm{\Delta }t_\mathrm{d},0)=1`$, yielding $`b_{iso}=b_e(d/D_s)`$ and, if the source’s speed dominates $`\mathrm{V}_{\mathrm{eff}}`$ (as it does for many pulsars), $`\mathrm{\Delta }t_\mathrm{d}=b_e/(1D_s/d)V_\mathrm{p}`$. We also define the isoplanatic scale at the source’s location using $`D_\varphi (0,0,\delta 𝐫_{𝐬,\mathrm{𝐢𝐬𝐨}})=1`$, resulting in $`\delta r_{s,iso}=b_e(1D_s/d)^1=V_\mathrm{p}\mathrm{\Delta }t_\mathrm{d}`$. The isoplanatic scale $`\delta r_{s,iso}`$ defines the separation at which two point sources would scintillate with a correlation coefficient of $`e^1`$. This scale determines whether DISS can resolve a source, as discussed in the Introduction. For reference, the Fried scale $`r_0`$, defined in the optical and IR literature, is related to our definitions using $`D_\varphi (b,0,0)=6.88(b/r_0)^{5/3}`$ (e.g. Goodman 1985), so $`r_03.2b_e`$. Also, in the radio case $`r_0\lambda ^{6/5}`$ while $`r_0\lambda ^{+6/5}`$ for optical/IR propagation.
In strong scattering, the isoplanatic scale is smaller, in some cases by three orders of magnitude or more, than the Fresnel scale. The Fresnel scale, at meter wavelengths, $`\sqrt{\lambda D}10^{11}`$ cm for kiloparsec distances.
The isoplanatic scale $`\delta r_{s,iso}`$ is smaller for sources that are scattered more heavily. This corresponds to more distant sources, sources viewed through regions of excess scattering, or sources observed at longer wavelengths. For a continuous source, such as one with a Gaussian brightness distribution, its size compared to $`\delta r_{s,iso}`$ determines the depth of modulation of the DISS. To apply this basic idea, however, averaging over time and over the receiver bandwidth must also be dealt with carefully because it too affects the depth of modulation. We now consider all these effects in what follows.
## 4 Resolving Sources with Scintillations
Several methods can be used to exploit the scintillation phenomenon in order to resolve sources. These are (1) measurement of the fractional modulation of the source through analysis of the intensity variance; (2) estimation of the intensity PDF, equivalent to analysis of all moments; and (3) use of cross-correlation functions for the DISS of separate sources to measure the spatial offsets of those sources.
### 4.1 Visibility and Intensity Statistics
The visibility function and the intensity are both second field moments. Estimates of these second moments from finite data sets fluctuate by amounts that are formally described by the fourth field moment. Encoded in these fluctuations is information about source structure and the intervening medium. We use a normalized fourth moment — a generalized modulation index (squared) — to cast intensity and visibility fluctuations in a similar form. The modulation index includes the effects of averaging over time and frequency and is calculated for an arbitrary source brightness distribution.
### 4.2 Autocorrelation Functions
To model realistic cases, we take into account averaging over time and the finite bandwidth of the narrowband signal. For simplicity, we refer to the finite bandwidth as “frequency averaging.’ To resolve DISS, the averaging intervals $`T`$ and $`B`$ must be smaller than the characteristic correlation scales of $`G`$, the DISS gain in time and frequency. These are usually called the DISS or ‘scintillation’ time scale and bandwidth, denoted $`\mathrm{\Delta }t_\mathrm{d}`$ and $`\mathrm{\Delta }\nu _\mathrm{d}`$, respectively.
Let $`\overline{\mathrm{I}}(𝐫,t)`$ be the intensity at location $`𝐫`$ calculated for a narrowband signal with bandwidth $`B`$ after averaging over the interval $`[tT/2,t+T/2]`$. Define also the averaged visibility function as $`\overline{\mathrm{\Gamma }}(𝐛,t)`$, calculated between two sites separated by baseline $`𝐛`$ and with time lag $`\tau `$ but with zero frequency lag. The autocorrelation functions (ACFs) of the intensity and visibility are
$`R_{\overline{\mathrm{I}}}(𝐛,\overline{\tau })\overline{\mathrm{I}}(𝐫,t)\overline{\mathrm{I}}(𝐫+𝐛,t+\overline{\tau })`$ (11)
$`R_{\overline{\mathrm{\Gamma }}}(𝐛,\overline{\tau })\overline{\mathrm{\Gamma }}(𝐫,t)\overline{\mathrm{\Gamma }}^{}(𝐫+𝐛,t+\overline{\tau })`$ (12)
### 4.3 Modulation Indices
To quantify fluctuations of $`\overline{\mathrm{I}}`$ and $`\overline{\mathrm{\Gamma }}`$, we calculate the modulation indices,
$`m_{\overline{\mathrm{I}}}^2(𝐛,\tau )`$ $``$ $`{\displaystyle \frac{R_{\overline{\mathrm{I}}}(𝐛,\tau )\mathrm{I}^2}{\mathrm{I}^2}}`$ (13)
$`m_{\overline{\mathrm{\Gamma }}}^2(𝐛,\tau )`$ $``$ $`{\displaystyle \frac{R_{\overline{\mathrm{\Gamma }}}(𝐛,\tau )|\mathrm{\Gamma }_\epsilon |^2}{\mathrm{I}^2}},`$ (14)
where we normalize by the mean intensity in both cases. As shown in Appendix B, the modulation index receives contributions from four terms. For compact sources, the dominant term is in fact caused by DISS and is given by
$`m_{\mathrm{ISS}}^2(𝐛,\overline{\tau })`$ $`=`$ $`I^2{\displaystyle 𝑑𝐫_{𝐬}^{}{}_{1}{}^{}𝑑𝐫_{𝐬}^{}{}_{2}{}^{}I_s(𝐫_{𝐬}^{}{}_{1}{}^{})I_s(𝐫_{𝐬}^{}{}_{2}{}^{})Q_{\mathrm{ISS}}(𝐛,\overline{\tau },𝐫_{𝐬}^{}{}_{2}{}^{}𝐫_{𝐬}^{}{}_{1}{}^{},T,B)},`$ (15)
where, for uniform averaging<sup>2</sup><sup>2</sup>2Uniform averaging over frequency corresponds to spectrometer passbands that are perfectly rectangular; real passbands, $`h(\nu )`$, have shapes similar to Gaussian functions and can be handled by replacing $`1|\delta \nu /B|`$ in Eq. 19 with the autocorrelation of $`h(\nu )`$ and extending the limits to $`\pm \mathrm{}`$. Computations show that this refinement produces no significant changes in our discussion. in $`T`$ and $`B`$,
$`Q_{\mathrm{ISS}}(𝐛,\overline{\tau },\delta 𝐫_𝐬,T,B)`$ $`=`$ $`(TB)^1{\displaystyle _T^{+T}}𝑑\tau ^{}\left(1\left|{\displaystyle \frac{\tau ^{}}{T}}\right|\right){\displaystyle _B^{+B}}𝑑\delta \nu \left(1\left|{\displaystyle \frac{\delta \nu }{B}}\right|\right)`$ (19)
$`\times \{\begin{array}{cc}\gamma _\mathrm{G}(𝐛,\tau ^{}+\overline{\tau },\delta \nu ,\delta 𝐫_𝐬)\hfill & \text{ intensity}\hfill \\ & \\ |\mathrm{\Delta }(\tau _i)|^2e^{ikd^1𝐛\delta 𝐫_𝐬}\gamma _\mathrm{G}(0,\tau ^{}+\overline{\tau },\delta \nu ,\delta 𝐫_𝐬)\hfill & \text{ visibility},\hfill \end{array}`$
with $`\gamma _\mathrm{G}`$ defined in Eq. 7. Visibility fluctuations are independent of baseline $`𝐛`$ for unresolved sources, for which the complex exponential $`1`$, while intensity variations depend more strongly on $`𝐛`$, even for sources unresolved by the baseline. The baseline-independent property for visibility fluctuations is similar to the conclusion found by Goodman & Narayan (1989). The difference between intensity and visibility statistics arises from the ordering of the time-averaging and cross-correlation operations in the two cases.
The modulation index, as presented here, includes only DISS fluctuations. There are additional contributions to the visibility or intensity variance from intrinsic source fluctuations and from additive noise. These are secondary to our discussion here, but are important in any practical application where the time-bandwidth product is low and where intrinsic source fluctuations are high. For pulsars, pulse-to-pulse amplitude variations are important when only a few pulses are included in any averaging. Appendix B gives full expressions for all contributions to intensity variations.
### 4.4 Number of Degrees of Freedom in Fluctuations
TB averaging and extended structure expressed in the integrals of $`\gamma _\mathrm{G}`$ in Eq. 15,19 diminish scintillation fluctuations. The modulation index of the averaged intensity or visibility depends on time averaging and source extension in similar ways because both increase the number of degrees of freedom in the integrated intensity. The number of degrees of freedom is
$`N_{\mathrm{dof}}=2m_{\mathrm{ISS}}^2=2N_{\mathrm{ISS}}2,`$ (20)
where $`N_{\mathrm{ISS}}`$ is the number of independent DISS fluctuations (“scintles”) that are averaged. For observations in the single speckle regime, where scintles are resolved in time and frequency, we expect $`1N_{\mathrm{ISS}}2`$.
### 4.5 Examples
In Figures 1-3 we show $`m_{\mathrm{ISS}}^2`$ plotted against integration time $`T`$ for different values of bandwidth $`B`$. These cases are for a point source. In the figures, we use $`T`$ and $`B`$ normalized by the scintillation time scale, $`\mathrm{\Delta }t_\mathrm{d}`$, and bandwdith, $`\mathrm{\Delta }\nu _\mathrm{d}`$. Figure 1 is the case for a thin screen with a square-law structure function (i.e. $`\alpha =2`$). Figure 2 is for a thin screen with $`\alpha =5/3`$, the form appropriate for a Kolmogorov medium if the scattering is strong but not ultrastrong (Cordes & Lazio 1991). Finally, Figure 3 is the Kolmogorov case for an extended medium with uniform statistics along the line of sight.
The differences between the plotted curves for the different media are subtle, but nonetheless significant if one were to use the modulation index (for given $`T,B`$, say) to try to determine what kind of medium was responsible for a given measurement. More importantly, the differences must be considered when assessing whether the source is extended.
To consider source-size effects, we adopt a circular Gaussian brightness distribution with source size, $`\sigma _r`$,
$`I_s(𝐫_𝐬)=\left(2\pi \sigma _r^2\right)^1\mathrm{exp}\left({\displaystyle \frac{|𝐫_𝐬|^2}{2\sigma _r^2}}\right).`$ (21)
For this case, the squared modulation index is
$`m_{\mathrm{ISS}}^2(𝐛,\overline{\tau })`$ $`=`$ $`(4\pi \sigma _r^2)^1{\displaystyle 𝑑\delta 𝐫_𝐬e^{(|\delta 𝐫_𝐬|/2\sigma _r)^2}Q_{\mathrm{ISS}}(𝐛,\overline{\tau },\delta 𝐫_𝐬,T,B)}.`$ (22)
Figures 4-6 show $`m_{\mathrm{ISS}}^2`$ plotted against source size in units of the isoplanatic scale, $`\delta r_{s,iso}`$ (defined in §3.2). The three figures are for the same square-law, thin-screen Kolmogorov, and uniform Kolmogorov media considered in Figures 1-3, which are results for point sources. The different curves are for different combinations of $`T/\mathrm{\Delta }t_\mathrm{d}`$ and $`B/\mathrm{\Delta }\nu _\mathrm{d}`$. One feature to note is that $`m_{\mathrm{ISS}}^2`$ falls off more rapidly with source size $`\sigma _r/\delta r_{s,iso}`$ than it does with $`T/\mathrm{\Delta }t_\mathrm{d}`$ or with $`B/\mathrm{\Delta }\nu _\mathrm{d}`$. This is because the source-size dependence results from a two-dimensional integration over the difference vector $`𝐫_{𝐬}^{}{}_{2}{}^{}𝐫_{𝐬}^{}{}_{1}{}^{}`$ in Eq. 15 as compared with one-dimensional integrals for TB averaging. These figures show again that any inference on source size must account not only for TB-averaging, but also for the type of medium underlying the measurements. Furthermore, the predicted contributions from TB-averaging rely on accurate measurements of the DISS time scale and bandwidth.
### 4.6 Interpretation of Modulation Index
Application of Eq. 15 is as follows. If $`m_{\mathrm{ISS}}^2=1`$ (within errors) then the source is unresolved by the DISS, the baseline $`𝐛`$ has not resolved the scattering disk, and the scintillations cannot have decorrelated over the averaging intervals $`T`$ and $`B`$. The DISS gain $`G`$ then has an exponential PDF associated with the two degrees of freedom in the scattered wavefield.
Alternatively, $`m_{\mathrm{ISS}}^2<1`$ can signify (1) resolution of the scattering disk by the baseline (in the case of intensity interferometry); (2) variation of the DISS over the averaging time or averaging bandwidth; or that (3) the source has been resolved by the DISS, i.e. that it is comparable to or larger than the isoplanatic scale of the DISS. To discriminate between these possibilities, auxiliary information is needed that characterizes the dependence of $`\gamma _\mathrm{g}`$ on its four arguments, $`𝐛,\tau ,\delta \nu `$, and $`\delta 𝐫_𝐬`$. Such information is obtained by making DISS and angular broadening measurements over a wide range of frequencies (e.g. Rickett 1990).
Complications in estimating $`m_{\mathrm{ISS}}^2`$ arise from the fact that scintillating sources fluctuate, on inverse-bandwidth time scales and on a variety of longer time scales, and there is additive noise in any real-world receiver system. We consider all such complications in Appendices B and C and also in the next few sections.
### 4.7 PDF of the Averaged DISS Gain
While the modulation index of visibility fluctuations may allow inference of source structure, an analysis of the full probability density function (PDF) may be more sensitive. Here we investigate the PDF for several cases.
First consider a scintillating point source with no TB averaging of $`G`$. As is well known, the intensity PDF is a one-sided exponential in the limit of no additive noise because $`G`$ is a chi-square random variable (RV) with two degrees of freedom, $`\chi _2^2`$ (e.g. Goodman 1985). TB averaging and extended sources increase the number of degrees of freedom, and therefore decrease $`m_{\mathrm{ISS}}^2`$, as discussed in §4.4. If TB averaging and source superposition are viewed as combining statistically independent RV, $`G`$ is distributed as $`\chi _{2N_{\mathrm{ISS}}}^2`$,
$`f_G(G)`$ $``$ $`{\displaystyle \frac{(GN_{\mathrm{ISS}})^{N_{\mathrm{ISS}}}}{G\mathrm{\Gamma }(N_{\mathrm{ISS}})}}e^{GN_{\mathrm{ISS}}}U(G),`$ (23)
where $`\mathrm{\Gamma }(x)`$ is the gamma function and $`U(x)`$ is the unit step function. In detail, however, the intensity (or visibility) is the integral of variables that are statistically dependent, so the $`\chi _{2N_{\mathrm{ISS}}}^2`$ PDF is only an approximation to the true PDF of G.
The true PDF is calculated by solving a homogeneous Fredholm equation of the second kind (Press et al. 1992, pp. 779-785) that results from expanding the propagator $`g(𝐫,t,𝐫_𝐬,\nu )`$ onto an orthonormal set of eigenvectors $`\psi _n`$ (e.g. Goodman 1985, pp. 250-256) and requiring that the expansion coefficients be statistically independent (a Karhunen-Loève expansion). The general case, where time-bandwidth averaging and source extension must be considered, requires solution of the eigenvalue problem (see Appendix C)
$`(TB)^1{\displaystyle _{tT/2}^{t+T/2}}𝑑t^{}{\displaystyle _B^{+B}}𝑑\delta \nu \left(1{\displaystyle \frac{|\delta \nu |}{B}}\right){\displaystyle 𝑑𝐫_{𝐬}^{}{}_{1}{}^{}\left[I_s(𝐫_{𝐬}^{}{}_{1}{}^{})I_s(𝐫_{𝐬}^{}{}_{2}{}^{})\right]^{1/2}}`$ (24)
$`\gamma _\mathrm{g}(0,t^{\prime \prime }t^{},\delta \nu ,𝐫_{𝐬}^{}{}_{2}{}^{}𝐫_{𝐬}^{}{}_{1}{}^{})\psi _n(t^{},𝐫_{𝐬}^{}{}_{1}{}^{})=\lambda _n\psi _n(t^{},𝐫_{𝐬}^{}{}_{2}{}^{}),`$
where the $`\lambda _n`$ are the eigenvalues. The time and frequency integrands are slightly different because the wave propagator, $`g`$, is integrated over frequency before squaring of the wavefield, while time-averaging occurs after squaring.
For the simple example of a time average of $`G`$ at discrete times $`t_j=j\mathrm{\Delta }t,j=1,\mathrm{},N`$, the eigenvectors $`\psi _n(t_j)`$ and eigenvalues $`\lambda _n`$ are solutions of
$`{\displaystyle \underset{i=1}{\overset{N}{}}}\gamma _\mathrm{g}(0,t_it_j,0,0)\psi _n(t_i)=\lambda _n\psi _n(t_j),n=1,\mathrm{},N.`$ (25)
The PDF of $`G=N^1_i|g(t_i)|^2`$ is the $`N`$-fold convolution of one-sided exponential PDFs, each of which has mean $`\lambda _n/N`$, because each coefficient $`b_n`$ in the expansion of $`g(t)`$ is statistically independent and is a complex, Gaussian RV. The PDF can be written in the form $`f_G(G)=_nc_n\mathrm{exp}(GN/\lambda _n)`$ where the $`c_n`$ are functions of the eigenvalues, $`c_n=\lambda _n^1_{n^{}n}(1\lambda _n^{}/\lambda _n)^1`$. Frequency averaging behaves similarly.
Goodman (1985) has shown that exact PDFs calculated in this way are fairly well approximated by the $`\chi _{2N_{\mathrm{ISS}}}^2`$ PDF with the appropriate number of degrees of freedom given by Eq. 20. The approximate PDF has the same mean and variance as the exact PDF. In the limits $`N_{\mathrm{ISS}}1`$ and $`N_{\mathrm{ISS}}1`$, the two PDFs become identical. Also, if an observable is calculated as the sum of strictly independent DISS fluctuations with equal variances, then the exact PDF is the $`\chi _{2N_{\mathrm{ISS}}}^2`$ PDF.
Figure 7 shows exact PDFs for different instances of time averaging and finite source size. The results are for a Kolmogorov wavenumber spectrum ($`\alpha =5/3`$) but do not differ substantially for a square-law structure function. Figure 7a is a sequence of PDFs for time-averaging and for a point source, while Figure 7b is a sequence of PDFs for finite source sizes but with no time-averaging. Also shown in Figure 7c are approximate PDFs based on the $`\chi ^2`$ PDF for different numbers of degrees of freedom given by $`2N_{\mathrm{ISS}}`$, where values of $`N_{\mathrm{ISS}}`$ are chosen to yield PDFs of similar variance as in Figure 7a,b.
Comparison of the panels in Figure 7 indicates that time averaging and source extent produce similar forms for the PDF of G. This conclusion verifies the notion that, from a statistical point of view, TB-averaging and source extension produce like effects in scintillation fluctuations. The figure also supports the notion that one may use the approximate $`\chi ^2`$ PDF to make calculations rather than solving the Fredholm equation for every case. This is useful because in cases where time and frequency averaging as well as source extent are important, the Fredholm solution may not be obtainable.
### 4.8 PDF of Visibility Fluctuations
The PDF for $`G`$ derived in the previous section excludes contributions from source fluctuations such as those that arise from the amplitude modulated noise model. Here we give a nearly exact treatment that takes into account all fluctuations. For simplicity, we assume that observation baselines do not resolve either the source or the seeing disk from the scattering. In this case, we write the average visibility across baseline $`𝐛_{ij}`$ between two sites $`i`$ and $`j`$ as
$`\overline{\mathrm{\Gamma }}G\mathrm{I}+𝒩_i\delta _{ij}+X+C,`$ (26)
where $`I`$ is the mean source intensity, $`\delta _{ij}`$ is the Kronecker delta, $`𝒩`$ is the mean of $`𝒩_i|n_i|^2`$ (the background noise intensity), $`X`$ is a real Gaussian RV with zero mean, and $`C`$ is a complex Gaussian RV with zero mean. Source fluctuations are described by $`X`$ which includes noise fluctuations associated with $`m(𝐫_𝐬,t)`$ and amplitude fluctuations associated with $`a(𝐫_𝐬,t)`$ in the amplitude modulated noise model. $`C`$ includes additive radiometer noise combined with source noise fluctuations, but is uninfluenced by source amplitude fluctuations. Expressions for $`\sigma _X^2`$ and $`\sigma _C^2`$ are given in Appendix C.2.
The PDF for the visibility magnitude is calculated by successively integrating over the PDFs for the different, independent terms in Eq. 26, as done in Appendix C. The PDF for the scaled visibility magnitude, $`\gamma =|\overline{\mathrm{\Gamma }}|/\sigma _C`$, is
$`f_\gamma (\gamma )`$ $`=`$ $`{\displaystyle 𝑑Gf_G(G)𝑑Xf_X(X)\left[\gamma e^{\frac{1}{2}(\gamma ^2+G^2i^2)}I_0(\gamma i)\right]_{i=(\mathrm{I}+X/G)/\sigma _C}}`$ (27)
where $`I_0`$ is the modified Bessel function. The integrand factor in square brackets is the Rice-Nakagami PDF of a signal phasor added to complex noise (e.g. Thompson, Moran & Swenson 1991, p. 260).
To demonstrate the importance of various terms and factors, we first show, in Figure 8, the visibility PDF when we vary the bandwidth, taking into account that $`N_{\mathrm{ISS}}`$ varies as we do so. There is a tradeoff in discerning the underlying shape of the DISS gain PDF (which encourages use of narrow bandwidths) and maximizing S/N, which favors larger bandwidths: as we decrease the bandwidth there is less averaging of the scintillations but contributions from noise (the $`X`$ and $`C`$ terms) increase, thus widening the PDF.
We compare the $`X`$ and $`C`$ terms in Eq. 26 with the $`G\mathrm{I}`$ term, which dominates when the source is strong. First, we calculate the PDF for $`|\overline{\mathrm{\Gamma }}|`$ with various terms excluded. Figure 9 shows the visibility PDF when there are no DISS and no intrinsic fluctuations, just additive noise. Figure 10 shows the visibility PDF with the intrinsic fluctuations turned on for a fairly high S/N observation, but still with no DISS variations. These curves indicate that source fluctuations can contribute significantly to the shape of the PDF. Figure 11 shows the PDF for different source intensities (top panel) along with (bottom panel) the difference between the true PDF and the PDF where pulsar fluctuations are ignored. The difference vanishes when the source intensity is zero and $`1`$% for finite source intensities. The error is largest near the peak of the PDF.
Figure 12 shows the PDF when DISS is included with varying numbers of degrees of freedom, $`2N_{\mathrm{ISS}}`$, but with constant bandwidth. This case would apply to observations of sources with different intrinsic sizes or to a point source observed with varying amounts of time averaging. As $`N_{\mathrm{ISS}}\mathrm{}`$, the PDF tends toward a Gaussian form. The results indicate that the net PDF for $`|\overline{\mathrm{\Gamma }}|`$ is extremely sensitive to the number of degrees of freedom in the DISS. Also, it is evident that a pure point source can have statistics that mimic those given by an extended source if there is sufficient TB averaging.
Comparison of Figures 12 and 8 indicates that some of the changes in PDF shape evident in Figure 12 that might be due to, say, effects of source size, are indeed masked by any changes in resolution bandwidth in the signal processing.
#### 4.8.1 When is Pulsar Noise Important?
In §C.2 we give expressions for different terms in the visibility variance, including those involving pulsar fluctuations and one involving radiometer noise. All fluctuations diminish, of course, with increased averaging time. However, the relative sizes of the fluctuating terms are independent of the averaging time. In studies of the shape of the visibility or intensity PDF, pulsar fluctuations will be comparable to radiometer noise fluctuations when the pulsar signal strength satisfies
$`G\mathrm{I}\left({\displaystyle \frac{\mathrm{N}_i\mathrm{N}_j}{BW_I}}\right)^{1/2},`$ (28)
where $`\mathrm{I}`$ is the source flux density and $`W_I`$ is its characteristic time scale (either intrinsic or imposed by some sampling scheme); $`B`$ is the bandwidth. For a pulsar, $`W_I`$ would be the sample window in pulse phase and $`\mathrm{I}`$ the flux density in that window.
### 4.9 Cross-Correlation Functions
As described in Appendix B.7, when the intensity is measured for two sources and cross correlated, the total cross modulation index includes only an ISS term because noise and source fluctuations do not correlate. Defining the cross-correlation of the average intensities as in Eq. B30 and defining the cross modulation index as
$`m_{\mathrm{cross}}^2(𝐛,\overline{\tau })={\displaystyle \frac{C_{\overline{\mathrm{I}},12}(𝐛,\overline{\tau })I_1I_2}{I_1I_2}},`$ (29)
we find that $`m_{\mathrm{cross}}^2`$ is identically equal to $`m_{\mathrm{ISS}}^2`$ of Eq. 15.
The cross modulation index can be used in two ways to detect source structure or spatial offsets between two sources. First, if its maximum is smaller than unity, quenching of DISS by source structure is signified. Also, if the cross-modulation maximizes at a time lag $`\overline{\tau }0`$ for zero baseline ($`𝐛=0`$), that signifies a significant offset between sources. We show these cases in the Appendix B and will apply this method to pulsar observations in a later paper.
## 5 Inferring Source Size from Intensity & Visibility Statistics
Given measurements of the visibility of intensity that have time and frequency resolutions that resolve DISS, it is possible to place constraints on source size by comparing visibility fluctuations with those expected from the time and frequency resolutions and as a function of source structure. Here we outline general Bayesian procedures that first analyze visibility/intensity data and then are restricted to a simpler analysis of just the modulation index.
### 5.1 Likelihood Analysis of Visibility & Intensity Fluctuations
Given a set of visibility (or intensity) measurements
$`\left\{\overline{\mathrm{\Gamma }}_i(𝐛_i),i=1,N\right\}`$ (30)
we can normalize them by the rms off-source noise, $`\sigma _C`$, and use the PDF of $`\gamma \overline{\mathrm{\Gamma }}/\sigma _C`$ (c.f. Eq. 27) to calculate a likelihood function
$`={\displaystyle \underset{i}{}}f_\gamma (\gamma _i).`$ (31)
The likelihood function depends on numerous parameters that could be estimated by maximizing $``$. These include parameters that describe the source, $`𝚯_S`$, such as source structure and flux density; those that describe wave propagation through the ISM, $`𝚯_{\mathrm{𝐈𝐒𝐌}}`$, including distances, type of medium (e.g. Kolmogorov), and scintillation parameters ($`\mathrm{\Delta }\nu _\mathrm{d},\mathrm{\Delta }t_\mathrm{d}`$); and those that characterize the receiver and telescope system, $`𝚯_\mathrm{R}`$, including the telescope gain and system temperature. Many of these parameters will be known from auxiliary observations.
Denoting all parameters collectively as $`𝚯=(𝚯_\mathrm{S},𝚯_{\mathrm{ISM}},𝚯_\mathrm{R})`$, we identify the data probability $`P(𝒟|𝚯)`$ as the likelihood function and we calculate, in standard Bayesian fashion (e.g. Gregory & Loredo 1992), the posterior PDF for the parameters as
$`P(𝚯|𝒟)`$ $`=`$ $`{\displaystyle \frac{P(𝚯)P(𝒟|𝚯)}{𝑑𝚯P(𝚯)P(𝒟|𝚯)}}={\displaystyle \frac{}{𝑑𝚯}},`$ (32)
where $`P(𝚯)`$ is the prior PDF for the parameters and the denominator normalizes the PDF. The second equality follows if we assume that the parameters have a flat prior PDF. In the case where, a priori, we know many of the parameters, we adopt delta function priors and marginalize them by integrating over those parameters. For example, if we wish to derive the PDF of only source parameters, we would integrate over $`d𝚯_{\mathrm{ISM}}d𝚯_\mathrm{R}`$ to obtain $`P(𝚯_\mathrm{S}|𝒟)`$.
To apply this approach we need to solve the multidimensional Fredholm problem that includes source extent and time-bandwidth averaging if we want to use the exact PDF for scintillation gain. Alternatively, we could use the approximate PDF based on the $`\chi ^2`$ PDF (Eq. 23) by calculating the effective number of degrees of freedom associated with source extent and TB averaging. Another, simpler, approach is to analyze only the second moment of the visibility/intensity fluctuations, as we now consider.
### 5.2 Inferring Source Size from Modulation Indices
The squared modulation index $`m_{\mathrm{ISS}}^2`$ can be calculated through direct estimation of moments or by fitting a PDF shape to a histogram of visibilities (or average intensities). Such estimates of $`\widehat{m}_{\mathrm{ISS}}^2`$ are typically made from data that span a large number of scintles, $`N_{\mathrm{ISS},\mathrm{TOTAL}}1`$. This number is approximately
$`N_{\mathrm{ISS},\mathrm{TOTAL}}\left(1+\zeta {\displaystyle \frac{T_{\mathrm{tot}}}{\mathrm{\Delta }t_\mathrm{d}}}\right)\left(1+\zeta {\displaystyle \frac{B_{\mathrm{tot}}}{\mathrm{\Delta }\nu _\mathrm{d}}}\right),`$ (33)
where the total observing time and bandwidth may be written as $`T_{\mathrm{tot}}=N_TT`$ and $`B_{\mathrm{tot}}=N_\nu B`$, using $`T`$ and $`B`$ as the basic resolutions in time and frequency defined previously. The characteristic time and frequency scales for DISS are $`\mathrm{\Delta }t_\mathrm{d}`$ and $`\mathrm{\Delta }\nu _\mathrm{d}`$, respectively. The factor $`\zeta 0.20.3`$ takes into account that scintles are not packed tightly in the $`\nu t`$ plane. In Cordes (1986), I conservatively used $`\zeta =0.1`$ whereas a more accurate calculation yields the values presented here.
The fractional estimation error on $`\widehat{m}_{\mathrm{ISS}}^2`$ due to the finite number of scintles is $`\sigma _{m^2}2m_{\mathrm{ISS}}^2N_{\mathrm{ISS},\mathrm{TOTAL}}^{1/2}`$. Invoking the Central Limit Theorem for $`N_{\mathrm{ISS},\mathrm{TOTAL}}`$, we expect $`\widehat{m}_{\mathrm{ISS}}^2`$ to have PDF, $`N(\widehat{m}_{\mathrm{ISS}}^2,\sigma _{m^2})`$.
We write the likelihood function in terms of the PDF for $`\widehat{m}_{\mathrm{ISS}}^2`$ estimated from data and using the model for $`m_{\mathrm{ISS}}^2`$ (as given by Eq. 15):
$`=\left(2\pi \sigma _{\mathrm{\Gamma }^2}^{}{}_{}{}^{2}\right)^{1/2}\mathrm{exp}\left\{{\displaystyle \frac{1}{2\sigma _{m^2}}}\left[m_{\mathrm{ISS}}^2(𝚯_\mathrm{S})\widehat{m}_{\mathrm{ISS}}^2\right]^2\right\},`$ (34)
where $`𝚯_\mathrm{S}`$ is a vector of parameters that represents the source structure. Bayes’ theorem can then be applied according to Eq. 32 to derive the posterior PDF for $`𝚯_\mathrm{S}`$. For the Gaussian brightness distribution of Eq 21, the posterior PDF is simply a one-dimensional PDF $`f_{\sigma _r}(\sigma _r)`$ for the sole source parameter $`\sigma _r`$.
Uncertainties in the application of this inference scheme include the systematic errors associated with not knowing the true form of the structure function for the medium and also the statistical errors in the measured modulation index, $`\widehat{m}_{\mathrm{ISS}}^2`$ and in the DISS parameters $`\mathrm{\Delta }\nu _\mathrm{d}`$ and $`\mathrm{\Delta }t_\mathrm{d}`$. We address these uncertainties in Paper II.
## 6 Summary and Conclusions
We have derived a general methodology for analyzing diffractive interstellar scintillation fluctuations that is applicable to single aperture and interferometric observations. In this paper, we considered only the strong scattering regime where the scattered wavefield has Gaussian statistics. The method explicitly takes into account time-bandwidth averaging that is often used in the statistical analysis of such observations. Such averaging modifies the statistics in a way that is identical to the effects of extended source structure. We show that time-bandwidth averaging and extended source structure both increase the number of degrees of freedom in the scintillations from the minimum value of two that describes the fully modulated, Gaussian wavefield of the scintillations.
Our methodology can be applied to any radio source in the strong scattering regime, including compact active galactic nuclei and gamma-ray burst afterglows. In another paper, we will address sources of these types and we will also consider scintillations in the weak and transition scattering regimes.
In Paper II we apply our results to the recent VLBI observations of the Vela pulsar by Gwinn et al. (1997) and find that the scintillation statistics may be accounted for fully by time-bandwidth averaging. Any contribution from extended source structure is less than an upper limit of about 400 km at the 95% confidence interval. This upper limit on the transverse extent is substantially larger than the size expected from conventional models that place radio emission well within the light cylinder of the pulsar and close to the surface of the neutron star.
I thank Z. Arzoumanian, S. Chatterjee, C. R. Gwinn, H. Lambert, M. McLaughlin, and B. J. Rickett for useful discussions and H. Lambert and B. J. Rickett for making available their numerically-derived autocovariance functions for Kolmogorov media. This research was supported by NSF grant 9819931 to Cornell University and by NAIC, which is managed by Cornell University under a cooperative agreement with the NSF.
APPENDICES
## Appendix A Scintillating Amplitude Modulated Noise Model
Here we derive a general statistical model that incorporates incoherent summing in the source and wave propagation throughout the interstellar medium.
The narrowband (scalar) electric field incident on an aperture and selected by a feed antenna and by a bandpass receiver may be written in the form
$`E_\mathrm{\Delta }(t)=Re\left\{\epsilon (t)\mathrm{exp}(i\omega _0t)\right\},`$ (A1)
where $`\omega _0`$ is the center frequency and the complex, baseband wavefield is $`\epsilon `$ (e.g. Thomas 1969). The baseband field is often explicitly extracted through quadrature mixing schemes in heterodyned radio receivers (e.g. Thompson, Moran & Swenson 1991, p. 150).
Early work on pulsars modeled $`\epsilon (t)`$ as amplitude modulated noise (AMN) with additive background and receiver noise:
$`\epsilon (t)=a(t)m(t)+n(t),`$ (A2)
where $`m`$ and $`n`$ are complex, Gaussian wavefields that describe intrinsic source noise and additive noise, respectively. The factor $`a(t)`$ is a real modulation function that describes source variations on time scales much longer than the reciprocal center frequency or reciprocal bandwidth; otherwise the statistics of $`a(t)`$ are arbitrary. The additive noise has a modulation index $`m_N=[|m|^4/|m|^2^21]^{1/2}=1.`$
The AMN model was first presented by Rickett (1975), who attributed complex, Gaussian statistics to $`m(t)`$. For this case the modulation index is also unity, $`m_M=1`$. Cordes (1976b) considered Poissonian shot-noise statistics for $`m(t)`$ based on physical models for pulsar emission; for Poissonian noise, $`m_M1`$.
Empirical tests on pulsars (Cordes 1976a; Hankins & Boriakoff 1978; Cordes & Hankins 1979; Bartel & Hankins 1982) show consistency of $`m(t)`$ with Gaussian statistics on time scales as short as $`1\mu s`$. Tests in the time domain resort to investigation of the relative amplitudes of various terms in the autocorrelation function of the intensity (Rickett 1975; Cordes 1976a; Bartel & Hankins 1982). Tests in the frequency domain, with the same conclusion, use the autocorrelation function of the spectrum (Cordes & Hankins 1979). The model predicts that there is frequency structure in the spectrum of a single pulse with characteristic bandwidth equal to the reciprocal of the time duration of $`a(t)`$. This frequency structure averages out as multiple pulses are summed.
Tests on OH and H<sub>2</sub>O masers (Evans et al. 1972; Moran 1981) show that maser emission also conforms to the AMN picture. It is expected that any radio source can be described by AMN because large numbers of particles contribute to the observed signals and most, if not all, natural sources involve incoherent superposition of radiation from incoherent or coherent emissions from individual radiators. Thus AMN should apply to gamma-ray burst sources and the most compact AGNs that show intra-day variability.
### A.1 Amplitude Modulated Noise for Extended Sources
We model extended sources as follows. First, the baseband field produced by a point source at location $`(𝐫_𝐬,z=0)`$ is
$`\epsilon _s(𝐫_𝐬,t)=a(𝐫_𝐬,t)m(𝐫_𝐬,t),`$ (A3)
where $`a(𝐫_𝐬,t)`$ is the (real) amplitude modulation and $`m(𝐫_𝐬,t)`$ is complex Gaussian noise (Rickett 1975). Here and everywhere, vectors are two dimensional and perpendicular to the line of sight. The quantity $`\epsilon _s`$ is the field emitted per unit area at the source, uninfluenced by propagation (either through free space or a turbulent medium), except that we include the dependence on distance from the source on its mean amplitude, for simplicity. For a steady source, $`a(𝐫_𝐬,t)`$, is constant in time. For pulsars, it describes the periodic envelope of pulses that modulates the underlying noise process, $`m(𝐫_𝐬,t)`$. The corresponding mean intensity is, using $`Aa^2`$,
$`I_s(𝐫_𝐬,t)=|\epsilon _s(𝐫_𝐬,t)|^2=A(𝐫_𝐬,t).`$ (A4)
In most of this paper we assume stationary statistics, so $`I_s(𝐫_𝐬,t)I_s(𝐫_𝐬)`$.
The total measured field is the integral over source components
$`\epsilon (𝐫,t)={\displaystyle 𝑑𝐫_𝐬\epsilon _s(𝐫_𝐬,t)g(𝐫,t,\nu ,𝐫_𝐬)},`$ (A5)
where we include a multiplicative propagation factor, $`g(𝐫,t,\nu ,𝐫_𝐬)`$, defined in the next section.
The noise, with stationary statistics, has a correlation function
$`\mathrm{\Gamma }_{2m}(𝐫_{𝐬}^{}{}_{1}{}^{},𝐫_{𝐬}^{}{}_{2}{}^{},t_1,t_2)=m(𝐫_{𝐬}^{}{}_{1}{}^{},t_1)m^{}(𝐫_{𝐬}^{}{}_{2}{}^{},t_2)=\delta (𝐫_{𝐬}^{}{}_{1}{}^{}𝐫_{𝐬}^{}{}_{2}{}^{})\mathrm{\Delta }(t_2t_1),`$ (A6)
where the asterisk denotes conjugation and $`\mathrm{\Delta }(\tau )`$ is an Hermitian function having unit amplitude, $`\mathrm{\Delta }(0)=1`$, and width approximately equal to the inverse of the receiver bandwidth. Angular brackets denote an ensemble average, except where noted. The noise fourth moment is the standard dual sum of products for a complex Gaussian process,
$`\mathrm{\Gamma }_{4m}(𝐫_{𝐬}^{}{}_{1}{}^{},𝐫_{𝐬}^{}{}_{2}{}^{},𝐫_{𝐬}^{}{}_{3}{}^{},𝐫_{𝐬}^{}{}_{4}{}^{},t_1,t_2,t_3,t_4)`$ $`=`$ $`m(𝐫_{𝐬}^{}{}_{1}{}^{},t_1)m^{}(𝐫_{𝐬}^{}{}_{2}{}^{},t_2)m(𝐫_{𝐬}^{}{}_{3}{}^{},t_3)m^{}(𝐫_{𝐬}^{}{}_{4}{}^{},t_4)`$ (A7)
$`=`$ $`\mathrm{\Gamma }_{2m}(𝐫_{𝐬}^{}{}_{1}{}^{},𝐫_{𝐬}^{}{}_{2}{}^{},t_1,t_2)\mathrm{\Gamma }_{2m}(𝐫_{𝐬}^{}{}_{3}{}^{},𝐫_{𝐬}^{}{}_{4}{}^{},t_3,t_4)`$
$`+`$ $`\mathrm{\Gamma }_{2m}(𝐫_{𝐬}^{}{}_{1}{}^{},𝐫_{𝐬}^{}{}_{4}{}^{},t_1,t_4)\mathrm{\Gamma }_{2m}^{}(𝐫_{𝐬}^{}{}_{2}{}^{},𝐫_{𝐬}^{}{}_{3}{}^{},t_2,t_3).`$
### A.2 Propagation Through a Thin Diffracting Screen
Consider the following geometry: a point source at $`(𝐫_𝐬,0)`$, a thin screen at $`(𝐫^{},D_s)`$ and an observer at $`(𝐫,d=D_s+D)`$. The screen changes only the phase of incident waves. Under the narrowband approximation ($`\mathrm{\Delta }\nu \nu `$) (so that all phase factors may be considered constant over the band) the propagated baseband field for a point source at $`𝐫_𝐬`$ is (e.g. Goodman 1985)
$`\epsilon (𝐫,t,𝐫_𝐬)=(i\lambda \overline{D})^1{\displaystyle 𝑑𝐫^{}e^{i\varphi (𝐫^{},tc^1𝒟_{23})}e^{ik𝒟_{13}}\epsilon _s(𝐫_𝐬,tc^1𝒟_{13})},`$ (A8)
where $`\overline{D}(D_{s}^{}{}_{}{}^{1}+D^1)^1`$ and the integral is normalized so that a screen with zero phase yields simply a delayed version of the emitted field, $`\epsilon _s(𝐫_𝐬,tc^1d)`$. Under the paraxial approximation (transverse scales much smaller than line-of-sight distances),
$`𝒟_{13}`$ $`=`$ $`𝒟_{12}+𝒟_{23}`$
$`𝒟_{12}`$ $``$ $`D_s+{\displaystyle \frac{|𝐫^{}𝐫_𝐬|^2}{2D_s}}`$
$`𝒟_{23}`$ $``$ $`D+{\displaystyle \frac{|𝐫𝐫^{}|^2}{2D}}.`$
We assume further that variations in propagation times, $`c^1𝒟_{13}`$ and $`c^1𝒟_{23}`$ (as a function of relevant source locations $`𝐫_𝐬`$ and screen exit points $`𝐫^{}`$) are negligible compared to the characteristic variation time scales for $`\varphi (𝐫^{},t)`$ and $`\epsilon (𝐫_𝐬,t)`$. $`c^1\mathrm{\Delta }\nu (𝒟_{13}d)1`$. Typically, $`\epsilon _s`$ varies on time scales of order the reciprocal bandwidth (e.g. 100 $`\mu s`$ or less) while $`g`$ varies, due to the changing geometry, on time scales of seconds to hours or more, for the situations we wish to consider. Therefore, for $`\epsilon _s`$ to be factored out of the integral, we require $`c^1\mathrm{\Delta }\nu (𝒟_{13}d)1`$. This simply means that any time smearing from differential arrival times must be less than the time resolution of the signal. Though we assume in the remainder that the inequality is satisfied, we point out that there are many instances where it is not, corresponding to the well known ‘pulse broadening’ effect (e.g. Rickett 1990). As a rule of thumb, when pulse broadening is important, scintillations are difficult to resolve in time and frequency. And when scintillations are important, the pulse broadening can be too small to be important, as we assume here.
We can now write the propagated field for a single point source as
$`\epsilon (𝐫,t,𝐫_𝐬)=\epsilon _s(𝐫_𝐬,td/c)g(𝐫,tD/c,\nu ,𝐫_𝐬),`$ (A9)
where $`g`$ is the propagator,
$`g(𝐫,t,\nu ,𝐫_𝐬)=(i\lambda \overline{D})^1{\displaystyle 𝑑𝐫^{}\mathrm{exp}\left\{i\left[\frac{k}{2}\left(D_{s}^{}{}_{}{}^{1}|𝐫^{}𝐫_𝐬|^2+D^1|𝐫𝐫^{}|^2\right)+\varphi (𝐫^{},t)\right]\right\}},`$ (A10)
and we use $`k=2\pi c^1\nu `$. The normalization of $`g`$ yields it to be simply a unit modulus phase factor when the screen phase is zero and also it yields $`|g|^2=1`$ when the screen phase has Gaussian statistics. In the following we will ignore the delays $`d/c`$ and $`D/c`$ in Eq. A9 in our notation. The time dependence of $`g(𝐫,t,\nu ,𝐫_𝐬)`$ arises from motions of source, observer and medium, which influence all terms in the exponent in the integrand. Absent any random phase screen (i.e. $`\varphi =0`$), $`g`$ is simply a complex phase factor that describes free-space propagation,
$`g(𝐫,t,\nu ,𝐫_𝐬)=e^{ik|𝐫_𝐬𝐫|^2/2d}.`$ (A11)
### A.3 Propagator Second Moment for a Thin Screen
The propagator’s second moment across a baseline $`𝐛`$, at two times separated by $`\tau `$ for two point sources at $`𝐫_{𝐬}^{}{}_{1,2}{}^{}`$ and at two frequencies separated by $`\delta \nu `$ is
$`\mathrm{\Gamma }_g(𝐛,\tau ,\delta \nu ,\delta 𝐫_𝐬)`$ $`=`$ $`g(𝐫,t,\nu ,𝐫_𝐬)g^{}(𝐫+𝐛,t+\tau ,\nu +\delta \nu ,𝐫_𝐬+\delta 𝐫_𝐬)`$ (A12)
$`=`$ $`e^{i\psi }\gamma _\mathrm{g}(𝐛,\tau ,\delta \nu ,\delta 𝐫_𝐬)`$
$`\psi `$ $`=`$ $`kd^1\{(𝐫_𝐬𝐫)\mathrm{\Delta }𝐑+{\displaystyle \frac{1}{2}}|\mathrm{\Delta }𝐑|^2\}`$ (A13)
$`\mathrm{\Delta }𝐑`$ $`=`$ $`[\delta 𝐫_𝐬𝐛+(𝐕_\mathrm{p}𝐕_{\mathrm{obs}})\tau ].`$ (A14)
There is no term $`\mathrm{\Delta }\nu `$ in $`\psi `$ because it is negligible according to the narrowband assumption made earlier, i.e. that $`c^1\mathrm{\Delta }\nu (𝒟_{13}d)1`$.
The real, second moment $`\gamma _\mathrm{g}`$ for zero frequency lag is
$`\gamma _\mathrm{g}(𝐛,\tau ,\delta \nu =0,\delta 𝐫_𝐬)`$ $`=`$ $`e^{\frac{1}{2}D_\varphi (𝐛,\tau ,\delta 𝐫_𝐬)}`$ (A15)
where the phase structure function $`D_\varphi `$ and its arguments are given by
$`D_\varphi (𝐛,\tau ,\delta 𝐫_𝐬)`$ $``$ $`[\varphi (𝐫,t)\varphi (𝐫+𝐛_{\mathrm{eff}},t)]^2=\left({\displaystyle \frac{b_{\mathrm{eff}}}{b_e}}\right)^\alpha `$ (A16)
$`𝐛_{\mathrm{eff}}`$ $`=`$ $`(D_s/d)𝐛+𝐕_{\mathrm{eff}}\tau +(D/d)\delta 𝐫_𝐬`$ (A17)
$`𝐕_{\mathrm{eff}}`$ $`=`$ $`(D_s/d)𝐕_{\mathrm{obs}}+(1D_s/d)𝐕_\mathrm{p}𝐕_\mathrm{m}`$ (A18)
The scaling exponent for the structure function, $`D_\varphi `$, is $`\alpha `$, which takes on a value $`\alpha =5/3`$ for a Kolmogorov spectrum in some instances. The $`e^1`$ scale of $`|\gamma _\mathrm{g}|^2`$ is $`b_e`$, $`𝐕_\mathrm{p}`$ is the pulsar velocity, $`𝐕_{\mathrm{obs}}`$ is the observer’s velocity and $`𝐕_\mathrm{m}`$ is the velocity of the scattering material in the ISM. The normalized autocovariance function for $`G=|g|^2`$, which we call the intensity “gain,” is (in the strong scattering, or Rayleigh, regime)
$`\gamma _\mathrm{G}(𝐛,\tau ,\delta \nu ,\delta 𝐫_𝐬)G(𝐫,t,\nu ,𝐫_{𝐬}^{}{}_{1}{}^{})G(𝐫+𝐛,t+\tau ,\nu +\delta \nu ,𝐫_{𝐬}^{}{}_{2}{}^{})1=|\gamma _\mathrm{g}(𝐛,\tau ,\delta \nu ,\delta 𝐫_𝐬)|^2,`$ (A19)
and the mean intensity is
$`I(𝐫,t)={\displaystyle 𝑑𝐫_𝐬I_s(𝐫_𝐬,t)}={\displaystyle 𝑑𝐫_𝐬A(𝐫_𝐬,t)}.`$ (A20)
Note that $`G=1`$. We emphasize that the form for $`\gamma _\mathrm{g}`$ in the second equalities of Eq. A12 and Eq. A19 relies on the assumption that the scattered wavefield is sufficiently randomized that Gaussian statistics apply.
For nonzero frequency lags, $`\gamma _\mathrm{g}`$ generally must be obtained through appropriate numerical integration (e.g. Lambert & Rickett 1999; Lee & Jokipii 1975). For the special case of a square-law structure function, a closed form expression is available for a thin screen (Chashei & Shishov 1976; Cordes et al. 1986; Gupta et al. 1994).
### A.4 Extension to An Arbitrarily Thick Medium
The results presented so far were derived explicitly for a thin screen. For a medium in which scattering occurs with variable strength all along the line of sight, the results may be extrapolated quite simply and generally. Assuming that the net measured field is still Gaussian, it is a matter of simple geometry to work out the form of the equivalent phase structure function. Following Lotova & Chashei (1981) and Cordes & Rickett (1998) we use the same results as in Eq. A12-A26 but make the replacements,
$`D_\varphi (𝐛,\tau ,\delta 𝐫_𝐬)`$ $`=`$ $`(\lambda r_e)^2f_\alpha {\displaystyle _0^d}𝑑sC_n^2(s)|𝐛_{\mathrm{eff}}(s)|^\alpha `$ (A21)
$`𝐛_{\mathrm{eff}}(s)`$ $`=`$ $`(s/d)𝐛+𝐕_{\mathrm{eff}}(s)\tau +(1s/d)\delta 𝐫_𝐬`$ (A22)
$`𝐕_{\mathrm{eff}}(s)`$ $`=`$ $`(s/d)𝐕_{\mathrm{obs}}+(1s/d)𝐕_\mathrm{p}𝐕_\mathrm{m}(s),`$ (A23)
where
$`f_\alpha ={\displaystyle \frac{8\pi ^2}{\alpha 2^\alpha }}{\displaystyle \frac{\mathrm{\Gamma }(1\alpha /2)}{\mathrm{\Gamma }(1+\alpha /2)}}.`$ (A24)
and $`C_n^2`$ is the coefficient in the wavenumber spectrum for electron density variations (Cordes & Lazio 1991; Armstrong, Rickett & Spangler 1995). For the case $`C_n^2(s)\delta (sD_s)`$, we retrieve the thin-screen results of the previous section.
### A.5 Visibility Function
The ensemble mean visibility function is the product of the true source visibility and the propagator’s second moment (ignoring a phase factor)
$`\mathrm{\Gamma }_\epsilon (𝐛,t,\tau _i)=\epsilon (𝐫,t)\epsilon ^{}(𝐫+𝐛,t+\tau _i)=\mathrm{\Delta }(\tau _i)\gamma _\mathrm{g}(𝐛,0,0,0)\mathrm{\Gamma }_s(𝐛,t),`$ (A25)
where we have designated the interferometer lag as $`\tau _i`$. This must match any geometrical time delays to within the reciprocal of the receiver bandwidth. The source visibility is the usual Fourier transform of the brightness distribution,
$`\mathrm{\Gamma }_s(𝐛,t)={\displaystyle 𝑑𝐫_𝐬e^{+ikd^1𝐫_𝐬𝐛}I_s(𝐫_𝐬,t)}.`$ (A26)
The time lag $`\tau _i`$ in Eq. A25 is zero to within a very small time (of order the reciprocal bandwidth; see Thompson, Moran & Swenson 1991). By contrast, the time lag in Eq. A12 extends over very long times, seconds to hours, that characterize DISS fluctuations.
## Appendix B Intensity & Visibility Fluctuations
Here we consider variations in the time-averaged intensity and visibility for the model of Appendix A. We derive the modulation fractions of these quantities taking into account scintillations and intrinsic and additive noise. Our expressions will include any averaging of the scintillation modulation over frequency as well as time.
We define the general fourth moment for the narrowband field $`\epsilon `$ (sans additive noise)
$`R_{4\epsilon }(𝐛,t_1,t_2,t_3,t_4)`$ $`=`$ $`\epsilon (𝐫,t_1)\epsilon ^{}(𝐫,t_2)\epsilon (𝐫+𝐛,t_3)\epsilon ^{}(𝐫+𝐛,t_4),`$ (B1)
where $`\epsilon (𝐫,t)`$ is given by Eq. A5. Expanding out, $`R_{4\epsilon }`$ involves fourth-order moments of $`a(𝐫_𝐬,t)`$ and $`m(𝐫_𝐬,t)`$ from the AMN model of Appendix A and of the propagator $`g(𝐫,t,\nu ,𝐫_𝐬)`$. We use the fact that $`a,m`$, and $`g`$ are statistically independent. Also, $`m`$ and $`g`$ are complex gaussian processes so their fourth moments are dual sums of products of their second moments, as in Eq. A7.
The fourth moment of of $`a(𝐫_𝐬,t)`$ ends up as the second moment of its square, $`Aa^2`$, which we assume has the form
$`A(𝐫_{𝐬}^{}{}_{1}{}^{},t_1)A(𝐫_{𝐬}^{}{}_{2}{}^{},t_2)=A(𝐫_{𝐬}^{}{}_{1}{}^{},t_1)A(𝐫_{𝐬}^{}{}_{2}{}^{},t_2)[1+m_A^2\rho _A(𝐫_{𝐬}^{}{}_{2}{}^{}𝐫_{𝐬}^{}{}_{1}{}^{},t_2t_1)].`$ (B2)
This form assumes stationary statistics for $`A`$ with a correlation function $`\rho _A`$ and modulation index $`m_A`$. As we show in paper II, assuming stationary statistics is not restrictive, even for pulsars which are highly nonstationary across pulse phase but appear to have stationary statistics when a fixed pulse phase is considered. For other radio sources with intrinsic variations much longer than those of pulsars, we may consider $`A(𝐫_𝐬,t)`$ to be constant in time, with $`m_A=0`$, at least over a typical observation time of minutes.
Our treatment also includes any variations of the scintillation propagator, $`g`$, across the bandwidth of the narrowband signal, $`\epsilon (𝐫,t)`$. By partitioning $`\epsilon `$ into subbands in which the propagator is piecewise constant and between which the emitted signal $`\epsilon _s`$ is statistically independent, we easily can incorporate finite bandpass effects while being consistent with our earlier assumption about the narrowband signal.<sup>3</sup><sup>3</sup>3The net effect is that intensities from the subbands add while the intensity autocorrelation function discussed later involves an average over frequency lag.
### B.1 Autocorrelation Functions of Time Average Quantities
The time-averaged intensity,
$`\overline{\mathrm{I}}(𝐫,t)=T^1{\displaystyle _{tT/2}^{t+T/2}}𝑑t^{}I(𝐫,t^{}),`$ (B3)
has autocorrelation function
$`R_{\overline{\mathrm{I}}}(𝐛,\tau )\overline{\mathrm{I}}(𝐫,t)\overline{\mathrm{I}}(𝐫+𝐛,t+\tau )`$ $`=`$ $`T^2{\displaystyle _{tT/2}^{t+T/2}𝑑t_a𝑑t_b\mathrm{R}_\mathrm{I}(𝐛,t_bt_a+\tau )}`$ (B4)
$`=`$ $`T^1{\displaystyle _T^{+T}}𝑑\tau ^{}\left(1{\displaystyle \frac{|\tau ^{}|}{T}}\right)\mathrm{R}_\mathrm{I}(𝐛,\tau ^{}+\tau ).`$ (B5)
The integrand is given by
$`\mathrm{R}_\mathrm{I}(𝐛,\tau )=R_{4\epsilon }(𝐛,t_1,t_1,t_1+\tau ,t_1+\tau ).`$ (B6)
Similarly, we consider the visibility,
$`\mathrm{\Gamma }_\epsilon (𝐛,t,\tau _i)=\epsilon (𝐫,t)\epsilon ^{}(𝐫+𝐛,t+\tau _i)`$ (B7)
and its time average,
$`\overline{\mathrm{\Gamma }}(𝐛,t,\tau _i)=T^1{\displaystyle _{tT/2}^{t+T/2}}𝑑t^{}\overline{\mathrm{\Gamma }}(𝐛,t^{},\tau _i).`$ (B8)
The autocorrelation, analogous to Eq. B5, is
$`R_{\overline{\mathrm{\Gamma }}}(𝐛,\overline{\tau };\tau _i)=T^1{\displaystyle _T^{+T}}𝑑\tau ^{}\left(1{\displaystyle \frac{|\tau ^{}|}{T}}\right)R_{\overline{\mathrm{\Gamma }}}(𝐛,\tau ^{}+\overline{\tau };\tau _i),`$ (B9)
and involves the integrand
$`R_{\overline{\mathrm{\Gamma }}}(𝐛,\overline{\tau };\tau )=R_{4\epsilon }(𝐛,t_1,t_1+\overline{\tau },t_1+\overline{\tau }+\tau ,t_1+\tau ).`$ (B10)
Note that we distinguish here between the lag associated with the definition of the visibility, $`\tau _i`$, and the lag $`\overline{\tau }`$ with which we consider the autocorrelation of the visibility.
### B.2 Modulation Indices
We are most interested in the normalized variances of the time-average intensity and visibility. These are defined in terms of the autocorrelation functions as
$`m_{\overline{\mathrm{I}}}^2(𝐛,\tau )`$ $``$ $`{\displaystyle \frac{R_{\overline{\mathrm{I}}}(𝐛,\tau )\mathrm{I}^2}{\mathrm{I}^2}}`$ (B11)
$`m_{\overline{\mathrm{\Gamma }}}^2(𝐛,\tau )`$ $``$ $`{\displaystyle \frac{R_{\overline{\mathrm{\Gamma }}}(𝐛,\tau )|\mathrm{\Gamma }_\epsilon |^2}{\mathrm{I}^2}},`$ (B12)
where we normalize by the mean intensity in both cases.
The total modulation index squared for the intensity or visibility is the sum of three main terms:
$`m^2(𝐛,\overline{\tau })=m_{\mathrm{ISS}}^2(𝐛,\overline{\tau })+m_{\mathrm{PSR}}^2(𝐛,\overline{\tau }),+m_{\mathrm{NOISE}}^2(𝐛,\overline{\tau }),`$ (B13)
where $`m_{\mathrm{ISS}}^2`$ measures the contribution from scintillations only, $`m_{\mathrm{PSR}}^2`$ measures the contribution from source amplitude fluctuations combined with scintillations, and $`m_{\mathrm{NOISE}}^2`$ measures source noise fluctuations. Later (§B.6) we will also consider the effects of additive radiometer fluctuations. For pulsars, $`m_{\mathrm{PSR}}^2`$ includes pulse shape variations and noise fluctuations. For sources that are steady over an observation span of minutes to hours (or more), $`m_{\mathrm{PSR}}^2=0`$. The ‘noise’ term, $`m_{\mathrm{NOISE}}^2`$, depends on source structure and corresponds to the output of an intensity interferometer that is proportional to the square of the visibility function (e.g. Hanbury-Brown 1974, pp. 48-49).
We can write these terms as
$`m_{\mathrm{ISS}}^2(𝐛,\overline{\tau })`$ $`=`$ $`\mathrm{I}^2{\displaystyle 𝑑𝐱𝑑𝐲\mathrm{I}_{𝐫_𝐬}(𝐱)\mathrm{I}_{𝐫_𝐬}(𝐲)Q_{\mathrm{ISS}}(𝐛,\overline{\tau },𝐲𝐱,T,B)}`$ (B14)
$`m_{\mathrm{PSR}}^2(𝐛,\overline{\tau })`$ $`=`$ $`\mathrm{I}^2{\displaystyle 𝑑𝐱𝑑𝐲\mathrm{I}_{𝐫_𝐬}(𝐱)\mathrm{I}_{𝐫_𝐬}(𝐲)Q_{\mathrm{PSR}}(𝐛,\overline{\tau },𝐲𝐱,T,B)}`$ (B15)
$`m_{\mathrm{NOISE}}^2(𝐛,\overline{\tau })`$ $`=`$ $`\mathrm{I}^2{\displaystyle 𝑑𝐱𝑑𝐲\mathrm{I}_{𝐫_𝐬}(𝐱)\mathrm{I}_{𝐫_𝐬}(𝐲)Q_{\mathrm{NOISE}}(𝐛,\overline{\tau },𝐲𝐱,T,B)}`$ (B16)
The ‘$`Q`$’ functions are defined in terms of integrals over time lag, like that in Eq. B5 and over similar frequency-lag integrals, that we denote as
$`X(y)_{y,Y}Y^1{\displaystyle _Y^{+Y}}𝑑y\left(1{\displaystyle \frac{|y|}{Y}}\right)X(y).`$ (B17)
### B.3 Intensity Fluctuations
For intensity fluctuations, the Q functions are
$`Q_{\mathrm{ISS}}(𝐛,\overline{\tau },\delta 𝐫_𝐬,T,B)`$ $`=`$ (B18)
$`\gamma _\mathrm{G}(𝐛,\tau ^{}+\overline{\tau },\delta \nu ,\delta 𝐫_𝐬)_{\tau ^{},T;\delta \nu ,B},`$
$`Q_{\mathrm{PSR}}(𝐛,\overline{\tau },\delta 𝐫_𝐬,T,B)`$ $`=`$
$`m_A^2\{\rho _A(\delta 𝐫_𝐬,\tau ^{}+\overline{\tau })_{\tau ^{},T}+\rho _A(\delta 𝐫_𝐬,\tau ^{}+\overline{\tau })\gamma _\mathrm{G}(𝐛,\tau ^{}+\overline{\tau },\delta \nu ,\delta 𝐫_𝐬)_{\tau ^{},T;\delta \nu ,B}`$
$`+R_\mathrm{\Delta }(\overline{\tau },T)e^{ikd^1𝐛\delta 𝐫_𝐬}\rho _A(\delta 𝐫_𝐬,0)[\gamma _\mathrm{G}(𝐛,0,0,0)+\gamma _\mathrm{G}(0,0,\delta \nu ,\delta 𝐫_𝐬)_{\delta \nu ,B}]\},`$
$`Q_{\mathrm{NOISE}}(𝐛,\overline{\tau },\delta 𝐫_𝐬,T,B)`$ $`=`$ $`R_\mathrm{\Delta }(\overline{\tau },T)e^{ikd^1𝐛\delta 𝐫_𝐬}\left[\gamma _\mathrm{G}(𝐛,0,0,0)+\gamma _\mathrm{G}(0,0,\delta \nu ,\delta 𝐫_𝐬)_{\delta \nu ,B}\right].`$ (B20)
We have made use of the lag-integrated noise correlation
$`R_\mathrm{\Delta }(\tau ,T)T^1{\displaystyle _T^T}𝑑\tau ^{}\left(1{\displaystyle \frac{\tau ^{}}{T}}\right)|\mathrm{\Delta }(\tau ^{})|^2.`$ (B21)
Recall that $`\mathrm{\Delta }(\tau )`$ is a function with unit maximum amplitude \[$`\mathrm{\Delta }(0)=1`$\] and width of order the reciprocal bandwidth, $`B^1`$. For integration times $`TB^1`$, we have
$`R_\mathrm{\Delta }(\tau ,T){\displaystyle \frac{W_\mathrm{\Delta }}{T}}U(T\tau )U(T+\tau ),`$ (B22)
where $`U(x)`$ is the unit step function and $`W_\mathrm{\Delta }𝑑\tau |\mathrm{\Delta }(\tau )|^2B^1`$ is the characteristic time scale of the noise fluctuations.
Note that for $`TW_\mathrm{\Delta }`$, the terms involving $`R_\mathrm{\Delta }`$ in $`Q_{\mathrm{PSR}}`$ and $`Q_{\mathrm{NOISE}}`$ are much smaller than $`Q_{\mathrm{ISS}}`$. Also, for radio sources other than pulsars, $`Q_{\mathrm{PSR}}0`$ because $`m_A=0`$.
### B.4 Intensity Interferometry
We can relate our results to those of Hanbury-Brown and Twiss (e.g. Hanbury Brown 1974), who used the intensity autocorrelation function to determine the magnitude of the source visibility function for optical stars. The term of interest in their work corresponds to our $`Q_{\mathrm{NOISE}}`$, in particular the second term which involves $`\gamma _\mathrm{G}(0,0,\delta \nu ,\delta 𝐫_𝐬)_{\delta \nu ,B}.`$ The first term with $`\gamma _\mathrm{G}(𝐛,0,0,0)`$ vanishes for baselines much larger than the Fried scale. Also, the apertures used by Hanbury Brown and Twiss were larger than the Fried scale, causing aperture averaging that we have not treated but which is analogous to time-bandwidth averaging. For ground-based optical observations of stars, Hanbury-Brown and Twiss used a bandwidth such that scintillations were constant over the band and the stars they observed were much smaller than the isoplanatic scale. In this case, $`\gamma _\mathrm{G}(0,0,\delta \nu ,\delta 𝐫_𝐬)1`$ and the effect was maximized. Note also, however, that the amplitude of the effect scales with $`R_\mathrm{\Delta }(\overline{\tau },T)W_\mathrm{\Delta }/T(BT)^1`$, which is small for significant time-bandwidth averaging.
### B.5 Visibility Fluctuations
For the visibility we have, using $`𝚫𝐕𝐕_\mathrm{p}𝐕_{\mathrm{obs}}`$,
$`Q_{\mathrm{ISS}}(𝐛,\overline{\tau },\delta 𝐫_𝐬,T,B)`$ $`=`$ $`|\mathrm{\Delta }(\tau _i)|^2\gamma _\mathrm{G}(0,\tau ^{}+\overline{\tau },\delta \nu ,\delta 𝐫_𝐬)e^{ikd^1𝐛[\delta 𝐫_𝐬+𝚫𝐕(\tau ^{}+\overline{\tau })]}_{\tau ^{},T;\delta \nu ,B},`$ (B23)
$`Q_{\mathrm{PSR}}(𝐛,\overline{\tau },\delta 𝐫_𝐬,T,B)`$ $`=`$
$`m_A^2\{|\mathrm{\Delta }(\tau _i)|^2e^{ikd^1𝐛\delta 𝐫_𝐬}\rho _A(\delta 𝐫_𝐬,\tau ^{}+\tau )\gamma _\mathrm{G}(𝐛,0,0,0)_{\tau ^{},T}`$
$`+\rho _A(\delta 𝐫_𝐬,\tau ^{}+\tau )\gamma _\mathrm{G}(0,\tau ^{}+\overline{\tau },\delta \nu ,\delta 𝐫_𝐬)e^{ikd^1𝐛[𝚫𝐕(\tau ^{}+\overline{\tau })]}_{\tau ^{},T;\delta \nu ,B}`$
$`+R_\mathrm{\Delta }(\overline{\tau },T)\rho (\delta 𝐫_𝐬,0)[1+\gamma _\mathrm{G}(𝐛,\tau _i,\delta \nu ,\delta 𝐫_𝐬)_{\delta \nu ,B}]\}.`$
$`Q_{\mathrm{NOISE}}(𝐛,\overline{\tau },\delta 𝐫_𝐬,T,B)`$ $`=`$ $`R_\mathrm{\Delta }(\tau ,T)\left[1+\gamma _\mathrm{G}(𝐛,\tau _i,\delta \nu ,\delta 𝐫_𝐬)_{\delta \nu ,B}\right].`$ (B25)
### B.6 Effects of Additive Noise
Results given so far for intensity and visibility fluctuations have considered only the signal emitted by the source. Including the additive radiometer noise $`n`$ as in Eq. A2, we obtain an additional contribution to the total modulation index, that we denote $`m_{rad}^2`$: For the time-average intensity and visibility, the contribution is
$`m_{rad}^2(𝐛_{ij},\overline{\tau })=\{\begin{array}{cc}2\mathrm{I}^1\mathrm{N}_i\delta _{ij}R_\mathrm{\Delta }(\overline{\tau },T)\left(1+\frac{1}{2}\mathrm{I}^1\mathrm{N}_i\right)\hfill & \text{ intensity fluctuations}\hfill \\ & \\ \mathrm{I}^1\mathrm{N}_iR_\mathrm{\Delta }(\overline{\tau },T)\left(1+\frac{\mathrm{N}_j}{\mathrm{N}_i}+\mathrm{I}^1\mathrm{N}_j\right)\hfill & \text{ visibility fluctuations}\hfill \end{array}`$ (B29)
We have labelled the baseline with $`ij`$ indices to represent the $`i`$th and $`j`$th sites. The Kronecker delta indicates that the contribution for the intensity holds only for single-site measurements for which $`i`$ and $`j`$ are equal.
### B.7 Cross Correlations of Time-Average Intensities
In some circumstances, we are interested in the cross correlation function of the intensity between two sources that may or may not be scintillating together. Pulsars, for example, have different pulse components that may come from spatially different emission regions and it is possible to record or compute intensities for each component separately and cross-correlate them. We define the cross correlation as
$`C_{\overline{\mathrm{I}},12}(𝐛,\overline{\tau })\overline{\mathrm{I}}_1(𝐫,t)\overline{\mathrm{I}}_2(𝐫+𝐛,t+\tau ).`$ (B30)
The utility of the cross correlation is that it is affected by both the separations of the two sources and the size of each source. For example, the time lag at which the CCF maximizes is determined by the separation of the sources and by the effective velocity, $`𝐕_{\mathrm{eff}}`$ (Eq. 10,A18).
The cross correlation simplifies greatly if we assume that the amplitude-modulated noise in each source is statistically independent from the other. For pulsars, this is a reasonable assumption in many cases, though in others where there are drifting subpulse fluctuations that appear successively in different pulse components, this is an approximation. Letting $`I_{1,2}`$ be the (ensemble) mean intensity of each source, we find that
$`C_{\overline{\mathrm{I}}_{12}}(𝐛,\overline{\tau })`$ $``$ $`I_1(t)I_2(t+\overline{\tau })`$
$`+`$ $`{\displaystyle 𝑑𝐫_{𝐬}^{}{}_{1}{}^{}𝑑𝐫_{𝐬}^{}{}_{2}{}^{}A(𝐫_{𝐬}^{}{}_{1}{}^{},t)A(𝐫_{𝐬}^{}{}_{2}{}^{},t+\overline{\tau })\gamma _G(𝐛,\tau ^{}+\overline{\tau },\delta \nu ,𝐫_{𝐬}^{}{}_{2}{}^{}𝐫_{𝐬}^{}{}_{1}{}^{})_{\tau ^{},T;\delta \nu ,B}}.`$
For two point sources, one at $`𝐫_{𝐬}^{}{}_{1}{}^{}`$, another at $`𝐫_{𝐬}^{}{}_{2}{}^{}`$, that have stationary statistics, the normalized crosscovariance is
$`\gamma _{\overline{\mathrm{I}}_{12}}(𝐛,\overline{\tau }){\displaystyle \frac{C_{\overline{\mathrm{I}}_{12}}(𝐛,\overline{\tau })I_1I_2}{I_1I_2}}=\gamma _G(𝐛,\tau ^{}+\overline{\tau },\delta \nu ,𝐫_{𝐬}^{}{}_{2}{}^{}𝐫_{𝐬}^{}{}_{1}{}^{})_{\tau ^{},T;\delta \nu ,B}.`$ (B32)
To illustrate the utility of the crosscovariance consider the case where $`\gamma _G`$ is constant over the averaging intervals $`T`$ and $`B`$. The time lag that maximizes $`\gamma _{\overline{\mathrm{I}}_{12}}`$ is the solution of
$`{\displaystyle \frac{}{\overline{\tau }}}D_\varphi (𝐛,\tau ^{}+\overline{\tau },𝐫_{𝐬}^{}{}_{2}{}^{}𝐫_{𝐬}^{}{}_{1}{}^{})=0.`$ (B33)
For a thin screen this becomes (c.f. Eq. 6-7)
$`{\displaystyle \frac{}{\overline{\tau }}}\left|\left({\displaystyle \frac{D_s}{d}}\right)𝐛+𝐕_{\mathrm{eff}}\overline{\tau }+\left({\displaystyle \frac{D}{d}}\right)\delta 𝐫_𝐬\right|=0,`$ (B34)
which has the solution
$`\tau _{\mathrm{max}}={\displaystyle \frac{𝐕_{\mathrm{eff}}[(D_s/d)𝐛+(D/d)\delta 𝐫_𝐬]}{\mathrm{V}_{\mathrm{eff}}^{}{}_{}{}^{2}}}.`$ (B35)
The ability to estimate $`\tau _{\mathrm{max}}`$ with precision depends on its value relative to the characteristic DISS time, which is the width of $`\gamma _\mathrm{G}`$ as a function of $`\overline{\tau }`$. Defining $`\mathrm{\Delta }t_\mathrm{d}`$ using $`\gamma _\mathrm{G}(0,\mathrm{\Delta }t_\mathrm{d},0,0)=e^1`$, we find that $`\mathrm{\Delta }t_\mathrm{d}=\mathrm{}_d/\mathrm{V}_{\mathrm{eff}}`$, where $`\mathrm{}_d`$ is the characteristic diffraction scale, yielding
$`{\displaystyle \frac{\tau _{\mathrm{max}}}{\mathrm{\Delta }t_\mathrm{d}}}={\displaystyle \frac{𝐕_{\mathrm{eff}}[(D_s/d)𝐛+(D/d)\delta 𝐫_𝐬]}{\mathrm{}_d\mathrm{V}_{\mathrm{eff}}}}.`$ (B36)
For $`𝐛=0`$, we expect to identify $`\delta 𝐫_𝐬0`$ only if $`(D/d)\delta 𝐫_𝐬`$ is a sizable fraction of $`\mathrm{}_d`$ and also if the effective velocity is not orthogonal to $`\delta 𝐫_𝐬`$. The definition of ‘sizable’ depends on the number of independent ISS fluctuations used in any estimate of the cross correlation function, which is $`N_{\mathrm{ISS}}`$ given by Eq. 20. The error on $`\tau _{\mathrm{max}}`$ $`\mathrm{\Delta }t_\mathrm{d}N_{\mathrm{ISS}}^{}{}_{}{}^{1/2}`$, so a three-sigma measurement requires $`\delta 𝐫_{𝐬}^{}{}_{}{}^{}3N_{\mathrm{ISS}}^{}{}_{}{}^{1/2}(d\mathrm{}_d/D)`$, where $`𝐫_{𝐬}^{}{}_{}{}^{}𝐫_𝐬𝐕_{\mathrm{eff}}`$.
Smirnova, Shishov & Malofeev (1996) give a similar expression for $`\tau _{\mathrm{max}}`$ that is based on a medium with a square-law structure function. Our result is more general.
## Appendix C Probability Densities for Strong Scattering
Here we derive probability density functions for the DISS gain in strong scattering. We use an exact treatment based on Karhunen-Loève expansions that take into account source extent and time-bandwidth averaging. Then we derive the PDF for the visibility and intensity that takes into account all features of the amplitude modulated noise model of Appendix A.
### C.1 Exact Solution for the PDF of the Scintillation Gain
The time average intensity may be written as
$`\overline{\mathrm{I}}(𝐫,t)=T^1{\displaystyle _{t\pm T/2}}𝑑t^{}{\displaystyle 𝑑𝐫_𝐬I_s(𝐫_𝐬)|\overline{g}(r,t^{},\nu ,𝐫_𝐬)|^2}`$ (C1)
for a source with arbitrary brightness distribution $`I_s`$ modulated by DISS. The DISS modulation $`\overline{g}`$ has been bandwidth averaged in accord with the considerations of Appendix A,
$`\overline{g}(𝐫,t,\nu ,𝐫_𝐬)=B^1{\displaystyle _{\nu \pm B/2}}𝑑\nu ^{}g(𝐫,t,\nu ^{},𝐫_𝐬).`$ (C2)
Following the approach described by Goodman (1985; pp. 250-252), we expand $`I_s^{1/2}(𝐫_𝐬)\overline{g}(𝐫,t,\nu ,𝐫_𝐬)`$ onto a set of orthonormal basis vectors $`\psi _n(t,𝐫_𝐬)`$ with coefficients $`b_n`$. The orthonormality condition is
$`T^1{\displaystyle _{t\pm T/2}}𝑑t{\displaystyle 𝑑𝐫_𝐬\psi _n(t,𝐫_𝐬)\psi _n^{}^{}(t,𝐫_𝐬)}=\delta _{nn^{}}`$ (C3)
and the $`b_n`$ are given by
$`b_n=T^1{\displaystyle _{t\pm T/2}}𝑑t{\displaystyle 𝑑𝐫_𝐬I_s^{1/2}(𝐫_𝐬)\overline{g}(𝐫,t,\nu ,𝐫_𝐬)\psi _n(t,𝐫_𝐬)}.`$ (C4)
By requiring that $`b_nb_n^{}^{}=|b_n|^2\delta _{nn^{}}`$ (i.e. that the $`b_n`$ are statistically independent), the following eigenvalue problem results:
$`(TB)^1{\displaystyle _{tT/2}^{t+T/2}}𝑑t^{}{\displaystyle _B^{+B}}𝑑\delta \nu \left(1{\displaystyle \frac{|\delta \nu |}{B}}\right){\displaystyle 𝑑𝐫_{𝐬}^{}{}_{1}{}^{}\left[I_s(𝐫_{𝐬}^{}{}_{1}{}^{})I_s(𝐫_{𝐬}^{}{}_{2}{}^{})\right]^{1/2}}`$ (C5)
$`\gamma _\mathrm{g}(0,t^{\prime \prime }t^{},\delta \nu ,𝐫_{𝐬}^{}{}_{2}{}^{}𝐫_{𝐬}^{}{}_{1}{}^{})\psi _n(t^{},𝐫_{𝐬}^{}{}_{1}{}^{})=\lambda _n\psi _n(t^{},𝐫_{𝐬}^{}{}_{2}{}^{}),`$
where $`\lambda _n=|b_n|^2`$ are the eigenvalues. The time and frequency averaging are handled differently because the wave propagator, $`g`$, is integrated over frequency before squaring of the wavefield, while time-averaging occurs after squaring.
The expansion implies that
$`\overline{\mathrm{I}}(𝐫,t)={\displaystyle \underset{n}{}}|b_n|^2`$ (C6)
and that
$`I(𝐫,t)={\displaystyle \underset{n}{}}\lambda _n.`$ (C7)
The expansion coefficients are gaussian distributed because the integral Eq. C4 is a sum of gaussian variables. Therefore, each term in Eq. C6, $`|b_n|^2`$, is exponentially distributed and the intensity PDF is the convolution of each of these exponentials.
The convolution can be calculated through Fourier transforms and inverted using the residue theorem yielding, for nondegenerate eigenvalues,
$`f_I(I)={\displaystyle \underset{n=1}{\overset{N}{}}}c_ne^{I/\lambda _n}U(I),`$ (C8)
where $`U(I)`$ is the unit step function and the coefficients are given by
$`c_n=\lambda _{n}^{}{}_{}{}^{1}{\displaystyle \underset{n^{}n}{\overset{N}{}}}\left(1\lambda _n^{}/\lambda _n\right)^1.`$ (C9)
When only one eigenvalue is important, as it is for a point source with negligible time-bandwidth averaging, the PDF for I contains only a single term with mean $`\lambda =\mathrm{I}`$.
### C.2 Visibility PDF for the AMN Model
Here we present an alternative derivation of the visibility PDF that takes into account all source, propagation and additive-noise fluctuations.
The visibility function is the product of the narrowband fields from two sites (i and j)<sup>4</sup><sup>4</sup>4We take the product at the same time in order to keep notation simple. In practice, a delay must be introduced to account for the different optical path lengths to the two sites. Our notation assumes this has already been removed.,
$`\mathrm{\Gamma }(t)=\epsilon _i(t)\epsilon _j^{}(t),`$ (C10)
where $`\epsilon _{i,j}=g_{i,j}am+n_{i,j}`$ for a point source that produces an identical field at the two sites. The propagator and the additive noise are both different at the two sites, in general.
The instantaneous value of the visibility is
$`\mathrm{\Gamma }=g_ig_j^{}AM+n_in_j^{}+a(g_imn_j^{}+g_j^{}m^{}n_i).`$ (C11)
We have simplified the notation, using $`i,j`$ to label spatial location rather than using location and baseline vectors as we have used in previous sections. The first term is the scintillated pulsar signal, the second is due to additive noise at the two sites, while the third term represents cross products. If there were no scintillations, the visibility would be a noisy phasor (from the pulsar) combined with complex noise. Scintillations modify the source phasor to make it complex, in general. However, the pulsar noise (from $`A`$ and $`M`$) are in phase with respect to the scintillations.
The (ensemble-average) mean visibility is
$`\mathrm{\Gamma }=g_ig_j^{}A+N_i\delta _{ij}.`$ (C12)
In practice, a time average is used to approximate the ensemble average, with attendant errors. We use the following notation for the time-averaged visibility:
$`\overline{\mathrm{\Gamma }}=\mathrm{\Gamma }(t)_{BT}T^1{\displaystyle _{tT/2}^{t+T/2}}\mathrm{\Gamma }(t),`$ (C13)
where the subscript ‘BT’ on the angular brackets denotes time averaging of a bandlimited process with bandwidth B.
By expanding $`A`$ and $`M`$ into mean values and zero-mean fluctuations, e.g. $`A=A+\delta A`$ and $`M=M+\delta M`$, we can write the time-average visibility as
$`\overline{\mathrm{\Gamma }}=g_ig_j^{}_{BT}\mathrm{A}\mathrm{M}+\mathrm{N}_i\delta _{ij}+X+C,`$ (C14)
where $`\delta _{ij}`$ is the Kronecker delta. The first term is due to the source, the second term is the mean system noise for a single site observation ($`i=j`$), and the last two terms are fluctuations,
$`X`$ $`=`$ $`g_ig_j^{}\left[\delta A+\delta M\mathrm{A}+\delta A\delta M\right]_{BT}`$
$`C`$ $`=`$ $`a(g_imn_j^{}+g_j^{}m^{}n_i)_{BT}+n_in_j^{}_{BT}\mathrm{N}_i\delta _{ij}.`$
We separate $`X`$ and $`C`$ because in useful limiting cases, discussed in the next subsections, they become, respectively, real and complex processes. For $`i=j`$, $`C`$ also becomes real. Moreover, $`X`$ is in phase with the phasor term, $`g_ig_j^{}_{BT}\mathrm{A}\mathrm{M}`$, while $`C`$ is randomly phased.
If there are no intrinsic fluctuations, $`X`$ vanishes and $`C`$ then depends only on additive noise. However, the AMN model demands that there be source fluctuations even if there are no amplitude modulations. We let $`\mathrm{M}|m|^2=1`$ without any loss of generality.
Though we have assumed a point source to arrive at Eq. C14, the equation also applies to extended sources that are spatially incoherent. Spatial incoherence yields summation of contributions from different source elements that also imply Gaussian statistics.
All terms in $`X`$ and $`C`$ are uncorrelated, so variances of individual terms sum to yield the total variance. We assume time-bandwidth averaging such that $`BT1`$; thus $`X`$ and $`C`$ become Gaussian random variables (GRVs) by the Central Limit Theorem. However, we assume $`BT`$ is small enough so that the DISS factor, $`g_ig_j^{}_{BT}`$, is not a GRV. For now, holding $`g_i(t)`$ and $`g_j(t)`$ as fixed realizations of the DISS fluctuation (i.e. not averaging over an ensemble for these quantities), we find that $`X`$ has PDF <sup>5</sup><sup>5</sup>5$`N(0,\sigma _{X}^{}{}_{}{}^{2})`$ denotes a Gaussian PDF of a real variable with zero mean and variance $`\sigma _{X}^{}{}_{}{}^{2}`$ while $`N_c(0,\sigma _{C}^{}{}_{}{}^{2})`$ denotes a complex Gaussian quantity having real and imaginary parts with equal variances, $`\sigma _{C}^{}{}_{}{}^{2}`$., $`N(0,\sigma _{X}^{}{}_{}{}^{2})`$, while $`C`$ has PDF $`N_c(0,\sigma _{C}^{}{}_{}{}^{2})`$, where the variances are
$`\sigma _{X}^{}{}_{}{}^{2}`$ $`=`$ $`|X|^2=\mathrm{A}^2G_iG_j)_{BT}(BT)^1[BW_Am_A^2+m_M^2(1+m_A^2)]`$ (C15)
$`\sigma _{C}^{}{}_{}{}^{2}`$ $`=`$ $`|C|^2=(2BT)^1\left[\mathrm{A}\left(\mathrm{N}_iG_j_{BT}+\mathrm{N}_jG_i_{BT}\right)+\mathrm{N}_i\mathrm{N}_j\right],`$ (C16)
and we have used $`G_i|g_i|^2`$, etc. The forms of these variances are consistent with expressions given by Rickett (1975). In general, we can write the variances of the real and imaginary parts of $`X`$ and $`C`$ as $`\sigma _{X_{r,i}}^2=\frac{1}{2}\sigma _X^2(1\pm \rho _{G_{i,j}})`$ and $`\sigma _{C_{r,i}}^2=\frac{1}{2}\sigma _C^2(1\pm \delta _{i,j})`$. When the DISS is perfectly correlated between the two sites, the correlation coefficient $`\rho _{G_{i,j}}`$ (which is equal to $`\gamma _\mathrm{G}(𝐛,0,0,0)`$, c.f. Eq. A19) is unity and $`X`$ is real. As the DISS decorrelates between the sites, the $`\sigma _{X_i}\sigma _{X_r}`$. Only when the two sites are identical (e.g. for a single aperture measurement of intensity) is $`C`$ real. For all interferometers, $`C`$ is complex with equal variances of the real and imaginary parts.
Deriving the variances involves assumptions about the correlation times for the various signal terms and how they influence mean squares of the time averages. We have assumed that $`n`$ and $`m`$, the noise processes, have correlation times much smaller than that of the amplitude modulation, $`a`$, which in turn has a much smaller correlation time than the integration time used. This hierarchy is consistent with the fact that the noise correlation times are the reciprocal of the bandwidth used. For pulsars, data are often obtained by using only a small range of pulse phase but averaging over many pulse periods. Pulsar pulses are broad band but decorrelate on times about equal to the spin period. Our expression for $`\sigma _X`$ uses the correlation time $`W_A`$ for the amplitude modulation. This is effectively the width of the correlation function $`\rho _A`$ defined in Eq. B2. As applied in Paper II, we would take $`W_A`$ to be the width, $`\mathrm{\Delta }tP\mathrm{\Delta }\varphi _p`$, of the pulse window used in the analysis, where $`P`$ is the pulse phase and $`\mathrm{\Delta }\varphi _p`$ is the window width in pulse phase units (cycles). Then $`T=N_p\mathrm{\Delta }t`$. We have also included $`m_M`$, the modulation index of pulsar noise fluctuations. For amplitude modulated noise with Gaussian statistics, $`m_M1`$. By retaining it, we can see what changes in the statistics if we artifically turn off the noise fluctuations. $`\sigma _C`$ depends on $`\mathrm{M}`$ but not on $`m_M`$.
Note that the contribution of pulsar noise to visibility fluctuations relative to the contribution from additive noise is independent of the averaging time. Consider, for example, the ratio $`\sigma _X/\sigma _+`$, where $`\sigma _+`$ is the value of $`\sigma _C`$ when there is no source. We have $`\sigma _X/\sigma _+\mathrm{SNR}_0`$, where $`\mathrm{SNR}_0\mathrm{A}/\sqrt{\mathrm{N}_i\mathrm{N}_j}`$ is the ratio of source strength to system temperature, when both are in the same units (e.g. Janskys). Similarly, $`\sigma _C/\sigma _+1\sqrt{\mathrm{SNR}_0}`$. This result is at odds with Gwinn et al. (2000), who state that pulsar “self noise” can be ignored because it diminishes with averaging time. It does diminish but it cannot be ignored unless the signal to noise ratio is small.
The PDF for $`\overline{\mathrm{\Gamma }}`$ given $`g_i(t)`$ and $`g_j(t)`$ may be calculated by appropriate integration over the Gaussian PDFs for $`X`$ and $`C`$. We now consider some specific cases.
### C.3 DISS Perfectly Correlated Between Sites and Constant over $`BT`$
For perfectly correlated DISS (between sites $`i`$ and $`j`$), $`g_ig_j^{}_{BT}G_i_{BT}=G_j_{BT}G_{BT}`$ and $`G_iG_j)_{BT}G^2_{BT}`$; thus $`X`$ becomes real. If, moreover, the DISS modulation is constant over the averaging time $`T`$, then $`G_{BT}G=`$ constant and $`G^2_{BT}G^2`$. The visibility becomes
$`\overline{\mathrm{\Gamma }}=G\mathrm{A}+\mathrm{N}_i\delta _{ij}+X+C.`$ (C17)
We assume a point source (and strong, saturated scintillations) so that, with $`G`$ constant over $`T`$, the scintillation PDF is a one-sided exponential. The PDFs of individual elements of $`\overline{\mathrm{\Gamma }}`$ are
$`f_G(G)`$ $`=`$ $`e^GU(G)`$
$`f_X(X)`$ $`=`$ $`N(0,\sigma _{X}^{}{}_{}{}^{2})`$
$`f_C(C)`$ $`=`$ $`N_c(0,\sigma _{C}^{}{}_{}{}^{2})`$
$`\sigma _{X}^{}{}_{}{}^{2}`$ $`=`$ $`G^2\mathrm{A}^2(BT)^1\left[BW_Am_A^2+m_M^2(1+m_A^2)\right]`$
$`\sigma _{C}^{}{}_{}{}^{2}`$ $`=`$ $`(2BT)^1[(\mathrm{N}_i\mathrm{N}_j)++G\mathrm{A}(\mathrm{N}_i+\mathrm{N}_j)].`$ (C18)
#### C.3.1 PDF of the Complex Visibility
Let $`\overline{\mathrm{\Gamma }}=\overline{\mathrm{\Gamma }}_r+i\overline{\mathrm{\Gamma }}_i`$. The PDF of $`\overline{\mathrm{\Gamma }}`$ is
$`f_{\overline{\mathrm{\Gamma }}}(\overline{\mathrm{\Gamma }})`$ $`=`$ $`{\displaystyle 𝑑Gf_G(G)𝑑Xf_X(X)f_C(\overline{\mathrm{\Gamma }}_rG\mathrm{A}X,\overline{\mathrm{\Gamma }}_i)}`$ (C19)
$`=`$ $`{\displaystyle 𝑑Gf_G(G)𝑑X(2\pi \sigma _X)^{1/2}e^{X^2/2\sigma _X^2}(2\pi \sigma _C)^1e^{{\displaystyle \frac{1}{2\sigma _C^2}}\left[\left(\overline{\mathrm{\Gamma }}_rG\mathrm{A}X\right)^2+\overline{\mathrm{\Gamma }}_i^2\right]}}.`$
Performing the integral over $`X`$, we obtain
$`f_{\overline{\mathrm{\Gamma }}}(\overline{\mathrm{\Gamma }})`$ $`=`$ $`(2\pi \sigma _C)^1e^{\overline{\mathrm{\Gamma }}_i^2/2\sigma _C^2}{\displaystyle 𝑑Gf_G(G)(\sigma _X^2+\sigma _C^2)^{1/2}e^{{\displaystyle \frac{1}{2}}\left(\overline{\mathrm{\Gamma }}_rG\mathrm{A}\right)^2/2\left(\sigma _X^2+\sigma _C^2\right)}}.`$ (C20)
Note that for no signal ($`\mathrm{A}0`$), we get
$`f_{\overline{\mathrm{\Gamma }}}(\overline{\mathrm{\Gamma }})`$ $`=`$ $`(2\pi \sigma _C^2)^1e^{|\overline{\mathrm{\Gamma }}|^2/2\sigma _C^2},`$ (C21)
a circular Gaussian PDF. Our expression in Eq. C20 disagrees with Eq. 11 of Gwinn et al. (2000), which assigns equal variances to the real and imaginary parts of $`\overline{\mathrm{\Gamma }}`$. The variances are not equal, in general. Also, there is an extra factor of $`2\pi `$ in their equation.
#### C.3.2 PDF of the Visibility Magnitude
The PDF of $`|\overline{\mathrm{\Gamma }}|`$ can be calculated as
$`f_{|\overline{\mathrm{\Gamma }}|}(|\overline{\mathrm{\Gamma }}|)=|\overline{\mathrm{\Gamma }}|{\displaystyle _0^{2\pi }}𝑑\varphi f_{\overline{\mathrm{\Gamma }}}(|\overline{\mathrm{\Gamma }}|\mathrm{cos}\varphi ,|\overline{\mathrm{\Gamma }}|\mathrm{sin}\varphi ).`$ (C22)
Here we take a slightly different approach. It is convenient to scale the magnitude of the visibility by the rms of the complex term, $`\sigma _C`$. Using $`\gamma |\overline{\mathrm{\Gamma }}|/\sigma _C`$ and $`i\mathrm{A}/\sigma _C`$, the conditional PDF for constant $`G`$ and $`X=0`$ is
$`f_\gamma (\gamma |i,G)=\gamma e^{\frac{1}{2}(\gamma ^2+G^2i^2)}I_0(\gamma i),`$ (C23)
where $`I_0`$ is the modified Bessel function. This result is the well known Rice-Nakagami PDF for a real phasor of length $`i`$ combined with a complex Gaussian phasor (e.g. Thomson, Moran & Swenson 1991, Eq. 9.37). Integrating over the PDF for X, we have, for fixed $`G`$,
$`f_\gamma (\gamma |G)={\displaystyle 𝑑Xf_X(X)f_\gamma (\gamma |i+\frac{X}{G\sigma _C},G)}.`$ (C24)
Then, integrating over the PDF for $`G`$, we have the PDF for $`\gamma `$ that takes into account all fluctuations, including DISS,
$`f_\gamma (\gamma )={\displaystyle 𝑑Gf_G(G)f_\gamma (\gamma |G)}.`$ (C25)
For S/N $`\mathrm{}`$, $`\sigma _C0`$,
$`\overline{\mathrm{\Gamma }}`$ $`=`$ $`G\mathrm{A}+\mathrm{N}_i\delta _{ij}+X,`$ (C26)
and the PDF of $`|\overline{\mathrm{\Gamma }}|`$ for fixed $`G`$ becomes $`N(G\mathrm{A}+\mathrm{N}_i\delta _{ij},\sigma _{X}^{}{}_{}{}^{2})`$, with the PDF for $`\gamma `$ given by Eq. C25.
### C.4 Perfectly Correlated DISS but $`G`$ Constant over BT
Specializing to the case of a weak source for which $`\sigma _X\sigma _C`$, we have (using $`GG_{BT}`$)
$`\overline{\mathrm{\Gamma }}`$ $``$ $`G\mathrm{A}+\mathrm{N}_i\delta _{ij}+C`$
$`f_G(G)`$ $``$ $`{\displaystyle \frac{(GN_{\mathrm{ISS}})^{N_{\mathrm{ISS}}}}{G\mathrm{\Gamma }(N_{\mathrm{ISS}})}}e^{GN_{\mathrm{ISS}}}U(G),`$
where $`\mathrm{\Gamma }(x)`$ is the gamma function and $`U(x)`$ is the unit step function. $`G`$ is distributed approximately as $`\chi _{2N_{\mathrm{ISS}}}^2`$, a chi-square random variable with $`2N_{\mathrm{ISS}}`$ degrees of freedom, where $`N_{\mathrm{ISS}}m_{\mathrm{ISS}}^2`$ and where $`m_{\mathrm{ISS}}^2`$ is given by Eq. 20. The true PDF of $`G`$ is obtained by solving the appropriate Fredholm equation for the eigenvalues that determine the PDF, as described in the main text.
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# Topological gravity in genus 2 with two primary fields
## 1. Topological recursion relations
### 1.1. Notation
The correlators of the theory are denoted $`\tau _{k_1,a_1}\mathrm{}\tau _{k_n,a_n}_g`$. We denote $`\tau _{0,a}`$ by $`𝒪_a`$. The labels on the primaries are fixed in such a way that the puncture operator is $`𝒪_0`$. Let $`\eta _{ab}`$ be the intersection form, $`\eta ^{ab}`$ its inverse, and let $`𝒪^a=\eta ^{ab}𝒪_b`$. In the case of two primaries, the intersection form equals $`\eta _{ab}=\delta _{a+b,1}`$.
Let $`_g`$ be the genus $`g`$ potential on the large phase space:
(1.1)
$$_g=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{\begin{array}{c}k_1\mathrm{}k_n\\ a_1\mathrm{}a_n\end{array}}{}t_{k_1}^{a_1}\mathrm{}t_{k_n}^{a_n}\tau _{k_1,a_1}\mathrm{}\tau _{k_n,a_n}_g.$$
we use the summation convention with respect to the indices $`a_i`$ labelling the primaries.
Denote $`/t_k^a`$ by $`_{k,a}`$. The vector field $`=_{0,0}`$, corresponding to the puncture operator $`𝒪_0`$, plays a special role in the theory. The partial derivatives of the potential $`_g`$ are denoted
$$\tau _{k_1,a_1}\mathrm{}\tau _{k_n,a_n}_g=_{k_1,a_1}\mathrm{}_{k_n,a_n}_g.$$
### 1.2. The topological recursion relation in genus $`0`$
The simplest example of a topological recursion relation is obtained by taking the relation $`\psi _1=0`$ on the zero-dimensional moduli space $`\overline{}_{0,3}`$. The resulting topological recursion relation is the equation
(1.2)
$$\tau _{k,a}\tau _{\mathrm{},b}\tau _{m,c}_0=\tau _{k1,a}𝒪^d_0𝒪_d\tau _{\mathrm{},b}\tau _{m,c}_0.$$
Let $`\mathrm{\Theta }`$ be the power series
$$\mathrm{\Theta }(z)_a^b=\delta _a^b+\underset{k=0}{\overset{\mathrm{}}{}}z^{k+1}\tau _{k,a}𝒪^b_0;$$
it is an orthogonal matrix, in the sense that $`\mathrm{\Theta }^1(z)=\mathrm{\Theta }^{}(z)`$. Let $`𝒰`$ be the matrix with components $`𝒰_a^b=𝒪_a𝒪^b_0`$. The topological recursion relation (1.2) with $`m=0`$ may be rewritten as
(1.3)
$$_{k,a}\mathrm{\Theta }(z)=z\mathrm{\Theta }(z)_{k,a}𝒰.$$
Let $`_a(z)=_{k=0}^{\mathrm{}}z^k_{k,a}`$, and define vector fields $`\{D_{k,a}k0\}`$ on the large phase space by
$$D_a(z)=\underset{k=0}{\overset{\mathrm{}}{}}z^kD_{k,a}=\mathrm{\Theta }^1(z)_a^b_b(z).$$
For example, $`D_{0,a}=_{0,a}`$ and $`D_{1,a}=_{1,a}𝒰_a^b_{0,b}`$.
###### Lemma 1.1.
We have $`D_a(z)𝒰=_{0,a}𝒰`$ and $`D_a(z)\mathrm{\Theta }(w)=w\mathrm{\Theta }(w)_{0,a}𝒰`$. In particular, $`D_{k,a}𝒰=D_{k,a}\mathrm{\Theta }=0`$ if $`k>0`$.
###### Proof.
It follows easily from (1.2) that $`D_a(z)𝒰=_{0,a}𝒰`$; since
$$D_a(z)\mathrm{\Theta }(w)=w\mathrm{\Theta }(w)D_a(z)𝒰,$$
the result follows. ∎
###### Corollary 1.2.
The vector fields $`D_{k,a}`$ and $`D_{\mathrm{},b}`$ commute if both $`k`$ and $`\mathrm{}`$ are positive, while
$$[D_{k,a},_{0,b}]=𝒪_a𝒪_b𝒪^c_0D_{k1,c}.$$
###### Proof.
By Lemma 1.1,
$`D_a(w)D_b(z)`$ $`=D_a(w)\mathrm{\Theta }^1(z)_b^c_c(z)`$
$`=\mathrm{\Theta }^1(z)_b^cD_a(w)_c(z)z𝒪_a𝒪_b𝒪^c_0D_c(z)`$
$`=\mathrm{\Theta }^1(z)_b^c\mathrm{\Theta }^1(w)_a^d_d(w)_c(z)z𝒪_a𝒪_b𝒪^c_0D_c(z).`$
It follows that $`[D_a(w),D_b(z)]=𝒪_a𝒪_b𝒪^c_0(wD_c(w)zD_c(z))`$. ∎
This corollary leads to an algorithm for the calculation of $`D_a(z)𝒪_{a_1}\mathrm{}𝒪_{a_n}_g`$ by induction on $`n`$ in terms of $`D_a(z)_g`$, using the formula
(1.4)
$$D_a(z)𝒪_{a_1}\mathrm{}𝒪_{a_n}_g=\underset{i=1}{\overset{n}{}}_{0,a_1}\mathrm{}[D_a(z),_{0,a_i}]\mathrm{}_{0,a_n}_g+_{0,a_1}\mathrm{}_{0,a_n}D_a(z)_g.$$
### 1.3. The string equation in genus $`0`$ and coordinates on the large phase space
The genus $`0`$ string equation says that $`_1_0+\frac{1}{2}\eta _{ab}t_0^at_0^b=0`$, where $`_1`$ is the vector field
$$_1=\underset{k=0}{\overset{\mathrm{}}{}}t_{k+1}^a_{k,a}_{0,0}.$$
The string equation implies the following lemma.
###### Lemma 1.3.
The restriction of $`𝒰`$ to the small phase space $`\{t_k^a=0k>0\}`$ equals the identity, while for $`n>1`$, the restriction of $`^n𝒰`$ to the small phase space vanishes.
###### Proof.
The vector fields $`_{0,a}`$ commute with $`_1`$; it follows that
$$_1𝒰_{ab}=_1_{0,a}_{0,b}_0=_{0,a}_{0,b}_1_0=\eta _{ab}.$$
Written out explicitly, this equation says that
$$𝒰_a^b=\delta _a^b+\underset{k=0}{\overset{\mathrm{}}{}}t_{k+1}^c𝒪_a𝒪^b\tau _{k,c}_0.$$
Applying the operator $`^{n1}`$, $`n>0`$, we obtain
$$^n𝒰_a^b=\underset{k=0}{\overset{\mathrm{}}{}}t_{k+1}^c^{n1}𝒪_a𝒪^b\tau _{k,c}_0.$$
The lemma is an immediate consequence of these formulas. ∎
In conjunction with the genus $`0`$ topological recursion relation, this implies the following theorem.
###### Theorem 1.4.
Let $`u^a=𝒪^a_0`$. The functions $`u_n^a=^nu^a`$, $`n0`$, form a coordinate system in a neighbourhood of the small phase space, and
(1.5)
$$D_a(z)=\underset{n=0}{\overset{\mathrm{}}{}}((+z𝒰)^n𝒰)_a^b\frac{}{u_n^b}.$$
###### Proof.
Since $`u^b=𝒰_0^b`$, Lemma 1.1 implies that
$$D_a(z)u_n^b=\left(\mathrm{\Theta }^1(z)^n\mathrm{\Theta }(z)𝒰\right)_a^b=\left((\mathrm{\Theta }^1(z)\mathrm{\Theta }(z))^n𝒰\right)_a^b.$$
Since $`\mathrm{\Theta }^1(z)\mathrm{\Theta }(z)=+z𝒰`$ by (1.3), we conclude that $`D_a(z)u_n^b=((+z𝒰)^n𝒰)_a^b`$.
By Lemma 1.3, the restriction of $`(+z𝒰)^n𝒰`$ to the small phase space equals $`z^n`$. It follows that the restriction of $`D_{k,a}u_n^b`$ to the small phase space equals $`\delta _{k,n}\delta _a^b`$; hence the functions $`u_n^a`$ form a coordinate system in a neighbourhood of the small phase space. ∎
Note that $`(+z𝒰)^n𝒰=z^1p_{n+1}(z𝒰)`$, where $`p_{n+1}(f)=(+f)^nf`$ is the $`(n+1)`$st Faà di Bruno polynomial.
###### Corollary 1.5.
If $`D_{k,a}f=0`$ for $`k>n`$, then $`f/u_k^a=0`$ for $`k>n`$.
###### Corollary 1.6.
In terms of the coordinates $`u_n^a`$, the small phase space $`\{t_k^a=0k>0\}`$ is the submanifold
$$u_n^a=\{\begin{array}{cc}\delta _0^a\hfill & n=1,\hfill \\ 0\hfill & n>1.\hfill \end{array}$$
Theorem 1.4 shows that the large phase space may be defined for any Frobenius manifold $`M`$, as the infinite jet space $`J^{\mathrm{}}M`$ (i.e. Dubrovin’s “loop space”). This is seen by rewriting the matrix $`𝒰_b^a`$ as $`u^c𝒜_{bc}^a`$, where
(1.6)
$$𝒜_{bc}^a=\frac{𝒰_b^a}{u^c}$$
is the tensor describing the product on the tangent bundle of $`M`$.
An attractive feature of the vector fields $`D_{k,a}`$ is that they commute with $`_1`$:
$`[_1,D_a(z)]`$ $`=[_1,\mathrm{\Theta }^1(z)_a^b_b(z)]=[_1,\mathrm{\Theta }^1(z)_a^b]_b(z)+\mathrm{\Theta }^1(z)_a^b[_1,_b(z)]`$
$`=(z\mathrm{\Theta }^1(z)_a^b)_b(z)\mathrm{\Theta }^1(z)_a^b(z_b(z))=0.`$
By the genus $`0`$ string equation, $`_1u_n^a`$ vanishes for $`n>0`$, while $`_1u^a=\delta _0^a`$: it follows that in the coordinate system $`\{u_n^a\}`$, the vector field $`_1`$ is given by the formula
$$_1=\frac{}{u^0}.$$
In the coordinate system $`\{u_n^a\}`$, the string equation $`_1_g=0`$ says that $`_g`$ is independent of $`u^0`$.
Lemma 1.3 shows that $`𝒰`$ is invertible in a neighbourhood of the small phase space: denote its inverse by $`𝒞`$. We also see that its determinant $`\mathrm{\Delta }=det(𝒰)`$ equals $`1`$ on the small phase space.
### 1.4. The topological recursion relation in genus $`1`$
We now illustrate the way in which use of the vector fields $`D_{k,a}`$ simplifies the discussion of topological recursion relations, using as an example the topological recursion relation in genus $`1`$:
(1.7)
$$\tau _{k,a}_1=\tau _{k1,a}𝒪^b_0𝒪_b_1+\frac{1}{24}\tau _{k1,a}𝒪_b𝒪^b_0.$$
Multiplying by $`z^k`$ and summing over $`k`$, we obtain
$$_a(z)_1=\mathrm{\Theta }(z)_a^b𝒪_b_1+\frac{1}{24}z_a(z)\mathrm{Tr}(𝒰),$$
hence, by Lemma 1.1,
$$D_a(z)_1=𝒪_b_1+\frac{1}{24}zD_a(z)\mathrm{Tr}(𝒰)=𝒪_b_1+\frac{1}{24}z_{0,a}\mathrm{Tr}(𝒰).$$
This may be written as the sequence of differential equations
(1.8)
$$D_{k,a}_1=\{\begin{array}{cc}\frac{1}{24}_{0,a}\mathrm{Tr}(𝒰)\hfill & k=1,\hfill \\ 0\hfill & k>1.\hfill \end{array}$$
The equations (1.8) have the particular solution $`\frac{1}{24}\mathrm{log}(\mathrm{\Delta })`$. Let $`\psi =_1\frac{1}{24}\mathrm{log}(\mathrm{\Delta })`$; we see that $`D_{k,a}\psi =0`$ for all $`k>0`$. Hence, by Corollary 1.5, $`\psi `$ depends only on the coordinates $`u^a`$; by the string equation, it is independent of $`u^0`$. In this way, we recover a result of Dijkgraaf and Witten : there is a function $`\psi `$ of the coordinates $`\{u^a\}`$ such that $`_1=\frac{1}{24}\mathrm{log}(\mathrm{\Delta })+\psi `$.
### 1.5. The dilaton equation
The dilaton equation is another important constraint on the potentials of topological gravity. Let $`𝒟`$ be the vector field
$$𝒟=_{1,0}\underset{k=0}{\overset{\mathrm{}}{}}t_k^a_{k,a}.$$
The dilaton equation says that
$$𝒟_g=\{\begin{array}{cc}(2g2)_g,\hfill & g1,\hfill \\ \chi /24,\hfill & g=1,\hfill \end{array}$$
where $`\chi `$ is the Euler characteristic of the background.
###### Proposition 1.7.
In the coordinate system $`\{u_n^a\}`$, the dilaton vector field $`𝒟`$ equals
$$𝒟=\underset{n=1}{\overset{\mathrm{}}{}}nu_n^a\frac{}{u_n^a}.$$
###### Proof.
By the genus $`0`$ dilaton equation $`𝒟_0=2_0`$, we have $`𝒟u_n^a=nu_n^a`$, and the formula for $`𝒟`$ follows. ∎
## 2. The $`A_2`$ and $`^1`$ models in genus $`2`$
In genus $`2`$, there are two topological recursion relations . The first is
(2.1) $`\tau _{k,a}_2`$ $`=\tau _{k1,a}𝒪^b_0𝒪_b_2+\tau _{k2,a}𝒪^b_0\left(\tau _{1,b}_2𝒪_b𝒪^c_0𝒪_c_2\right)`$
$`+\tau _{k2,a}𝒪^b𝒪^c_0\left(\frac{7}{10}𝒪_b_1𝒪_c_1+\frac{1}{10}𝒪_b𝒪_c_1\right)`$
$`+\frac{13}{240}\tau _{k2,a}𝒪^b𝒪^c𝒪_c_0𝒪_b_1\frac{1}{240}\tau _{k2,a}𝒪^b_1𝒪_b𝒪^c𝒪_c_0`$
$`+\frac{1}{960}\tau _{k2,a}𝒪^b𝒪_b𝒪^c𝒪_c_0.`$
Using the topological recursion relations in genus $`0`$ and $`1`$, (2.1) may be rewritten as the sequence of differential equations
(2.2)
$$D_{k,a}_2=_{k,a},$$
where
$$_{k,a}=\{\begin{array}{cc}𝒪_a𝒪_b𝒪_c_0\left(\frac{7}{10}𝒪^b_1𝒪^c_1+\frac{1}{10}𝒪^b𝒪^c_1\right)\hfill & \\ +\frac{13}{240}𝒪_a𝒪_b𝒪_c𝒪^c_0𝒪^b_1\frac{1}{240}𝒪_a𝒪^b_1𝒪_b𝒪_c𝒪^c_0\hfill & k=2,\hfill \\ +\frac{1}{960}𝒪_a𝒪_b𝒪^b𝒪_c𝒪^c_0\hfill & \\ 𝒪_a𝒪_b𝒪_c_0\left(\frac{1}{20}𝒪^b_1𝒪^c𝒪^d𝒪_d_0+\frac{1}{480}𝒪^b𝒪^c𝒪^d𝒪_d_0\right)\hfill & k=3,\hfill \\ +\frac{1}{1152}𝒪_a𝒪^b𝒪^c𝒪_c_0𝒪_b𝒪^d𝒪_d_0\hfill & \\ \frac{1}{1152}𝒪_a𝒪^b𝒪^c_0𝒪_b𝒪_c𝒪^d_0𝒪_d𝒪^e𝒪_e_0\hfill & k=4,\hfill \\ 0\hfill & k>4.\hfill \end{array}$$
The other topological recursion relation in genus $`2`$ is,
(2.3)
$$\begin{array}{c}\tau _{k,a}\tau _{\mathrm{},b}_2=\tau _{k,a}𝒪_c_2𝒪^c\tau _{\mathrm{}1,b}_0+\tau _{k1,a}𝒪_c_0𝒪^c\tau _{\mathrm{},b}_2\hfill \\ \hfill \begin{array}{cc}& \tau _{k1,a}𝒪_c_0\tau _{\mathrm{}1,b}𝒪_d_0𝒪^c𝒪^d_2\hfill \\ & +3\tau _{k1,a}\tau _{\mathrm{}1,b}𝒪^c_0\left(\tau _{1,c}_2𝒪_c𝒪^d_0𝒪_d_2\right)\hfill \\ & +\frac{13}{10}\tau _{k1,a}\tau _{\mathrm{}1,b}𝒪_c𝒪_d_0𝒪^c_1𝒪^d_1\hfill \\ & +\frac{4}{5}\left(\tau _{k1,a}𝒪_c_1𝒪_d_1+\frac{1}{24}\tau _{k1,a}𝒪_c𝒪_d_1\right)\tau _{\mathrm{}1,b}𝒪^c𝒪^d_0\hfill \\ & +\frac{4}{5}\tau _{k1,a}𝒪^c𝒪^d_0\left(\tau _{\mathrm{}1,b}𝒪_c_1𝒪_d_1+\frac{1}{24}\tau _{\mathrm{}1,b}𝒪_c𝒪_d_1\right)\hfill \\ & \frac{4}{5}\tau _{k1,a}\tau _{\mathrm{}1,b}𝒪_c_0(𝒪^c𝒪_d_1𝒪^d_1+\frac{1}{24}𝒪^c𝒪_d𝒪^d_1)\hfill \\ & +\frac{1}{48}\tau _{k1,a}𝒪_c𝒪_d𝒪^d_0𝒪^c\tau _{\mathrm{}1,b}_1+\frac{1}{48}\tau _{k1,a}𝒪_c_1\tau _{\mathrm{}1,b}𝒪^c𝒪_d𝒪^d_0\hfill \\ & +\frac{23}{240}\tau _{k1,a}\tau _{\mathrm{}1,b}𝒪_c𝒪_d𝒪^d_0𝒪^c_1\frac{1}{80}\tau _{k1,a}\tau _{\mathrm{}1,b}𝒪_c_1𝒪^c𝒪^d𝒪_d_0\hfill \\ & +\frac{7}{30}\tau _{k1,a}\tau _{\mathrm{}1,b}𝒪_c𝒪_d_0𝒪^c𝒪^d_1+\frac{1}{576}\tau _{k1,a}\tau _{\mathrm{}1,b}𝒪_c𝒪^c𝒪_d𝒪^d_0.\hfill \end{array}\end{array}$$
Taking $`k`$ and $`\mathrm{}`$ equal to $`1`$ and using the topological recursion relations in genus $`0`$ and $`1`$, we obtain the system of differential equations
(2.4)
$$D_{1,1,a,b}_2=_{1,1,a,b},$$
where $`D_{1,1,a,b}=D_{1,a}D_{1,b}3𝒪_a𝒪_b𝒪^c_0D_{1,c}`$, and
$`_{1,1,a,b}`$ $`=\frac{13}{10}𝒪_a𝒪_b𝒪_c𝒪_d_0𝒪^c_1𝒪^d_1`$
$`+\frac{4}{5}\left(𝒪_a𝒪_c_1𝒪_d_1+\frac{1}{24}𝒪_a𝒪_c𝒪_d_1\right)𝒪_b𝒪^c𝒪^d_0`$
$`+\frac{4}{5}𝒪_a𝒪^c𝒪^d_0\left(𝒪_b𝒪_c_1𝒪_d_1+\frac{1}{24}𝒪_b𝒪_c𝒪_d_1\right)`$
$`\frac{4}{5}𝒪_a𝒪_b𝒪_c_0(𝒪^c𝒪_d_1𝒪^d_1+\frac{1}{24}𝒪^c𝒪_d𝒪^d_1)`$
$`+\frac{1}{48}𝒪_a𝒪_c𝒪_d𝒪^d_0𝒪^c𝒪_b_1+\frac{1}{48}𝒪_a𝒪_c_1𝒪_b𝒪^c𝒪_d𝒪^d_0`$
$`+\frac{23}{240}𝒪_a𝒪_b𝒪_c𝒪_d𝒪^d_0𝒪^c_1\frac{1}{80}𝒪_a𝒪_b𝒪_c_1𝒪^c𝒪^d𝒪_d_0`$
$`+\frac{7}{30}𝒪_a𝒪_b𝒪_c𝒪_d_0𝒪^c𝒪^d_1+\frac{1}{576}𝒪_a𝒪_b𝒪_c𝒪^c𝒪_d𝒪^d_0.`$
We now specialize to the case of the $`A_2`$ model. In this model, there are two primary fields $`𝒪_0`$ and $`𝒪_1`$, with intersection form $`\eta _{ab}=\delta _{a+b,1}`$. Denote the associated coordinates $`u=𝒪_0𝒪_0_0=^2_0`$ and $`v=𝒪_0𝒪_1_0`$. The matrix $`𝒰`$ is given by the formula
$$𝒰=\left[\begin{array}{cc}𝒰_0^0& 𝒰_0^1\\ 𝒰_1^0& 𝒰_1^1\end{array}\right]=\left[\begin{array}{cc}v& u\\ u^2& v\end{array}\right],$$
and $`_1=\frac{1}{24}\mathrm{log}(\mathrm{\Delta })`$. As was shown by Eguchi, Yamada and Yang , the genus $`2`$ potential of the $`A_2`$-model is given by the formula
(2.5) $`_2`$ $`=\frac{1}{1152}^2𝒪_a𝒪_b𝒪_c𝒪_d_0𝒞^{ab}𝒞^{cd}`$
$`\frac{1}{1152}^2𝒪_a𝒪_b_0𝒪_c𝒪_d𝒪_e𝒪_f_0𝒞^{ac}𝒞^{bd}𝒞^{ef}`$
$`\frac{1}{360}^2𝒪_a𝒪_b𝒪_c_0𝒪_d𝒪_e𝒪_f_0𝒞^{ad}𝒞^{be}𝒞^{cf}`$
$`+\frac{1}{360}^2𝒪_a𝒪_b_0𝒪_c𝒪_d𝒪_e_0𝒪_f𝒪_g𝒪_h_0𝒞^{ac}𝒞^{bf}𝒞^{dg}𝒞^{eh}.`$
It may be checked that this function solves the equations (2.2) and (2.4).
For an arbitrary theory of topological gravity, let $`_{2,0}`$ be the function on the large phase space given by formula (2.5). For all theories of topological gravity for which we know the genus $`2`$ potential, the function $`_{2,0}`$ appears to be a major contribution to this potential.
We now turn to the case of $`^1`$. As in the $`A_2`$-model, there are two primary fields $`𝒪_0`$ and $`𝒪_1`$, with intersection form $`\eta _{ab}=\delta _{a+b,1}`$. Again, denote the associated coordinates by $`u=𝒪_0𝒪_0_0=^2_0`$ and $`v=𝒪_0𝒪_1_0`$. The matrix $`𝒰`$ is now given by the formula
$$𝒰=\left[\begin{array}{cc}v& u\\ e^u& v\end{array}\right],$$
and $`_1=\frac{1}{24}\mathrm{log}(\mathrm{\Delta })\frac{1}{24}u`$.
The correlators $`\tau _{1,a_1}𝒪_{a_2}\mathrm{}𝒪_{a_n}_2`$ and $`𝒪_{a_1}𝒪_{a_2}\mathrm{}𝒪_{a_n}_2`$ vanish in the $`^1`$-model for dimensional reasons. It follows that the following solution to the equations (2.2) and (2.4) is the genus $`2`$ potential:
(2.6) $`_2=_{2,0}`$ $`\frac{1}{480}^3𝒪_a𝒪_b_0𝒞^{ab}+\frac{7}{5760}^3𝒪_a_0^2𝒪_b_0𝒞^{ab}`$
$`+\frac{11}{5760}^2𝒪_a𝒪_b_0^2𝒪_c𝒪_d_0𝒞^{ac}𝒞^{bd}.`$
The three additional terms reflect the fact that, unlike in the $`A_2`$-model, the function $`\psi (u)=\frac{1}{24}u`$ is nonzero in the $`^1`$-model.
The Toda conjecture of Eguchi and Yang (, , ) provides conjectural formulas for the functions $`𝒪_1𝒪_1_g`$, $`g>0`$, of the $`^1`$-model:
$$\underset{g=0}{\overset{\mathrm{}}{}}\lambda ^{2g}𝒪_1𝒪_1_g=\mathrm{exp}\left(\frac{2}{\lambda ^2}\left(\mathrm{cosh}(\lambda )1\right)\underset{g=0}{\overset{\mathrm{}}{}}\lambda ^{2g}_g\right).$$
In genus $`2`$, this yields the equation
(2.7)
$$𝒪_1𝒪_1_2=e^u\left(^2_2+\frac{1}{12}^4_1+\frac{1}{360}^6_0+\frac{1}{2}(^2_1+\frac{1}{12}^4_0)^2\right).$$
It is easily checked, using the explicit formula formula for $`_2`$, that this equation holds.
## 3. Models with two primaries
In this section, we consider topological gravity in a general background with two primary fields $`𝒪_0`$ and $`𝒪_1`$, and intersection form $`\eta _{ab}=\delta _{a+b,1}`$. It is not clear to what extent such a model, even if it possesses a consistent loop expansion, corresponds to a physical theory: it may be that only the $`A_2`$ and $`^1`$-models are physical theories. The fact that our equations remain consistent in this setting is nevertheless very suggestive.
Denote the associated coordinates $`u=𝒪_0𝒪_0_0`$ and $`v=𝒪_0𝒪_1_0`$. The genus $`0`$ sector is characterized by the function $`𝒪_1𝒪_1_0`$; by the string equation, this is a function of $`u`$ alone, and we denote it by $`\varphi (u)`$. The matrix $`𝒰`$ is given by the formula
$$𝒰=\left[\begin{array}{cc}v& u\\ \varphi (u)& v\end{array}\right].$$
In this section, the correlation functions $`\tau _{k_1,a_1}\mathrm{}\tau _{k_n,a_n}_g`$ are assumed to have the following form: they are holomorphic functions of $`\{(v,u)^2u(\mathrm{},0]\}`$, Laurent polynomials in $`\mathrm{\Delta }`$, and polynomial in the remaining coordinates $`\{^nv,^nun>0\}`$.
There is a universal differential equation in topological gravity relating the potentials $`_0`$ and $`_1`$. In the case of two primary fields, this equation says that
(3.1)
$$\frac{1}{24}\varphi ^{\prime \prime \prime }+\varphi ^{\prime \prime }\psi ^{}2\varphi ^{}\psi ^{\prime \prime }=0.$$
It turns out that this equation is also the necessary and sufficient condition for the system of equations (2.2) and (2.4) to have a solution. The necessity follows from the formula
$$D_{1,1,0,0}_{2,0}D_{2,0}_{1,1,0,0}=\frac{2}{15}(u)^3(4(v)^2+(u)^2\varphi ^{})(\frac{1}{24}\varphi ^{\prime \prime \prime }+\varphi ^{\prime \prime }\psi ^{}2\varphi ^{}\psi ^{\prime \prime }).$$
###### Theorem 3.1.
Suppose that $`\frac{1}{24}\varphi ^{\prime \prime \prime }+\varphi ^{\prime \prime }\psi ^{}2\varphi ^{}\psi ^{\prime \prime }=0`$. Then the equations (2.2) and (2.4) have the solution $`_{2,0}+_{2,1}`$, where $`_{2,0}`$ is given by (2.5), and
$`_{2,1}`$ $`=\frac{1}{576}((\frac{1}{2}_{0,a}_{0,b}\psi +\frac{4}{5}_{0,a}\psi _{0,b}\psi )𝒞^{ab}`$
$`+^2𝒪_a𝒪_b_0(\frac{6}{5}_{0,c}_{0,d}\psi \frac{1}{10}_{0,c}\psi _{0,d}\psi )𝒞^{ac}𝒞^{bd}`$
$`+(\frac{7}{10}^2𝒪_a𝒪_b𝒪_c_0_{0,d}\psi \frac{3}{10}𝒪_a𝒪_b𝒪_c_0_{0,d}\psi )𝒞^{ab}𝒞^{cd}`$
$`+^2𝒪_a𝒪_b_0𝒪_c𝒪_d𝒪_e_0_{0,f}\psi (\frac{3}{10}𝒞^{af}𝒞^{bc}𝒞^{de}\frac{23}{10}𝒞^{ac}𝒞^{bd}𝒞^{ef})`$
$`+\frac{1}{10}(u)^4\varphi ^{\prime \prime }\psi ^{\prime \prime }\mathrm{\Delta }^1).`$
This solution may be characterized by the property that its restriction to the small phase space, together with the restrictions of the functions $`_{1,a}(_{2,0}+_{2,1})`$, vanish.
All of the terms in the formula for $`_{2,0}+_{2,1}`$ except the last one $`\frac{1}{5760}(u)^4\varphi ^{\prime \prime }\psi ^{\prime \prime }\mathrm{\Delta }^1`$ are associated to Feynman graphs with propagator $`𝒞`$ and vertices $`^n_{a_1}\mathrm{}_{a_{k2}}𝒰_{a_{k1}a_k}`$ and $`^n_{a_1}\mathrm{}_{a_k}\psi `$. From this point of view, the last term is an instanton, which vanishes if $`\psi `$ is a linear function of $`u`$, that is, for the $`A_2`$ and $`^1`$-models.
One calculates that $`_{2,1}`$ is given by the explicit formula
$`_{2,1}`$ $`=\frac{1}{576}(\frac{1}{2}(u)^2\psi ^{\prime \prime \prime }+\frac{9}{5}(u)^2\psi ^{\prime \prime }\psi ^{}+\frac{13}{5}^2u\psi ^{\prime \prime }+\frac{7}{10}^2u(\psi ^{})^2`$
$`((v)^2+\frac{7}{5}(u)^2)(u)^2\varphi ^{}\psi ^{\prime \prime }\psi ^{}\mathrm{\Delta }^1+\frac{6}{5}(u)^4\varphi ^{\prime \prime }(\psi ^{})^2\mathrm{\Delta }^1`$
$`+(\frac{2}{5}(^2vv^2uv)v\frac{1}{10}(u)^4\varphi ^{\prime \prime })\psi ^{\prime \prime }\mathrm{\Delta }^1`$
$`+(\frac{12}{5}^3vv\frac{12}{5}^3uu\varphi ^{}\frac{7}{5}^2u(u)^2\varphi ^{\prime \prime })\psi ^{}\mathrm{\Delta }^1`$
$`+\frac{11}{5}(4^2v^2uvu\varphi ^{}(^2v+^2u\varphi ^{})((v)^2+(u)^2\varphi ^{})`$
$`+2(^2vu^2uv)(u)^2v\varphi ^{\prime \prime }\frac{1}{2}(u)^6(\varphi ^{\prime \prime })^2)\psi ^{}\mathrm{\Delta }^2).`$
Now let $`_2`$ be a general solution of (2.2) and (2.4). Write $`_2=_{2,0}+_{2,1}+f_2`$. By the equations (2.2), $`D_{k,a}f_2=0`$ for $`k>1`$; thus, $`f_2`$ is a function of the coordinates $`\{u,v,u\}`$.
###### Theorem 3.2.
Define the functions $`h_a=h_a(u)`$ by the formula $`h_a={\displaystyle \frac{f_2}{(u^a)}}|_{(v,u)=(1,0)}`$. Then
$$f_2=\frac{1}{2}u^a𝒰_a^bh_b=\frac{1}{2}u^au^b𝒜_{ab}^ch_c=\frac{1}{2}\left((v)^2+\varphi ^{}(u)^2\right)h_0(u)+vuh_1(u).$$
###### Proof.
Let $`\stackrel{~}{f}_2=\frac{1}{2}u^a𝒰_a^bh_b`$; then $`D_{1,1,a,b}\stackrel{~}{f}_2=0`$ and $`h_a={\displaystyle \frac{\stackrel{~}{f}_2}{(u^a)}}|_{(v,u)=(1,0)}`$.
Thus $`f_2\stackrel{~}{f}_2`$ satisfies the equations $`D_{k,a}(f_2\stackrel{~}{f}_2)=0`$ for $`k>1`$, and $`D_{1,1,a,b}(f_2\stackrel{~}{f}_2)=0`$, as well as the dilaton equation $`𝒟(f_2\stackrel{~}{f}_2)=2(f_2\stackrel{~}{f}_2)`$, and is thus determined by the restrictions of the partial derivatives $`_{1,a}(f_2\stackrel{~}{f}_2)`$ to the small phase space. But these vanish; we conclude that $`f_2=\stackrel{~}{f}_2`$. ∎
In the next section, we determine the functions $`h_a`$.
## 4. Virasoro constraints
We now show that the Virasoro constraints $`L_0Z=L_1Z=0`$ of Eguchi, Hori and Xiong , as generalized to arbitrary Frobenius manifolds by Dubrovin and Zhang , may be used to complete the determination of the genus $`2`$ potential in two-primary models of topological gravity.
### The constraint $`L_0Z=0`$
According to Dubrovin and Zhang , an Euler vector on a Frobenius manifold determines matrices $`\mu `$ and $`R[n]`$, $`n>0`$, which satisfy the commutation relations $`[\mu ,R[n]]=nR[n]`$ and the symmetry conditions $`\mu _{ab}+\mu _{ba}=0`$ and
$$R[n]_{ab}+(1)^nR[n]_{ba}=0.$$
The basis $`𝒪_a`$ of primary fields may be chosen in such a way that the matrix $`\mu `$ is diagonal
$$\mu _a^b=\delta _a^b\mu _a,$$
and $`\mu _0<\mu _a`$ for $`a0`$. Setting $`d_a=\mu _a\mu _0`$ and $`d=2\mu _0`$, we have $`\mu _a=d_ad/2`$.
For the Gromov-Witten invariants of a Kähler manifold $`X`$, the primaries $`𝒪_a`$ form a basis of the De Rham cohomology $`H^{}(X,)`$, and the number $`d_a`$ is the holomorphic degree of $`𝒪_a`$, that is $`𝒪_aH^{d_a,}(X,)`$. (In particular, $`d`$ equals the complex dimension of $`X`$.) In this case, $`R[1]`$ is the matrix of multiplication by $`c_1(X)`$, and $`R[n]=0`$ for $`n>1`$.
Introduce the vector field
$$_0=\underset{k=0}{\overset{\mathrm{}}{}}\left\{(\mu _a^b+k+\frac{1}{2})\stackrel{~}{t}_k^a_{k,b}+\underset{\mathrm{}=1}{\overset{k}{}}R[\mathrm{}]_a^b\stackrel{~}{t}_k^a_{k\mathrm{},b}\right\},$$
where $`\stackrel{~}{t}_k^a`$ are the shifted coordinates $`\stackrel{~}{t}_k^a=t_k^a\delta _{k,1}\delta _0^a`$. The Virasoro constraint $`L_0Z=0`$ in genus $`g=0`$ may be expressed as the following equation:
(4.1)
$$_0𝒰+𝒰+[\mu ,𝒰]+R[1]=0.$$
In genus $`g>0`$, the Virasoro constraint $`L_0Z=0`$ says that
(4.2)
$$_0_g+\frac{1}{4}\delta _{g,1}\mathrm{Tr}(\frac{1}{4}\mu ^2)=0.$$
These equations are known to hold for Gromov-Witten invariants .
Let $`=^a/u^a`$ be the Euler vector field, where
(4.3)
$$^a=(1d_a)u^a+R[1]_0^a.$$
Then (4.1) implies that
(4.4)
$$_0u^a+^a=0.$$
In calculating the action of the vector field $`_0`$ in the coordinate system $`\{u_n^a\}`$, we use (4.4) together with the commutation relation $`[,_0]=\frac{1}{2}(1d)`$.
In the case of two primary fields, we have $`\mu =\frac{1}{2}\left[\begin{array}{cc}d& 0\\ 0& d\end{array}\right]`$. Consider first the case in which $`d`$ equals $`1`$; then $`R[1]=\left[\begin{array}{cc}0& r\\ 0& 0\end{array}\right]`$. By (4.1), we see that $`\varphi =ce^{2u/r}`$; redefining $`u`$, we may assume that $`r=2`$, and we recover the $`^1`$-model. Since $`\mathrm{Tr}(\mu ^2\frac{1}{4})=0`$, we see from (4.2) that $`_0_1=0`$; (3.1), now shows that $`\psi =\frac{1}{24}u`$, consistent with the known form of $`_1`$ in the $`^1`$-model.
The equation $`_0_2=0`$ of (4.2) constrains the functions $`h_a(u)`$ of Theorem 3.2; if $`d=1`$, it forces them to have negative degree in $`e^u`$, and hence to vanish, as we have already observed,
If $`d1`$, the matrix $`R[1]`$ vanishes. By (4.1), we see that $`\varphi (u)=u^{(1+d)/(1d)}`$, up to a constant which we take to equal $`1`$. (For example, the $`A_2`$-model, has $`d=\frac{1}{3}`$ and $`\varphi (u)=u^2`$.) In genus $`1`$, the equation (4.2) shows that $`\psi (u)`$ is proportional to $`\mathrm{log}(u)`$; both (3.1) and (4.2) yield the same answer for this constant,
$$\psi (u)=\frac{d(3d1)}{24(d1)}\mathrm{log}(u).$$
Note that the $`A_2`$-model, for which $`d=\frac{1}{3}`$, has $`\psi =0`$. The equation $`_0_2=0`$ imposes the homogeneities $`h_a(u)=C_au^{((1+a)d3)/(1d)}`$.
### The constraint $`L_1Z=0`$
Let $`_1`$ be the vector field
$`_1`$ $`=(\mu _a\frac{1}{2})(\mu _a+\frac{1}{2})𝒪^a_0_{0,a}+{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\{(\mu _a+k+\frac{1}{2})(\mu _a+k+\frac{3}{2})\stackrel{~}{t}_k^a_{k+1,a}`$
$`+{\displaystyle \underset{\mathrm{}k+1}{}}2(\mu _a+k+1)R[\mathrm{}]_a^b\stackrel{~}{t}_k^a_{k+1\mathrm{},b}+{\displaystyle \underset{\mathrm{}_1+\mathrm{}_2k+1}{}}(R[\mathrm{}_1]R[\mathrm{}_2])_a^b\stackrel{~}{t}_k^a_{k+1\mathrm{}_1\mathrm{}_2,b}\}.`$
Let $`𝒱=𝒰`$; by (4.1), $`𝒱=𝒰+[\mu ,𝒰]+R[1]`$. The constraint $`L_1Z=0`$ in genus $`0`$ may written (Dubrovin and Zhang ; cf. Theorem 5.7 of )
(4.5)
$$_1𝒰+𝒱^2=0.$$
In particular, we see that
(4.6)
$$_1u^a+^b^c𝒜_{bc}^a=0.$$
In genus $`g>0`$,the constraint $`L_1Z=0`$ is
(4.7)
$$_1_g+\frac{1}{2}(\frac{1}{4}\mu ^2)^{ab}\left(\underset{h=1}{\overset{g1}{}}𝒪_a_h𝒪_b_{gh}+𝒪_a𝒪_b_{g1}\right)=0.$$
In the case of two primaries, this becomes
(4.8)
$$_1_g+\frac{1}{8}(1d^2)\left(\underset{h=1}{\overset{g1}{}}𝒪_0_h𝒪_1_{gh}+𝒪_0𝒪_1_{g1}\right)=0.$$
In calculating the action of the vector field $`_1`$ in the coordinate system $`\{u_n^a\}`$, we use (4.6) and the commutation relation
(4.9)
$$[_{0,a},_1]=\left((\mu +\frac{1}{2})(\mu +\frac{3}{2})\right){}_{a}{}^{b}D_{1,b}^{}+\left((\mu +\frac{1}{2})𝒱+𝒱(\mu +\frac{1}{2})\right){}_{a}{}^{b}D_{0,b}^{}.$$
In the case of two primaries, this implies that
$$[,_1]=\{\begin{array}{cc}(1d)\left(\frac{1}{4}(3d)D_{1,0}+vD_{0,0}+uD_{0,1}\right),\hfill & d1,\hfill \\ 2D_{0,1},\hfill & d=1.\hfill \end{array}$$
Using these formulas, we see that the case $`g=2`$ of (4.8) yields the equation
$$\begin{array}{c}0=_1_2+\frac{1}{4}(1d^2)\left(𝒪_0_1𝒪_1_1+𝒪_0𝒪_1_1\right)\hfill \\ \hfill \begin{array}{cc}\hfill =& 6\left((d+1)C_0+\frac{1}{5760}d(3d1)(3d5)(d2)\right)u^2vu\hfill \\ & +3C_1(d1)u^{(d2)/(1d)}((v)^2+\varphi ^{}(u)^2).\hfill \end{array}\end{array}$$
It follows that $`h_1=0`$ and
(4.10)
$$h_0=\frac{d(3d1)(3d5)(d2)}{5760(d+1)}u^{(d3)/(1d)}.$$
completing the determination of $`_2`$.
Our formula for $`_2`$ agrees with that of Dubrovin and Zhang , who apply the method of Eguchi and Xiong ; in other words, they use the constraints $`D_{k,a}_2=0`$, $`k>4`$, and $`L_nZ=0`$, $`n10`$.
### The higher Virasoro constraints
The higher Virasoro constraints are given by formulas involving a Lie algebra of vector fields $`_n`$, $`n1`$, on the large phase space, which satisfy the commutation relations
$$[_m,_n]=(mn)_{m+n}.$$
This Lie algebra is generated by $`_1`$ and $`_n`$, for any $`n>1`$.
Just as for $`_0`$ and $`_1`$, we can avoid using the explicit formula for $`_n`$. The Virasoro constraint $`L_nZ=0`$ in genus $`0`$ may be written
(4.11)
$$_n𝒰+𝒱^{n+1}=0.$$
In calculating the action of the vector field $`_2`$ in the coordinate system $`\{u_n^a\}`$, we use (4.11) and the commutation relation
(4.12)
$$[_{0,a},_n]=\underset{i=0}{\overset{n}{}}(𝖡_{n,i})_a^bD_{i,b},$$
where the matrices $`𝖡_{n,i}`$ are determined by the recursion
$$𝖡_{n,i}=(\mu +i+\frac{1}{2})𝖡_{n1,i1}+𝒱𝖡_{n1,i},$$
with initial condition $`𝖡_{1,i}=\delta _{i+1,0}`$.
In the case of two primaries, with $`n=2`$, this implies that
$$[,_2]=\{\begin{array}{cc}(1d)(\frac{1}{8}(3d)(5d)D_{2,0}+\frac{3}{4}(3d)vD_{1,0}+\frac{1}{4}(d^22d+9)uD_{1,1}\hfill & \\ +(\frac{3}{2}v^2+\frac{1}{2}(1+d)(3d)u\varphi (u))D_{0,0}+3uvD_{0,1}),\hfill & d1,\hfill \\ 6D_{1,1}+4e^uD_{0,0}+6vD_{0,1},\hfill & d=1.\hfill \end{array}$$
In genus $`g>0`$,the constraint $`L_2Z=0`$ is
$`_2_g`$ $`+(\mu _a\frac{3}{2})(\mu _a\frac{1}{2})(\mu _a+\frac{1}{2})\eta ^{ab}\left({\displaystyle \underset{h=1}{\overset{g1}{}}}\tau _{1,a}_h𝒪_b_{gh}+\tau _{1,a}𝒪_b_{g1}\right)`$
$`\frac{1}{2}(3\mu _a^2+3\mu _a\frac{1}{4})R[1]^{ab}\left({\displaystyle \underset{h=1}{\overset{g1}{}}}𝒪_a_h𝒪_b_{gh}+𝒪_a𝒪_b_{g1}\right)=0.`$
It may be verified that $`_2`$ satisfies this equation.
Since the differentials operators $`L_n`$, $`n>1`$, lie in the Lie algebra generated by $`L_1`$ and $`L_2`$, it follows that the Virasoro conjecture holds to genus $`2`$ for two-primary models.
### The Belorousski-Pandharipande equation
The Belorousski-Pandharipande equation is a differential equation satisfied by the genus $`2`$ potential, analogous to the equation (3.1) in genus $`1`$; it may be expressed as saying that a certain cubic polynomial in the coordinates $`\{u^a\}`$ vanishes. It turns out that in the case of backgrounds with two primaries, the equation gives a second (and thus rigorous) derivation of the above formula for $`C_0`$, but leaves $`C_1`$ undetermined.
Taking Theorem 3.2 and the equations $`_1_2=0`$ and $`_0_2=0`$ into account, the Belorousski-Pandharipande equation reduces to a single equation
$$\varphi ^{}h_0^{}\frac{1}{2}\varphi ^{\prime \prime }h_0\frac{1}{48}\psi ^{\prime \prime \prime \prime }\frac{3}{5}\psi ^{\prime \prime \prime }\psi ^{}+\frac{9}{10}(\psi ^{\prime \prime })^2=0.$$
With $`\varphi (u)=u^{(1+d)/(1d)}`$ and $`\psi (u)=\frac{d(3d1)}{24(d1)}\mathrm{log}(u)`$, the function $`h_0`$ of (4.10) satisfies this equation.
## Acknowledgments
The authors thank B. Dubrovin and Y. Zhang for communication of their results prior to publication, and for a number of interesting conversations. The second author thanks Y. Zhang and Qinghua University for its hospitality while working on this paper.
The research of the first author is supported in part by special priority area #707 of the Japanese Ministry of Education. The research of the second author is supported in part by NSF grant DMS-9704320, and by the special year “Geometry of String Theory” of RIMS. The research of the third author is supported in part by NSFC grant 19925521 and by a startup grant from Beijing University.
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# Lévy Flights from a Continuous-Time Process
## I Introduction
Random walk processes leading to subdiffusive or superdiffusive behavior are adequate for describing various physical situations. Thus, the continuous-time random walk (CTRW) model of Scher and Montroll was a milestone in understanding of photoconductivity in strongly disordered and glassy semiconductors, while the Lévy-flight models are adequate for description of transport in heterogeneous catalysis , self-diffusion in micelle systems , reactions and transport in polymer systems under conformational motion , transport processes in heterogeneous rocks , and for description of behavior of dynamical systems . The closely related models appear in description of economic time series . The Lévy-related statistics were observed in hydrodynamic transport , and in the motion of gold nanoclusters on graphite . The mixed models were proposed , in which the slow temporal evolution (described by Scher-Montroll CTRW) is combined with the possibility of Lévy-jumps, so that in general both sub- or superdiffusive behavior can arise .
The continuous-time random walks (CTRW) first introduced by Montroll and Weiss correspond to a stochastic model in which steps of a simple random walk take place at times $`t_i`$, following some random process with non-negative increments: $`\tau _i=t_it_{i1}0`$. In a mathematical language one says that CTRW is a process subordinated to random walks under the operational time defined by the process $`\left\{t_i\right\}`$. It is typically thought that a CTRW-scheme alone can not describe any superdiffusive process, so that the introduction of very long jumps is an inevitable part of building a model leading to superdiffusive behavior.
Let us first discuss a typical CTRW approach. Let us consider a one-dimensional situation under which a particle from time to time makes a jump to a neighboring lattice site separated from the initial one by a distance $`a`$. The time $`\tau `$ between the two jumps is distributed according to some waiting-time distribution, represented by the probability density function (PDF) $`p(\tau )`$. If the mean waiting time $`\overline{\tau }`$ exists, the particle’s behavior is diffusive, with diffusion coefficient $`D=a^2/2\overline{\tau }`$. If the corresponding moment diverges, the particle’s behavior becomes subdiffusive, with $`r^2(t)t^\alpha `$, with $`\alpha <1`$ depending on the PDF $`p(\tau )`$. The subdiffusive behavior is indicated by vanishing of the diffusion coefficient $`D`$. It seems impossible to obtain within this scheme any type of a superdiffusive behavior unless one allows for infinitely long jumps with $`a^2\mathrm{}`$. The superdiffusive behavior is indicated by divergence of the diffusion coefficient $`D`$. If $`a^2`$ stays finite this can be the case only if $`\overline{\tau }`$ vanishes. Since $`\tau >0`$ and $`\overline{\tau }=_0^{\mathrm{}}\tau p(\tau )𝑑\tau ,`$ vanishing of the mean waiting time means that $`p(\tau )=\delta (\tau )`$, a marginal, degenerate situation.
On the other hand the consideration presented above shows only that the waiting-time distribution is not an adequate tool for description of superdiffusive CTRW. In what follows we show that superdiffusive CTRW with bounded step lengths are just as possible as the subdiffusive ones. Our considerations will be rather formal and do not follow from any particular physical model. On the other hand, the fact that Lévy-flights can stem from a process subordinated to simple random walks has many important implications. Thus, as we proceed to show, the fast dynamics of a free process can coexist in such models with simple exponential relaxation to a normal Boltzmann equilibrium distribution, if the behavior of an ensemble of random walkers under restoring force is considered. This shows that the relation between Lévy dynamics and the nonextensive thermodynamics described by nonclassical entropy functions is much looser than typically assumed.
The combinations of the superdiffusive Lévy-flights with the typical CTRW operational time leads to paradoxical diffusion behavior, having some parallels in transport on polymer chains. Moreover, the existence of a subordination model leading to Lévy flights can be useful in understanding of statistical implications of the processes described by fractional generalizations of diffusion and Fokker-Planck equations .
The article is organized as follows: In Sec. 2 we discuss general properties of subordinated random processes. In Sec.3 and 4 the processes subordinated to symmetric and asymmetric random walks are considered, these leading to symmetric and asymmetric Lévy-flights. The dualism between the Lévy-flights and the Scher-Montroll CTRW is discussed in Sec.5. Sections 6 and 7 discuse the models leading to paradoxical diffusion behavior. The relaxation to equilibrium is considered in Sec. 8.
## II The subordination of random processes
As already mentioned, a Scher-Montroll CTRW process is a simple random walk whose steps take place at times $`t_i`$ governed by a random process with nonnegative independent increments, so that
$$P(x,t)=\underset{n}{}P_{RW}(x,n)p_n(t),$$
(1)
where $`P_{RW}(x,n)`$ is a probability distribution to find a random walker at point $`x`$ after $`n`$ steps (i.e. the binomial distribution), and $`p_n(t)`$ is the probability to make exactly $`n`$ steps up to time $`t`$. For both $`t`$ and $`n`$ large, when the binomial distribution can be approximated by a Gaussian one, and when the corresponding sum can be changed to an integral, Eq.(1) reads:
$$P(x,t)_0^{\mathrm{}}\frac{1}{\sqrt{2\pi n}}\mathrm{exp}\left(\frac{x^2}{2n}\right)p_t(n,t)𝑑n.$$
(2)
In a classical Scher-Montroll CTRW $`p_t(n,t)`$ corresponds to a random process in which $`n`$ typically grows sublinearly in $`t`$. Thus, the overall process is subdiffusive.
Note that a description of CTRW-process given by Eq.(2) is an example of subordination, see Sec. X.7 of Ref.: If $`\left\{X(T)\right\}`$ is a Markov process with continuous transition probabilities and $`\left\{T(t)\right\}`$ a process with non-negative independent increments, then $`\left\{X(T(t))\right\}`$ is said to subordinate to $`\left\{X(t)\right\}`$ using the operational time $`T`$. In this case
$$P(x,t)=_0^{\mathrm{}}P_x(x,T)p_T(T,t)𝑑T.$$
(3)
In what follows we call the integral transform, Eq.(3) a subordination transformation, changing from time scale $`t`$ to a time-scale $`T`$. For example, in the Scher-Montroll case the operational time of a system is given by the number of steps of the RW, and is a random function of the physical time $`t`$ whose typical value grows sublinearly in $`t`$.
The operational time can also grow superlinearly with $`t`$. Such a process can not be described by a waiting-time distribution, and needs a complimentary description. Let us consider a random process, where the density of events fluctuates strongly. Let us subdivide the time axis into intervals of duration $`\mathrm{\Delta }t`$ and let us consider the number $`n`$ of jumping events within each interval. The value $`\rho =n/\mathrm{\Delta }t`$ defines the density of jump events. Now, if the mean density of events exists, its inverse gives us exactly the mean waiting time of a jump, and a process described by a finite density of events is a normal diffusive one. The divergence of a mean waiting time (like in Scher-Montroll CTRW) correspond to vanishing density. On the other hand, if one considers a strongly fluctuating density $`\rho (t)`$ whose first moment diverges, the mean waiting time vanishes and a process that subordinates a random walk process under such operational time can be superdiffusive. At longer times, the distribution of the number of events tends to one of the Lévy-stable laws: the typical number of events can grow superlinearly in time. A simple example of such process was already known to Feller, see Chap. X.7 of Ref.. He considers a process subordinated to simple random walks under the operational time governed by a fully asymmetric Lévy stable law of index 1/2. The corresponding PDF at time $`t`$ is given by
$`P(x,t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{\sqrt{2\pi n}}}\mathrm{exp}\left({\displaystyle \frac{x^2}{2n}}\right){\displaystyle \frac{t}{\sqrt{\pi }n^{3/2}}}\mathrm{exp}\left({\displaystyle \frac{t^2}{2n}}\right)𝑑n`$ (4)
$`=`$ $`{\displaystyle \frac{t}{\pi (t^2+x^2)}},`$ (5)
i.e. is a Cauchy Lévy-flight.
Let us now discuss a simple analogy describing the relation between the Scher-Montroll CTRW and Lévy-flights. This analogy makes clear many of the findings we are going to discuss below. Imagine a physical clock producing ticks following with frequency 1, which govern the behavior of a random walker. Imagine a switch situated at $`0`$, so that returning to the origin, the walker can trigger some physical process (the analogy with the Glarum model of relaxation, Ref., is evident!). The times between the subsequent returns are distributed according to a fully asymmetric Lévy stable law of index 1/2 used in a previous example. Imagine now another random walker performing its motion (a step per physical unit time) independently from the first one. Imagine a movie camera, taking frame-per-frame pictures of the positions of this second random walker at the moments when the first walker is at the origin and thus triggers the switch. Watching the movie taken by the camera, we immediately recognize that the second walker performs the Cauchy Lévy-flights. Imagine, that a clock is posed in a frame and also filmed. In this case its image will show exactly the operational time of the system; the spectator’s watch measures the physical time. Imagine an opposite situation: the first walker triggers the motion of the second one, and the camera is triggered by the physical clock, as a normal movie camera is. The process we recognize at the film is then the Scher-Montroll CTRW. We can take a Scher-Montroll movie also using another trick (which cannot be performed in a real time, but needs a record of return times). Let us take a record of subsequent return times of a first random walker (numbers $`n_1,`$ $`n_2,`$ …) and trigger our camera in such a way that it makes $`n_1`$ frames during the first second, $`n_2`$ frames during the second one, etc. If we film a normal random walker with a camera prepared in such a way, the movie will show us the Montroll-Weiss CTRW. An image of the physical clock will again show the operational time of the system, and again, looking at his watch, the spectator can measure the physical time between two events.
Let us use our camera triggered by returns of a random walker to film other processes taking place in the outer world. The film which is watched afterwards under constant speed shows us a possible world: The causality relations and thermodynamical time arrow are those of our usual world. On the other hand, a movie of a world undergoing continuous evolution, in which ”natura non facit saltus” holds, will show us a revolutionary world of ”great leaps” and abrupt changes (but following the same logics of development). The second camera (fed by a prescribed $`n`$-sequence) will show us the world of almost full stagnation seldomly interrupted by a bounded, local movement, a world developing in a slow time of old Asiatic despoty. We shall keep this analogy in mind when discussing the physical implications of subordination.
Let us consider a system which evolves according to a Markovian dynamics and whose state tends to a normal Boltzmann equilibrium under relaxation. In a system under action of outer forces, the transition probabilities between the states of the system (sites $`i`$ between which the random walk takes place) which are characterized by their energies $`E_i`$, are not independent. They are connected through the corresponding Boltzmann-factors, so that in equilibrium during any period of time $`\mathrm{\Delta }t`$ the mean numbers of forwards and of backwards jumps between any two sites $`i`$ and $`j`$ fulfill the condition
$$n_{ij}(\mathrm{\Delta }t)/n_{ji}(\mathrm{\Delta }t)=\mathrm{exp}\left[\left(E_iE_j\right)/kT\right],$$
(6)
where $`k`$ is the Boltzmann constant and $`T`$ is the system’s temperature. The condition Eq.(6) guarantees detailed balance in equilibrium, independentl of what the real dynamics of a system is. For simple RWs, where only transitions between the neighboring states are allowed, the corresponding transition rates with respect to the operational time of the system can be introduced. For a random a walker moving under the influence of a weak constant force $`F`$ the probabilities of the forward and backward jumps per unit time $`w_+`$ and $`w_{}`$ are connected through $`w_+/w_{}=\mathrm{exp}(Fa/kT)`$. The Markovian nature of RW then leads to the fact that the values of $`w_+`$ and $`w_{}`$ do not depend on whether the system is in equilibrium or not. For $`F`$ small one can take, say, $`w_+=w_0(1+Fa/kT)`$ and $`w_{}=w_0(1Fa/kT)`$ with $`w_0=1/2\tau `$.
Note that subordination, describing a transition from a physical time to an operational time of the system, does not change its equilibrium properties. Such subordination can be considered as random modulation of the transition rate $`w_0`$ by some independent process (say closing and opening the channels), and is fully irrelevant for thermodynamics (i.e. thermostatics) of the system. On the other hand, it strongly influences its kinetics, so that a question can be posed, what kinds of kinetics are compatible with the relaxation to a normal Boltzmann distribution under arbitrary subordination transformation of time. We address this question in Sec.8, after the free diffusion properties of superdiffusive CTRW will be discussed.
## III Symmetric Lévy flights from CTRW
Let us first concentrate on the symmetric random walk case. Let us consider a random process in which the number of events per given time is unbounded and follows, for example, a power-law distribution, $`p_n(t)tn^{1\alpha }`$ with $`0<\alpha 1`$ (this corresponds to the typical number of events scaling as $`nt^{1/\alpha }`$). Let us find the asymptotic behavior of $`P(x,t)`$ for $`t`$ large. Since the jumps during different intervals are uncorrelated, the PDF of $`n`$ for longer times converges to a fully asymmetric Lévy-stable law
$$p(n,t)t^{1/\alpha }L(n/t^{1/\alpha };\alpha ,\gamma )$$
(7)
with the asymmetry parameter $`\gamma =\alpha `$ (here the values of $`\gamma =\pm \alpha `$ correspond to the strongly asymmetric PDF that vanish identically for large positive (negative) $`x`$ values , while $`\gamma =0`$ corresponds to symmetric distributions; the notation in one of Ref.). Note that the Fourier-transforms of Lévy-stable laws are known: up to the translation $`P(k,t)`$ is equal to
$$f(\kappa )=\mathrm{exp}\left[\left|\kappa \right|^ae^{i\pi \gamma /2}\right]$$
(8)
(for $`0<\alpha <2`$, $`\alpha 1`$). The PDF is a real function, thus $`f(\kappa )=f^{}(k)`$. The corresponding function is analytical everywhere except for $`\kappa =0`$, so that the PDF is given by
$$L(x;\alpha ,\gamma )=\frac{1}{\pi }\text{Re}_0^{\mathrm{}}e^{ix\zeta \zeta ^\alpha e^{i\pi \gamma /2}}𝑑\zeta .$$
(9)
From Eq.9 the series expansions for $`L(y;\alpha ,\gamma )`$ follow, see Sec. XVII.6 of Ref.. In the case $`\alpha <1`$ one can move the path of integration to the negative imaginary axis (since the integrand tends to zero when $`\text{Im}\zeta \mathrm{}`$ due to the dominance of the linear term), which allows then for elementary integration after Taylor-expansion of $`\mathrm{exp}(A\zeta ^\alpha )`$. For $`1<\alpha <2`$ this dominance is no more the case, but the integrand still vanishes for $`\text{Im}\zeta \mathrm{}`$ in the case of symmetric distributions, while $`(i\left|\zeta \right|)^\alpha =\left|\zeta \right|^a(\mathrm{cos}\frac{\pi }{2}\alpha i\mathrm{sin}\frac{\pi }{2}\alpha )\mathrm{}`$ for $`1<\alpha <2`$. Thus the series which represents Lévy distributions for $`0<\alpha <1`$ and symmetric Lévy distribution also for $`1<\alpha <2`$ reads:
$$L(y;\alpha ,\gamma )=\frac{1}{\pi y}\underset{k=1}{\overset{\mathrm{}}{}}(1)^k\frac{\mathrm{\Gamma }(k\alpha +1)}{k!}\mathrm{sin}\left(\frac{k\pi }{2}(\gamma \alpha )\right)y^{\alpha k}$$
(10)
In general the Lévy-stable laws for $`1<\alpha <2`$ are given by another expansion,
$$L(y;\alpha ,\gamma )=\frac{1}{\pi y}\underset{k=1}{\overset{\mathrm{}}{}}(1)^k\frac{\mathrm{\Gamma }(1+k/\alpha )}{k!}\mathrm{sin}\left(\frac{k\pi }{2}(\gamma \alpha )\right)y^k$$
(11)
which also holds for asymmetric laws.
One can easily obtain the form of $`x`$-distributions by immediate integration: using Eq.(2) and a scaling form of a Lévy-distribution
$$p(x,t)_0^{\mathrm{}}\frac{1}{\sqrt{2\pi n}}\mathrm{exp}\left(\frac{x^2}{2n}\right)L(\frac{n}{t^{1/\alpha }};\alpha ,\alpha )\frac{dn}{t^{1/\alpha }}.$$
(12)
Using Eq.(10) and performing a term-by-term integration, we arrive to the series of integrals of the form
$`I_\mu (\zeta )`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{\sqrt{2\pi \xi }}}e^{\frac{\zeta ^2}{2\xi }}\xi ^\mu 𝑑\xi `$ (13)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }}}\left({\displaystyle \frac{\zeta ^2}{2}}\right)^{1/2\mu }\mathrm{\Gamma }(\mu 1/2).`$ (14)
For integral of the $`k`$-th term in Eq.(10) we have $`\mu =1+\alpha k`$. Let us concentrate first on the case $`0<\alpha <1`$. Using well-known relations for $`\mathrm{\Gamma }`$-function: $`\mathrm{\Gamma }(z+1)=z\mathrm{\Gamma }(z)`$ (Eq. (6.1.15) of Ref.) and $`\mathrm{\Gamma }(2z)=(2\pi )^{1/2}2^{2z1/2}\mathrm{\Gamma }(z)\mathrm{\Gamma }(z+1/2)`$ (Eq. (6.1.18) of ) we get
$$p(\zeta )=\frac{1}{\pi }\underset{k=1}{\overset{\mathrm{}}{}}(1)^k\frac{\mathrm{\Gamma }(2k\alpha +1)}{k!}\mathrm{sin}\left(k\pi \alpha \right)\left(\frac{\sqrt{2}}{\zeta }\right)^{2\alpha 1},$$
(15)
which represents a series expansion for a symmetric Lévy-stable law of index $`2\alpha `$, Eq.(10), for the scaled variable $`\zeta /\sqrt{2}`$. This corresponds to a form $`p(x,t)=t^{1/2\alpha }L(x/\sqrt{2}t^{2\alpha };2\alpha ,0)`$ of the $`x`$-distribution.
We note that taking Fourier-transform of the both parts of for symmetric RWs,
$$L(ax,2\alpha ,0)=_0^{\mathrm{}}\frac{1}{\sqrt{2\pi n}}\mathrm{exp}\left(\frac{x^2}{2n}\right)L(n;\alpha ,\alpha )𝑑n,$$
(16)
where $`a`$ is an unimportant scaling factor, we get:
$$\mathrm{exp}(A\left|k\right|^{2\alpha }t)=_0^{\mathrm{}}\mathrm{exp}\left(k^2n\right)L(n;\alpha ,\alpha )𝑑n,$$
(17)
which holds for any real $`k`$ (i.e. for any positive $`k^2`$), where $`A`$is a number factor. This gives us a general expression for a Laplace-transform of an asymmetric Lévy-distribution with $`\alpha <1`$: A Laplace-transform of $`L(n,\alpha ,\alpha )`$ is $`\mathrm{exp}(A\left|k\right|^\alpha t)`$. From this fact an important result follows:
$$L(ax;\alpha \beta ,0)=_0^{\mathrm{}}n^{1/\beta }L(x/n^{1/\beta };\beta ,0)L(n;\alpha ,\alpha )𝑑n:$$
(18)
A Lévy distribution with index $`\alpha \beta `$ is subordinated to a Lévy-distribution with index $`\beta <\alpha `$ under the operational time given by an asymmetric Lévy law of index $`\alpha <1`$. To see this, consider the characteristic functions of both sides of Eq.(18) and use Eq.(17).
$$\mathrm{exp}(A\left|k\right|^{\alpha \beta })=_0^{\mathrm{}}e^{\left|\kappa \right|^\beta n}L(n;\alpha ,\alpha )𝑑n,$$
(19)
see Sec.X.7 of Ref.. Eq.(17) corresponds to a special case of $`\beta =2`$ of Eq.(19). The distributions $`L(n;\alpha ,\alpha )`$ thus coincide with inverse Laplace transforms of stretched-exponentials. For example for $`L(n;1/2,1/2)`$ one readily gets:
$$p(n,t)=^1\left\{\mathrm{exp}(tu^{1/2})\right\}=\frac{t}{2\sqrt{\pi }n^{3/2}}\mathrm{exp}\left(\frac{t^2}{4n}\right)\text{,}$$
(20)
which differs only by a scale for the time-unit from a distribution used in the example Eq.(5).
## IV Asymmetric Lévy-flights
Imagine a random walker moving under the influence of a weak constant force $`F`$. Such force introduces an asymmetry into the walker’s motion, since the probabilities of the forward and backward jumps, $`w_+`$ and $`w_{}`$ are now weighed with the corresponding Boltzmann-factors, $`w_+/w_{}=\mathrm{exp}(Fa/kT)`$. For $`F`$ small one can take $`w_+=1/2+Fa/2\tau kT`$ and $`w_+=1/2+Fa/2\tau kT`$. For $`t`$ large such random walks lead to the Gaussian distribution of the particles’ positions
$$P_{RW}(x,t)=\frac{1}{\sqrt{2\pi n}}\mathrm{exp}\left(\frac{(xvn)^2}{2n}\right)$$
(21)
whose center moves with a constant velocity $`v=\mu F=Fa^2/2\tau kT`$. Note that our RW fulfil the Einstein’s relation between the mobility $`\mu `$ and diffusion coefficient $`D`$: $`\mu =D/kT`$. The PDF of a random process which subordinates biased RW under an operational time following the asymmetric Lévy-law is given by:
$`P(x,t){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{\sqrt{2\pi n}}}\mathrm{exp}\left({\displaystyle \frac{(xvn)^2}{2n}}\right)`$ (22)
$`\times L({\displaystyle \frac{n}{t^{1/\alpha }}};\alpha ,\alpha ){\displaystyle \frac{dn}{t^{1/\alpha }}}.`$ (23)
Using the series expansion, Eq.(10) and performing the term-by-term integration leads to the series of the integrals of the type:
$`I_\mu (\zeta ,\omega )`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{\sqrt{2\pi \xi }}}e^{\frac{(\zeta \omega \xi )^2}{2\xi }}\xi ^\mu 𝑑\xi `$ (24)
$`=`$ $`{\displaystyle \frac{2\mathrm{exp}(\zeta \omega )}{\sqrt{2\pi }}}\left({\displaystyle \frac{\zeta ^2}{\omega ^2}}\right)^{1/4\mu /2}K_{1/2\mu }\left(\zeta \omega \right)`$ (25)
for $`\omega 0.`$ For integral of the $`k`$-th term in Eq.(10) we again have $`\mu =1+\alpha k`$. Let us concentrate first on the case $`0<\alpha <1`$. For $`\zeta \omega `$ small, $`v`$ cancels (see the expansion 9.6.9 of Ref., $`K_\nu (z)\frac{1}{2}\mathrm{\Gamma }(\nu )(\frac{1}{2}z)^\nu `$ ($`\nu >0`$), note that $`K_\nu (z)=K_\nu (z)`$), so that the corresponding distribution tends to be a function of $`\zeta `$ only, it coincides with one for $`\omega =0`$, Eq.(14) so that a symmetric Lévy-stable law of index $`2\alpha `$, Eq.(10) emerges. On the other hand, for $`v0`$ and $`x`$ large the overall distributions follow from the expansion of $`K`$ for large values of the argument which reads: $`K_v(z)\sqrt{\frac{\pi }{2z}}e^z`$ (Eq. 9.7.2 of Ref.). The corresponding integral then tends to $`\frac{1}{\nu }\left(\zeta /\omega \right)^\mu `$, so that the corresponding PDF reproduces the PDF of the density of events (up to rescaling). This last form is also the asymptotic from corresponding to the behavior of Eq.(23) for large $`t`$.
Hence, the distribution $`P(x,t)`$ tends to a fully asymmetric one of index $`\alpha `$ for $`x`$ and $`t`$ large. In this case the distribution shows scaling with a scaling parameter $`\xi =x/(vt)^\alpha `$. We see that in this case the motion under the influence of a constant force is superdiffusive, so that $`x(Ft)^{1/\alpha }`$, and its dependence on the outer force is nonlinear. Thus, the model shows a behavior that differs considerably from a linear-response assumption of Refs.. This absence of a linear response regime is parallel to the CTRW-findings (see Ref. for a review) and mirrors the fact that only for normal diffusion a sweep with constant velocity and a drift under a constant force result in the same pattern of motion, see Ref. .
The case $`\alpha =1/2`$ again results in a closed expression:
$`P(x,t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{\sqrt{2\pi n}}}\mathrm{exp}\left({\displaystyle \frac{(xvn)^2}{2n}}\right)`$ (27)
$`\times {\displaystyle \frac{t}{\sqrt{2\pi n^3}}}\mathrm{exp}\left({\displaystyle \frac{t^2}{2n}}\right)dn`$
$`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{vt}{\sqrt{x^2+t^2}}}e^{vx}K_1\left(\sqrt{v^2\left(x^2+t^2\right)}\right)`$ (28)
(2.3.16.1 of Ref.). For $`v,x`$ and $`t`$ small, the corresponding distribution tends to a Cauchy-law. On the other hand, for $`t`$ large we can take approximately:
$$P(x,t)\frac{1}{\sqrt{2\pi }}\frac{\sqrt{v}t}{\left(x^2+t^2\right)^{3/4}}e^{v\left(x\sqrt{x^2+t^2}\right)}$$
(29)
The second moment of this distribution diverges, but the position of the maximum of $`P(x,t)`$, determining the typical particle position at time $`t`$, tends to grow as $`x_{\mathrm{max}}=\frac{2}{3}t^2`$ for $`t`$ large. Thus, the typical behavior of $`x(t)`$ under constant force is superlinear.
Note that in the case $`1<\alpha <2`$ the distribution of the particle’s displacement for the case $`v=0`$ will tend to a Gaussian, but in the case $`v>0`$ it still tends to a fully asymmetric Lévy one. On the other hand, in this case the distribution of the particle’s position possesses the first moment which grows linearly with time, thus the situation under $`\alpha >1`$ shows the linear response behavior. Since the second moment of the distribution is absent, the fluctuations are strong, and the width of such distribution is of the order of the typical value of $`x`$ itself.
## V The dualism between the sub- and the superdiffusive CTRW
There exists a clear dualism between the normal, subdiffusive CTRW and a superdiffusive one. The corresponding concepts are illustrated in discrete time by Fig.1, where we return to a situation discussed in Sec.2. Imagine a clock producing ticks following with frequency 1, marking the physical time of a system. Imagine a system which is triggered not by each tick of a physical clock, but follows some waiting-time distribution, $`\psi (n)`$. This means that after our random walker has jumped, the next jump will take place after $`n`$ ticks of a clock, where the number $`n`$ is chosen according to a power-law distribution, say $`\psi (n)(n+1)^{1+\gamma }`$. The number $`n`$ fluctuates strongly, so that the sequence of jumps (corresponding to a randomly decimated sequence of ticks) shows lacunae of different duration. Fig. 1a) shows a realization of such a sequence for the case $`\gamma =0.75`$. The lacuna starting in the middle of Fig. 1a) at $`t=54`$ ends at $`t=161`$. The mean number of jumps during the time $`t`$ grows sublinearly with $`t`$, namely as $`t^{3/4}`$. Let us denote the corresponding subordination transformation as time-expanding transformation (TET) of index $`\gamma `$. According to the procedure described above, the corresponding sequence does not have any intervals where the density of events is larger than one. The process subordinated to random walks under such operational time (normal CTRW) is subdiffusive.
Let us now consider the sequence of jumps of a walker as ticks marking relevant time epochs of a system (i.e. associate each jump with a tick of a physical clock). From this point of view, the ticks of initial clocks follow extremely inhomogeneously, so that the number of such ticks within a physical time unit varies according to $`p(n)(n+1)^{1+\gamma }`$. Fig.1b) illustrates this situation: Here we took 100 jumps from the realization shown in Fig.1a) and rescaled each of the corresponding time intervals to the unit length. The ticks of initial clock (shown as bars) follow inhomogeneously and show the intervals of high concentration (but no lacunae). The number of such events grows superlinearly in time. The corresponding subordination transformation will be called a ”time-squeezing transformation” (TST) of index $`\gamma `$. The process subordinated to random walks under such operational time is superdiffusive and corresponds to Lévy-flights. Note that both TST and TET are the probability distributions $`P(n,t)`$ of the operational time $`n`$ for a given physical time $`t`$, i.e. are positive, integrable functions of $`n`$.
Let us now combine the two processes. For example, let us first generate the superlinear sequence using algorithm described above (with $`p(n)(n+1)^{1+\gamma }`$) and then decimate it randomly according to the waiting-time distribution $`\psi (n)(n+1)^{1+\gamma }`$. The typical number of events during time interval $`t`$ grows in this case linearly with $`t`$, but the corresponding sequence of events is extremely inhomogeneous, showing both lacunae and accumulation intervals on all scales. This process is shown in a bar-code-like picture in Fig.2. It will be discussed in more detail in Sections 6 and 7. We can also proceed other way around, and apply the transformations in the opposite way, namely, first generating a sublinearly growing, lacunary operational time and then filling the lacunae according to a Lévy distribution. As we proceed to show, these two ways of constructing the event-time sets are not equivalent. The process subordinated to RW under such inhomogeneous operational time is a kind of a continuous-time Lévy-flight, and not a normal RW.
The example discussed above shows that transformations leading to sub- or superlinear operational time behavior (dual to each other in the sense described above) are not inverse of each other. Let us discuss a possibility of a subordination transformation transforming a Lévy-stable distribution of index $`\beta `$ (for example, a Gaussian distribution) into one with ones with index $`\gamma `$, in the sense that
$$L(ax;\gamma ,0)=_0^{\mathrm{}}n^{1/\beta }L(x/n^{1/\beta };\beta ,0)S(n,t)𝑑n,$$
(30)
where $`S(n,t)`$ is supposed to be a probability distribution of the number of steps $`n`$ done up to time $`t`$. Taking Fourier-transform of both parts of Eq.(30) and changing to a variable $`u=\left|k\right|^\beta `$ we get:
$$\mathrm{exp}(A\left|u\right|^\alpha )=_0^{\mathrm{}}e^{un}S(n,t)𝑑n$$
(31)
with $`\alpha =\gamma /\beta `$. From Eq.(31) it follows that $`S(n,t)`$ are the inverse Laplace transforms of stretched-exponentials $`\mathrm{exp}(Au^\alpha )`$. Note that according to the Bernstein’s theorem, a function $`f(x)`$ is a Laplace-transform of a probability distribution if and only if it is completely monotone (i.e. it is infinitely differentiable and $`(1)^nf^{(n)}(x)0`$ for all derivatives $`f^{(n)}`$) and $`f(0)=1.`$ The last condition is always fulfilled. Note that according to Criterion 2 discussed on p.441 of vol.IIII of Ref. a function $`f(x)=e^{\psi (x)}`$ is a completely monotone function if and only if $`\psi `$ is a positive function with a completely monotone derivative. In our case $`\psi (x)=Au^\alpha .`$ For $`0<\alpha <1`$ one has: $`g(x)=\psi ^{}(x)=A\alpha u^{\alpha 1}>0`$, and the higher derivatives (defined on the interval $`0<x<\mathrm{}`$) are: $`g^{}(x)=A\alpha (\alpha 1)u^{\alpha 2}<0,`$ $`g^{\prime \prime }(x)=A\alpha (\alpha 1)(\alpha 2)u^{\alpha 3}>0`$, $`g^{\prime \prime \prime }(x)=A\alpha (\alpha 1)(\alpha 2)(\alpha 3)u^{\alpha 3}<0`$, etc., so that $`(1)^ng^{(n)}(x)0`$, and thus the function $`g`$ is completely monotone. Thus $`S(n,t)`$ is a probability distribution (namely the one which we have found above by explicit calculation). On the other hand, for $`\alpha >1`$ the function $`g(x)`$ is not completely monotone, so that $`S(n,t)`$ is not a probability distribution. Thus, there is no random process which defines the operational time in such a way that the Lévy-flight of index $`\alpha _1`$ will be transformed into a Lévy-flight with index $`\alpha _2>\alpha _1`$. The absence of an inverse of a TST belonging to a class of subordination transformations has a deep physical interpretation: a TST is a coarse-graining procedure (see Fig.1): the information about the internal steps of the process gets lost. One can not anticipate that the transformation inverse to a coarse-graining belongs to the same class as the direct one.
Note also that the fact that the TET and TST are not inverse of each other is mirrored by the fact that within the formalism based on the fractional Fokker-Planck equations (FFPE), the first one corresponds to an additional fractional time derivative in the l.h.s. of the FFPE, while the second one is represented by a fractional spatial derivative, see Refs.. Note also that the noncommutativity mentioned above shows that the order of application of these derivatives is fixed and cannot be arbitrarily changed.
## VI The ”paradoxical” diffusion
A process subordinated to a Lévy-CTRW under TET (a time transform leading to subdiffusive CTRW) was considered in detail in Ref.. We now know that this process subordinates normal random walks under a combination of TST and TET of different indices $`\beta `$ and $`\gamma `$. The overall behavior of the process is superdiffusive for $`\gamma <\beta `$ and subdiffusive for $`\gamma >\beta `$. This is easy to understand since the scaling considerations show that the operational time grows superlinearly with physical time in the first case and that the behavior is sublinear in the second case. Note that the index $`\mu `$ of the corresponding Lévy-flight is exactly $`2\beta `$, so that this behavior is exactly the one obtained in Ref.. In the case $`\beta =\gamma `$ the operational time grows linearly with the physical one: Ref. suggests that it falls into the diffusion universality class. On the other hand this diffusion is a very special one: We will call the process subordinated to RW under such operational time a paradoxical diffusion. The random process defining an operational time stemming from a combination of TST and TET of the same index $`\gamma `$ has interesting properties: $`n`$ typically grows proportional to $`t`$; on the other hand, neither a well-defined density, nor a well-defined mean waiting-time exists.
Let us first discuss the situation mentioned in the beginning of the section: a RW subordinated to Lévy-distributed operational time, driven by a sublinear one. The PDF of the corresponding random walks has power-law tails, namely, exactly those of a Lévy-distribution of index $`\gamma `$. On the other hand, the overall width of the corresponding curve grows as $`\mathrm{\Lambda }\sqrt{t}`$. Moreover, the whole distribution scales a as a function of dimensionless displacement $`\xi =x/\mathrm{\Lambda }`$: the overall behavior is somewhat similar to one found on a polymer chain with bridges. The overall form of the function can be found using the well-known expression for $`p(n,u)`$, the Laplace-transform of the probability $`p(n,t)`$ to make exactly $`n`$ steps up to time $`t`$. Such a process corresponds to a directed motion under the same operational time as CTRW. For the ordinary renewal process one has $`p(n,u)=\frac{1}{u}\left[1\psi (u)\right]\psi ^n(u)`$, with $`\psi (u)1u^\gamma `$ . For $`u`$ small ($`t`$ large) this form corresponds to
$$p(n,u)u^{\gamma 1}\mathrm{exp}(nu^\gamma ).$$
(32)
Considering paradoxical diffusion as a process subordinated to Lévy-flights of index $`2\gamma `$ under operational time given by $`p(n,t)`$, we get for $`P(k,u)`$, the Fourier-Laplace transform of $`P(x,t),`$
$$P_\gamma (k,u)=_0^{\mathrm{}}e^{\left|k\right|^{2\alpha }n}p(n,u)𝑑n\frac{u^{\gamma 1}}{\left|k\right|^{2\gamma }+u^\gamma }.$$
(33)
The scaling nature of the distribution is immediately evident, the nature of its power-law tails follows from the asymptotic analysis for $`k`$ small: The tail of $`P_\gamma (\xi )`$ stems from those of $`L(x,2\gamma ,0)`$ and has a power-law asymptotics $`P_\gamma (\xi )\xi ^{12\gamma }`$ ($`\gamma <1`$). Note that such a distribution was obtained in Ref. as a solution of a fractional diffusion equation, describing a random process incorporating Lévy-jumps taking place under sublinear operational time. As an example let us consider the distribution $`P_{1/2}(x,t)`$, i.e. one for $`\gamma =1/2`$. This distribution has a simple analytical form, which can be obtained by an inverse Laplace-Fourier transformation of Eq.(33). The inverse Laplace transform of Eq.(33) is one given in 3.21 of Ref.and reads: $`P_{1/2}(k,t)=\mathrm{exp}(k^2t)\text{erfc}(\left|k\right|t^{1/2})`$. The inverse (cosine)-Fourier transform of this function is given by No. 10.6 of Ref. and reads:
$$P_{1/2}(x,t)=\frac{1}{2\sqrt{t}}\pi ^{3/2}\mathrm{exp}(x^2/4t)\text{Ei}(x^2/4t),$$
(34)
where $`\text{Ei}(x)`$ is the exponential integral, see Eq. 5.1.2 of Ref.. The corresponding function is a scaling function of $`\xi =`$ $`x/t^{1/2}`$; its behavior for $`\xi `$ large follows from asymptotic expansion of $`\text{Ei}(x)=\mathrm{E}_1(x)=x^1e^x[11/x+\mathrm{}]`$, so that asymptotically $`P_{1/2}(\xi )`$ shows the $`\xi ^2`$-like tail, similar to one of Cauchy-distribution. For $`\xi \mathrm{}`$ the distribution $`P_{1/2}(k,t)`$ shows week (logarithmic) singularity (following from Eq.(5.1.11) of Ref.), a sign of strong lacunarity of the corresponding operational time. The asymptotic analysis of Eq.(33) shows that such integrable singularities appear in the center of distribution for $`0<\gamma 1/2`$: the behavior for $`\xi 0`$ is given by $`P_\gamma (\xi )\xi ^{2\gamma 1}`$, for $`\gamma =1/2`$ $`P_\gamma (\xi )`$ diverges logarithmically, as we have already seen in Eq.(34).
The distribution $`P_{1/2}(\xi )`$ is plotted in Fig.3 together with the Gaussian distribution (i.e. the distribution $`P_1(\xi )`$ of the same class, the one corresponding to a normal diffusion) and with the distribution stemming from the inverse order of application of TET and TST to a simple diffusion, which is discussed in detail in the next section. All distributions are normalized in such a way that their quartiles coincide. Note that the quartiles of $`P_{1/2}(\xi )`$ are situated at $`\pm 0.841`$.
## VII Non-commutativity of time-subordination
Applying the transformations other way around, i.e. considering a process subordinated to Scher-Montroll CTRW under Lévy-time, we get a process which is different from one discussed above. Let us start from a simple example.
Let us note that the TET of index 1/2 (corresponding to an inverse Laplace-transform of the function $`e^{n\sqrt{u}}/\sqrt{u}`$) is given by
$$Q_{1/2}(n,t)=\frac{1}{\sqrt{\pi t}}e^{n^2/4t},$$
(35)
i.e. corresponds to a part of a Gaussian distribution for $`n>0`$, so that $`n`$ typically grows as $`t^{1/2}`$. The corresponding TST is given by a distribution, Eq.(20), $`R_{1/2}(T,n)=\frac{n}{2\sqrt{\pi }T^{3/2}}e^{n^2/4T}`$. The subordination of these two processes is described by a function
$`S_{1/2}(T,t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{\sqrt{\pi t}}}e^{n^2/4t}{\displaystyle \frac{n}{2\sqrt{\pi }T^{3/2}}}e^{n^2/4T}𝑑n`$ (36)
$`=`$ $`{\displaystyle \frac{2}{\pi t}}\sqrt{{\displaystyle \frac{t}{T}}}\left({\displaystyle \frac{T}{t}}+1\right)^1,`$ (37)
which is a probability distribution with the tail decaying as $`T^{3/2}`$ (as a tail of a stable distribution of index 1/2) and with the square-root singularity at zero. Note that this distribution is just a solution of a fractional Liouville equation describing directed motion under such an operational time, just like Eq.(34) is the solution of a fractional diffusion equation. This is a process subordinated to a Lévy one under sublinear time growth.
We now show that the $`Q`$\- and $`R`$\- distributions leading to the paradoxical diffusion are not commutative: An operatational time resulting from a $`RQ`$ transformation has a different distribution from one stemming from a $`QR`$-one. For example, the distribution $`S_{1/2}(T,t)`$ given by Eq.(37) is $`S_{1/2}(T,t)=QR=Q(n,t)R(T,n)𝑑n`$. Let us calculate a conjugated disrtibution, $`S_{1/2}^{}(T,t)=RQ=R(n,t)Q(T,n)𝑑n`$, one describing a process subordinated to a sublinear growth under the operational time growing according to a Lévy distribution. The distribution $`S_{1/2}^{}(T,t)`$ is given by:
$`S_{1/2}^{}(T,t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{\sqrt{\pi n}}}e^{T^2/4n}{\displaystyle \frac{t}{2\sqrt{\pi }n^{3/2}}}e^{t^2/4n}𝑑n`$ (38)
$`=`$ $`{\displaystyle \frac{2t}{\pi }}{\displaystyle \frac{1}{t^2+T^2}},`$ (39)
i.e. corresponds to a positive part of a Cauchy-distribution. Note that even such a robust scaling property of a probability distribution as a nature of its power-law tail is different from one of the conjugated counterpart.
The plausible scaling consideration here is as follows. The distribution $`Q(T,n)`$ has all moments, so that for $`n`$ large the value of $`T`$ is well-defined and is of the order of $`n^\alpha `$, $`\alpha <1`$. On the other hand, the distribution of $`n`$ as a function of $`t`$ is broad and shows a power-law tail $`P(n,t)t^{1/\alpha }(n/t^{1/\alpha })^{1\alpha }tn^{1\alpha }`$. Changing now variable from $`n`$ to $`Tn^\alpha `$ we get the asymptotics of the PDF of $`T`$in a form: $`P(T,t)tT^2`$, independently on $`\alpha `$. We note thus that the probability distribution subordinating a subliner continuous-time directed motion under the Lévy-distributed operational time of the same index has a power-law tail decaying as $`T^2`$, i.e. is similar to a Cauchy-distribution.
The process subordinated to a Gaussian RW under operational time defined by $`S_{1/2}^{}(T,t)`$ is also not a normal diffusion, but represents a marginal situation of a distribution whose second moment diverges logarithmically. The corresponding PDF shows power-law tails of a $`x^3`$ type. This PDF is given by:
$$P_{1/2}^{}(x,t)=_0^{\mathrm{}}\frac{1}{\sqrt{2\pi T}}e^{x^2/2T}\frac{2t}{\pi }\frac{1}{t^2+T^2}𝑑T.$$
(40)
Changing to a new variable $`\zeta =x^2/2T`$ and then introducing a scaling variable $`\xi =x/\sqrt{t}`$ we get the PDF $`P(x,t)`$ as a scaling function of $`\xi `$:
$$P_{1/2}^{}(\xi )=\frac{1}{\pi ^{3/2}}\left|\xi \right|_0^{\mathrm{}}\frac{\zeta ^{1/2}e^\zeta }{\zeta ^2+\xi ^4/4}𝑑\zeta .$$
(41)
For $`\xi `$ large the corresponding integral decays as $`(2/\pi )\xi ^3`$. Note that Eq.(41) can be expressed in terms of Fresnel sine- and cosine integrals, $`S(x)`$ and $`C(x)`$, so that $`P(\xi )`$ can be obtained in a closed form:
$`P_{1/2}^{}(\xi )`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }}}\{\mathrm{sin}\left({\displaystyle \frac{\xi ^2}{2}}\right)[12S\left({\displaystyle \frac{\left|\xi \right|}{\sqrt{\pi }}}\right)]`$ (43)
$`+\mathrm{cos}\left({\displaystyle \frac{\xi ^2}{2}}\right)[12C\left({\displaystyle \frac{\left|\xi \right|}{\sqrt{\pi }}}\right)]\},`$
see Eq.(2.3.7.10) of Ref.. The corresponding distribution is also plotted in Fig.3 as a dashed line. Note that the distribution shows a cusp-singularity at $`\xi =0`$. The value of $`P(\xi )`$ in this point is $`1/\sqrt{\pi }=0.564.`$.. The quartiles of this distribution are situated at $`\pm 0.621`$.
## VIII Relaxation phenomena under temporal subordination
The fact that the Lévy dynamics can follow from a temporal subordination is important if one wants to analyze the possible thermodynamical implications of the Lévy-flight transport. Imagine an ensemble of thermodynamical systems (say, Brownian particles in a harmonic potential) which was put out of equilibrium and then let relax. As discussed in Sec.2, such relaxation will lead to a stationary state corresponding to a normal equilibrium Boltzmann distribution. Since this distribution is time-independent, it would not change under temporal subordination, so that the systems with Lévy dynamics may have very ordinary thermodynamical equilibrium states and thus be described by normal Gibbs-Boltzmann entropy. The non-Boltzmann nature of the equilibrium found in Ref. was connected with the fact that the linear response was considered, as proposed by Ref., an assumption at variance with the findings of Sec. 4. Let us now discuss the relaxation to this equilibrium.
A system slightly outside of the equilibrium can be considered as evolving under the influence of the linear restoring force. In the operational time of the system (marked by the number $`n`$ of jumps) this relaxation will be described by a Fokker-Planck equation. For an overdamped particle in a harmonic potential we get, for example:
$$\frac{P}{n}=\frac{}{x}\left(\gamma kxP+D\frac{}{x}P\right)$$
(44)
Note that the values of $`\gamma `$ and $`D`$ fulfill the Einstein’s relation, $`\gamma =D/kT`$. The Green’s function of Eq.(44) has a form of a Gaussian distribution and reads:
$`G(x,n|x_0,n_0)`$ $`=`$ $`\sqrt{{\displaystyle \frac{\gamma }{2\pi D(1e^{2\gamma (nn_0)})}}}`$ (46)
$`\times \mathrm{exp}\left({\displaystyle \frac{\gamma k(xe^{\gamma (nn_0)}x_0)^2}{2D(1e^{2\gamma (nn_0)})}}\right),`$
see Sec. 5.4 of Ref. . This equation gives us e.g. the PDF at time $`n`$ in a system, in which the particles were all situated at $`x=x_0`$ at $`t=t_0`$. It is easy to see that the first two central moments $`M_1=x`$ and $`M_2=(xx)^2`$ relax exponentially to their equilibrium values, so that
$$x(n)=x_0\mathrm{exp}(\tau ^1n)$$
(47)
and
$$\sigma ^2(n)=\frac{D}{k\gamma }(1\mathrm{exp}(2\tau ^1n)),$$
(48)
being a typical pattern of relaxation of a system with only one relaxation time $`\tau =(k\gamma )^1`$. Since all higher moments of a Gaussian distribution are the combinations of the lower two, they also relax to their equilibrium values in a (multi-)exponential fashion. Let us start from the Fourier-transform of Eq.(46) and to note that under subordination
$`P(k,t)`$ $`=`$ $`{\displaystyle \mathrm{exp}\left(ikx^{^{}}e^{\gamma n}Dk^2(1e^{2\gamma n})/2\gamma \right)}`$ (50)
$`\times t^{1/\alpha }L(n/t^{1/\alpha },\alpha ,\alpha )dn.`$
Let us moreover expand the exponential term in a Taylor series in $`k`$: the coefficients of this series give the moments of the corresponding distribution. From Eq.(50) it follows then that the $`i`$-th moment is a combination of integrals of the type
$$\mathrm{\Phi }(t)=_0^{\mathrm{}}\mathrm{exp}(\lambda n)t^{1/\alpha }L(n/t^{1/\alpha },\alpha ,\alpha )𝑑n$$
(51)
with $`\lambda =m\gamma `$, $`0mi`$. Using the fact that a Laplace-transform of a fully asymmetric Lévy-distribution is a stretched exponential function, we get:
$$\mathrm{\Phi }(t)=\mathrm{exp}(A(\lambda t^{1/\alpha })^\alpha )=\mathrm{exp}(A\lambda ^\alpha t).$$
(52)
This means that the exponential relaxation under Lévy dynamics stays a simple exponential relaxation (only the corresponding relaxation time changes). For example, the first moment of the distribution (the particle’s position) still relaxes exponentially to its equilibrium value of zero. On the other hand, the dependence of the relaxation time on the outer parameters (say, temperature) entering through the values of $`\gamma `$ and $`D`$ can change considerably. Thus, the superdiffusive Lévy-flights dynamics in the force-free case can coexist with standard thermodynamics and with very simple relaxation patterns as soon as the case of a harmonic force is concerned.
Let us consider the relaxation in a harmonic potential under ”paradoxical” diffusion. Here again we can use the moment expansion, Eq.(50), and put down the expression for the characteristic function of the overall distribution:
$`P(k,t)`$ $`=`$ $`{\displaystyle \mathrm{exp}\left(ikx^{^{}}e^{\gamma n}Dk^2(1e^{2\gamma n})/2\gamma \right)}`$ (54)
$`\times S_\alpha (n,t)dn.`$
Note that the moments of the corresponding distribution are the combinations of the functions:
$$\mathrm{\Phi }(t)=_0^{\mathrm{}}\mathrm{exp}(\lambda T)S_\alpha (T,t)𝑑T.$$
(55)
Note that $`S_\alpha (n,t)`$ is a PDF of a process subordinated to a Lévy distribution under TET:
$$S_\alpha (T,t)=𝑑\tau \tau ^{1/\alpha }L_\alpha (T/\tau ^{1/\alpha },\alpha ,\alpha )Q_\alpha (\tau ,t)𝑑\tau $$
(56)
Thus, a Laplace transform of $`S`$ according to its outer time-variable is a stretched-exponential, so that
$$\mathrm{\Phi }(t)=_0^{\mathrm{}}p(\tau ,t)\mathrm{exp}(A\lambda ^\alpha \tau )𝑑\tau .$$
(57)
Let us take a Laplace-transform of this expression. Using Eq.(32) we get:
$$\mathrm{\Phi }(u)=_0^{\mathrm{}}u^{\alpha 1}\mathrm{exp}(\tau u^\alpha )\mathrm{exp}(A\lambda ^\alpha \tau )𝑑\tau =\frac{u^{\alpha 1}}{u^\alpha +A\lambda ^\alpha }.$$
(58)
For small $`u`$ (long times) this corresponds to a power-law decay of $`\mathrm{\Phi }(t)`$ of a form $`\mathrm{\Phi }(t)t^\alpha `$ for $`t>>\lambda ^1`$. Thus, the relaxation in the case of paradoxical diffusion resembles those in normal CTRW and is dominated by large lacunae. In the case when the processes are subordinated other way around, i.e. according to $`S_\alpha ^{}(T,t)`$, the decay at longer times follows the universal $`t^1`$-law: for example for $`\alpha =1/2`$ we get:
$`\mathrm{\Phi }(t)`$ $`=`$ $`{\displaystyle \frac{2t}{\pi }}{\displaystyle _0^{\mathrm{}}}\mathrm{exp}(\lambda T){\displaystyle \frac{1}{t^2+T^2}}𝑑T=`$ (59)
$`=`$ $`{\displaystyle \frac{2\lambda }{\pi }}\left[\mathrm{sin}(\lambda t)\text{ci}(\lambda t)\mathrm{cos}(\lambda t)\text{si}(\lambda t)\right],`$ (60)
see Eq.(2.3.7.11) of Ref. (here the integral sine- and cosine-functions, $`\text{si}(x)=_x^{\mathrm{}}\frac{\mathrm{sin}x}{x}𝑑x`$ and $`\text{ci}(x)=_x^{\mathrm{}}\frac{\mathrm{cos}x}{x}𝑑x,`$ are used). For $`\lambda t1`$ we get:
$$\mathrm{\Phi }(t)\frac{2}{\pi }(\lambda t)^1,$$
(61)
which asymptotic behavior is universal for all Lévy-driven CTRWs of the same index.
## IX Conclusions
A broad range of physical processes can be described as processes subordinated to a random walk under some operational time. In particular, such subordination leads to anomalous transport properties, the well-known example being the Scher-Montroll continuous-time random walks, a process in which the operational time (given by the number of steps) is sublinear in the physical time $`t`$. Here we have considered the processes subordinated to a diffusive process under operational time governed by a Lévy-distribution with index $`0<\alpha <1`$, namely the operational time superlinear in physical one. We have shown that in the absence of outer forces this subordination leads exactly to Lévy-flights. The response of such a system to weak outer force is strongly nonlinear. Interestingly enough the relaxation patterns in such systems are simpler than expected. Thus, we show that the behavior in the presence of a weak harmonic force corresponds to a simple exponential relaxation to a normal Boltzmann distribution. The combination of super- and sublinear operational times (i.e. Lévy-flights under sublinear operational time or the Scher-Montroll CTRW under Lévy-time) correspond to the ”paradoxical” diffusion, a random process which in a force-free case leads to the probability-distributions of the particle’s displacements, which show the power-law tails and lack the second moment. The width of the distribution, on the other hand, grows proportionally to the square-root of time, showing a typically diffusive behavior. Some physical implications of these findings have been discussed.
## X Acknowledgments
The author is indebted to S. Jespersen, Prof. A. Blumen and Prof. J. Klafter for fruitful discussions. Financial support by the Deutsche Forschungsgemeinschaft through the SFB 428 and by the Fonds der Chemischen Industrie is gratefully acknowledged.
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# Quantum Confinement Transition in a d-wave Superconductor
## I introduction
In recent years, due to remarkable experimental progress , the cuprate superconductors have revealed a host of mysterious phases as their chemical doping is varied. Indeed, it would seem as though these materials exhibit many of the wide range of behaviors possible for low-dimensional, highly-correlated electron systems. Centrally located within the phase diagram and adjacent to many of these puzzling regions is the d-wave superconductor. Beginning in this well-understood phase, one may develop theoretical descriptions of other, non-superconducting phases. Of particular interest are the $`T=0`$ quantum phases, both in the very underdoped and heavily overdoped regimes. The schematic situation is shown in Fig.1.
When describing a 2-dimensional superconductor, topological defects in the Cooper pair wavefunction (BCS vortices) are of particular importance. Being bosonic, once they proliferate, they condense at $`T=0`$, destroying superconductivity . In this way, a description of quantum phases with strong pairing correlations but lacking the phase coherence that is superconductivity emerge quite naturally as *vortex condensates* . If the superconductor is d-wave, there is the additional complication of low-energy quasiparticles. As recently emphasized , there is a statistical interaction between these spin-carrying quasiparticles and the vortices, making the resulting theory strongly-interacting.
Singlet-paired superconductors can be recast in a spin-charge separated form : the condensate carries all the charge but no spin, while the quasiparticles are electrically neutral with spin $`1/2`$. Most other well-understood phases of electrons (such as the Fermi liquid) are spin-charge confined. It was recently argued that many puzzling aspects of the cuprate phase diagram could be understood in terms of the fractionalization and confinement of electrons. In this approach, the regions containing the pseudo-gap (and superconducting) phase are characterized by the presence of spin-charge separation (electron fractionalization) , while the heavily overdoped regions are spin-charge confined. Between the two, a quantum confinement transition might cause critical behavior over wide regions of the high $`T_c`$ phase diagram. A key feature of the cuprates is the close proximity between a d-wave superconducting phase and a Mott insulating phase. Here, we work at half-filling and look at direct transitions between these two phases, viewing them as confinement transitions. We seek to answer two broad questions regarding the nature of such a transition.
First, in terms of phenomenology, what sorts of states might we find when separate spin- and charge- carrying excitations are confined to form electrons? As we shall see, a remarkable feature of superconductivity with one electron per unit cell is that in the dual theory, the vortices are fully frustrated. When the vortices proliferate and condense, this frustration leads to the existence of multiple vortex condensates which break spatial symmetries. In particular, we find vortex condensates which destroy superconductivity (and at half-filling, describe Mott insulators) as well as vortex condensates which preserve superconductivity. As a consequence of the vortex frustration, we find direct transitions from superconducting states to insulating states which spontaneously break rotational and/or translational symmetries, as well as the existence of superconducting states which have non-trivial spatial structure. Although we work at exactly one electron per unit cell, where the vortex theory is *fully* frustrated, it is hoped that even away from half-filling the qualitative features of our results will remain valid, in particular, the tendency toward spatial modulation near half-filling. In general, we hope that our explorations of frustrated vortex systems can yield insights into quantum phases of electrons which are complicated by the presence of competing interactions.
Second, as a specific example, we look at the critical properties of the the confinement transition between a spatially-modulated d-wave superconductor and a Mott insulator with the same broken translational symmetry. Characterized by the presence of long-ranged statistical interactions, which affect the confinement of spin and charge, this quantum critical point should have interesting universal properties. Within a special region of parameter space, we explore this transition analytically using renormalization group (RG) methods.
Before we begin to address these questions, we first lay out the basics of the model under consideration. This model was introduced in Ref. and many of its justifications and consequences can be found therein. Here, we provide only a whirlwind tour of its derivation and usefulness.
## II the model: $`Z_2`$ gauge theory
We begin by formally writing the electron creation operator as a product of *two* operators, one of which carries the spin of the electron and the other, the charge. These operators are defined with the singlet-paired superconductor in mind. If we write the Cooper pair creation operator as $`e^{i\phi _r}`$, we construct our spinless charge e boson (called a “chargon”) from the Cooper pair as:
$$b_r^{}=s_re^{i\phi _r/2}e^{i\varphi _r},s_r=\pm 1.$$
(1)
The chargon is “half a Cooper pair” in the sense that the *square* of $`b_r^{}`$ creates a Cooper pair. The neutral spin one-half particle (called a “spinon”) is obtained by removing the charge from the electron:
$$f_{r\alpha }^{}=b_rc_{r\alpha }^{}.$$
(2)
As we shall see shortly, this spinon can be thought of as a neutralized BCS quasiparticle. With these definitions, we may perform a change of variables in a suitable Hamiltonian describing electrons and Cooper pairs, resulting in a theory of chargons and spinons. However, the Hilbert space of chargons and spinons is much larger than that of electrons; for instance, the state with a single spinon but no chargons can be written down, but this state is unphysical and should be removed from the working Hilbert space. In other words, we may make this change of variables only if we additionally impose a *constraint* that the sum of the number of chargons, $`N_r`$ (canonically conjugate to the chargon phase, $`[\varphi _r,N_r^{}]=i\delta _{rr^{}}`$ ), and the number of spinons, $`\rho _r=f_{r\alpha }^{}f_{r\alpha }`$, on each site is an even integer:
$$(1)^{N_r+\rho _r}=1.$$
(3)
This constraint can be implemented within a Euclidean path integral representation, resulting in a theory of spinons and chargons coupled to a $`Z_2`$ gauge field . It should be noted that the constraint used here is not the same as Gutzwiller projection, and does not disallow doubly-occupied sites.
For an odd number of electrons per unit cell and d-wave pairing correlations, the appropriate action in the $`Z_2`$ gauge theory is:
$`S`$ $`=`$ $`S_c+S_s+S_B,`$ (4)
$`S_c`$ $`=`$ $`t_c{\displaystyle \underset{<ij>}{}}\sigma _{ij}(b_i^{}b_j+h.c.),`$ (5)
$`S_s`$ $`=`$ $`{\displaystyle \underset{<ij>}{}}\sigma _{ij}(t_{ij}^s\overline{f}_if_j+t_{ij}^\mathrm{\Delta }f_if_j+c.c.)`$ (7)
$`{\displaystyle \underset{i}{}}\overline{f}_if_i,`$
$`S_B`$ $`=`$ $`i{\displaystyle \frac{\pi }{2}}{\displaystyle \underset{i,j=i\widehat{\tau }}{}}(1\sigma _{ij}),`$ (8)
where i and j label sites on a cubic space-time lattice. The Ising gauge field minimally coupled to the chargons and spinons, $`\sigma _{ij}`$, can take values $`\pm 1`$, and $`S_B`$ is a Berry’s phase term.
One may arrive at this action by making the above-mentioned change of variables in a Hubbard-type Hamiltonian, as described in Ref. . Alternatively, this model can be taken as a starting point for describing systems with local singlet pairing correlations as well as Mott insulating tendencies. To exhibit the reasonableness of this model, consider the limits of infinite and vanishing $`t_c`$. For $`t_c\mathrm{}`$, the bosonic chargons will condense and the $`Z_2`$ gauge field will become frozen with $`\sigma _{ij}=1`$, which frees the spinons. This phase is simply the d-wave superconductor. The action reduces to S = $`S_s`$, which is just the Bogoliubov-deGennes action, with the spinons becoming the BCS d-wave quasiparticles. In the opposite limit, $`t_c0`$, the chargons are gapped into an insulating state. At $`t_c=0`$, the chargons may be trivially integrated out. The remaining action is just $`S=S_s+S_B`$. It is shown in Ref. that the partition function for this remaining spin theory is formally equivalent to that of the Heisenberg antiferromagnetic spin model. Therefore, we see the attractiveness of this model for the cuprate system, which also exhibits both superconductivity and antiferromagnetism. Many other additional properties of this action between these two limits are elucidated in Ref. , in particular, the presence of both spin-charge *confined* and *deconfined* phases .
The charge sector in Eqn. 4 is described in terms of the bosonic chargons, minimally coupled to a $`Z_2`$ gauge field. In two spatial dimensions, vortices in the boson many-body wavefunction are point-like. This allows for a particularly elegant dual description where the vortex rather than the chargon is the fundamental degree of freedom. In this duality, the condensate of chargons (the superconductor) is the vacuum of vortices; the condensate of vortices is an electronic insulator, where the chargons are gapped. Within the vortex theory, the superconductor is trivial (being just the vacuum) and is therefore a good place to plant our feet. From this vantage, we look out of the superconductor at the neighboring insulating phases. The duality transformation, on the lattice, in the presence of the $`Z_2`$ gauge field, has been explicitly implemented in Ref.. The full resulting action at half-filling is:
$`S`$ $`=`$ $`S_s+S_v+S_a+S_{CS},`$ (9)
$`S_v`$ $`=`$ $`t_v{\displaystyle \underset{<ij>}{}}\mu _{ij}cos(\theta _i\theta _j+{\displaystyle \frac{a_{ij}}{2}}),`$ (10)
$`S_a`$ $`=`$ $`{\displaystyle \frac{\kappa }{8\pi ^2}}{\displaystyle \underset{\mathrm{}}{}}|\mathrm{\Delta }\times a_{ij}2\pi \widehat{\tau }|^2,`$ (11)
$`S_{CS}`$ $`=`$ $`{\displaystyle i\frac{\pi }{4}(1\underset{\mathrm{}}{}\sigma )(1\mu _{ij})}.`$ (12)
The spinon action $`S_s`$ is unchanged. Here, $`e^{i\theta _i}`$ creates an $`\frac{hc}{2e}`$ vortex and the flux of the U(1) gauge field, $`a_{ij}`$, is the total electrical current. In particular, a flux of $`2\pi `$ through a spatial plaquette represents a charge of $`e`$. The terms $`S_v`$ and $`S_a`$ together form the usual dual vortex representation for charge $`2e`$ Cooper pairs except that here, the vortices are minimally coupled to the additional ($`Z_2`$) gauge field $`\mu _{ij}=\pm 1`$. The BCS vortex and the spinon are relative semions; upon circling a vortex, the spinon wavefunction picks up a minus sign. The term $`S_{CS}`$ is the $`Z_2`$ analog of a Chern-Simons term for the two $`Z_2`$ gauge fields and mediates this statistical vortex-spinon interaction. The spinons “see” a $`Z_2`$ flux $`_{\mathrm{}}\sigma =1`$ attached to each $`\frac{hc}{2e}`$ vortex, while the vortices see a flux of $`_{\mathrm{}}\mu =(1)^{J_f}`$. This flux attachment may be familiar to many in the context of the Quantum Hall Effect, where the gauge fields involved are for the U(1) group. Because of the anomalous “$`ff`$” terms in the action, spinon number is not conserved, and the usual Chern-Simons term cannot be used.
In the superconducting state, we are in the vacuum of vortices. The spinons see no flux and are free to propagate independently of the chargons. However, when single vortices condense, the long-range statistical interaction between the BCS vortex and the spinon drives spin-charge confinement. In the language of Ref., the condensation of $`hc/2e`$ vortices is accompanied by a condensation of the visons (vortices in the Ising field, $`\sigma `$), leading to a confined phase of electrons.
We wish here to explore in some detail the nature of this confinement transition, where the freely-propagating spin and charge excitations are “glued together” to form the electron. Aspects and implications of this quantum critical point pertaining to the high $`T_c`$ phase diagram have been introduced in Ref. . First, we will use Landau theory to find phases related to the d-wave superconductor by a second-order phase transition. Then, we will consider a special case where we recover a U(1) symmetry for the spinons and will use quantum field theory methods to extract some analytic critical properties of the transition between deconfined and confined phases.
## III dual vortex theory at half-filling
Concentrating on the vortices for the time being and working at half-filling, the dual theory for the charge sector becomes
$$_v=t_v\mathrm{cos}(\theta _i\theta _j+\frac{a_{ij}}{2})+\frac{1}{2}\left|\times \text{a}2\pi \widehat{\tau }\right|^2.$$
(13)
To obtain a low-energy effective theory, we work with a “soft-spin” model where the vortex creation operator $`e^{i\theta }`$ is replaced by a complex field $`\mathrm{\Phi }`$. In the interest of exploring the simplest case we set the charge per unit cell to be exactly e. In the dual theory, this corresponds to setting
$$(\times \text{a})_{\widehat{\tau }}=2\pi .$$
(14)
In this section, we drop fluctuations of the gauge field a, and consider a Landau mean-field approach. This is justified when the on-site repulsion between the electrons, U, is large. The vortices now see exactly $`(\stackrel{}{}\times \stackrel{}{\frac{a}{2}})_{\widehat{\tau }}=\pi `$ flux per spatial plaquette, and we are left with the two-dimensional fully-frustrated quantum XY model:
$`S={\displaystyle }d\tau \{{\displaystyle \underset{\stackrel{}{r}}{}}|_\tau \mathrm{\Phi }_\stackrel{}{r}|^2{\displaystyle \underset{\stackrel{}{r},\stackrel{}{r}^{}}{}}t_{\stackrel{}{r}\stackrel{}{r}^{}}(\mathrm{\Phi }_\stackrel{}{r}^{}\mathrm{\Phi }_\stackrel{}{r}^{}+c.c.)`$ (15)
$`+{\displaystyle \underset{\stackrel{}{r}}{}}[m^2|\mathrm{\Phi }_\stackrel{}{r}|^2+u(|\mathrm{\Phi }_\stackrel{}{r}|^2)^2]\},`$ (16)
where $`\stackrel{}{r}`$ labels sites on the 2d square lattice dual to the original electron lattice and the sign of $`t_{r,r^{}}`$ around a plaquette is $`1`$. The sites of the dual lattice are at the centers of the plaquettes of the original lattice, and in units of the lattice constant ($`a=1`$), $`\stackrel{}{r}=(x,y)`$ with x and y integers.
We proceed, following closely the work of others on the fully frustrated quantum Ising model , by choosing the gauge (to be used in the remainder of this paper) seen in Figure 2. We may diagonalize the kinetic piece of this action to find two low-energy modes, living at $`(k_x,k_y)=(0,0)`$ and $`(\pi ,0)`$, respectively. In real space, these (unnormalized) eigenvectors are:
$`\chi _\stackrel{}{r}^0`$ $`=`$ $`(1+\sqrt{2})e^{i\pi y},`$ (17)
$`\chi _\stackrel{}{r}^\pi `$ $`=`$ $`e^{i\pi x}\left[(1+\sqrt{2})+e^{i\pi y}\right],`$ (19)
(x,y integers).
For the purpose of characterizing the low-energy behavior of this vortex system, we consider fields which are linear combinations of these two low energy modes,
$$\mathrm{\Phi }(\stackrel{}{r},\tau )=\mathrm{\Psi }_0(\stackrel{}{r},\tau )\chi _\stackrel{}{r}^0+\mathrm{\Psi }_\pi (\stackrel{}{r},\tau )\chi _\stackrel{}{r}^\pi .$$
(20)
We now have two complex fields, $`\mathrm{\Psi }_0(\stackrel{}{r})`$ and $`\mathrm{\Psi }_\pi (\stackrel{}{r})`$, which describe the low-energy configurations of our vortex system. The phase transitions of the system can be explored within Ginzburg-Landau theory. The Ginzburg-Landau Hamiltonian for the two-vortex system must preserve all the symmetries of the original lattice Hamiltonian, namely: discrete $`\widehat{x}`$ and $`\widehat{y}`$ translations, rotations by $`\frac{\pi }{2}`$, and the vortex U(1) symmetry ($`\mathrm{\Phi }e^{i\alpha }\mathrm{\Phi }`$), as well as hermiticity. In terms of our two complex vortex fields, these symmetry transformations take a simpler form when expressed in terms of the fields
$$\varphi _1=\mathrm{\Psi }_0+i\mathrm{\Psi }_\pi ,\varphi _2=\mathrm{\Psi }_0i\mathrm{\Psi }_\pi ,$$
(21)
as follows:
$`T_{\widehat{x}}:`$ $`\varphi _1`$ $`\varphi _2,`$ (22)
$`T_{\widehat{y}}:`$ $`\varphi _1`$ $`i\varphi _2`$ (23)
$`\varphi _2`$ $`i\varphi _1,`$ (24)
$`R_{\frac{\pi }{2}}:`$ $`\varphi _1`$ $`e^{i\pi /4}\varphi _1`$ (25)
$`\varphi _2`$ $`e^{i\pi /4}\varphi _2,`$ (26)
$`U(1):`$ $`\varphi _a`$ $`e^{i\alpha }\varphi _a\text{(for a=1 \& 2)}.`$ (27)
Allowed terms for the action include
$`(I):`$ $`(|\varphi _1|^2)^n+(|\varphi _2|^2)^n,`$
$`(II):`$ $`(|\varphi _1|^2|\varphi _2|^2)^n,`$
$`(III):`$ $`[(\varphi _1^{}\varphi _2)^4+(\varphi _1\varphi _2^{})^4]^n,`$
(with arbitrary positive integer, n), and combinations of these terms. Expanding in powers of the fields, we take as our Landau-Ginzburg action:
$`S_{LG}={\displaystyle d^2x𝑑\tau \underset{a=1,2}{}\left[|_\mu \varphi _a|^2+r|\varphi _a|^2\right]}`$ (28)
$`+u_4(_a|\varphi _a|^2)^2+v_4|\varphi _1|^2|\varphi _2|^2v_8[(\varphi _1^{}\varphi _2)^4+h.c.],`$ (29)
where $`\tau `$ has been rescaled to set the vortex velocity $`\text{v}_v=1`$. The terms labeled by $`u_4`$ and $`v_4`$ are the only allowed quartic terms, and are invariant under independent U(1) transformations on $`\varphi _1`$ and $`\varphi _2`$. We have kept the $`v_8`$ term because it is the lowest-order term which breaks this symmetry down to the global $`U(1)`$ of Eqn. 27. This model will be employed to construct a description of various phases proximate to the d-wave superconductor within mean field theory.
We wish to characterize the various states of this vortex system. It is important to emphasize at this point that not all vortex condensates destroy superconductivity. Superconductivity is destroyed when the dual U(1) symmetry of the vortex theory (Eqn. 27) is broken. Therefore, it is possible to have non-trivial vortex condensates which are superconducting. This leads to two scenarios for the superconductor-insulator transition at half-filling. First, we may consider superconductors which are described by a vacuum of vortices; superconductivity is then destroyed when single vortices proliferate and condense (in a way which breaks the dual U(1)). Alternatively, the superconducting state could itself be a U(1)-preserving vortex condensate which then undergoes a transition which breaks the dual U(1), killing superconductivity.
In the following sections, we explore the phases of our dual vortex model using the Landau-Ginzburg action of Eqn. 29. Due to the frustration of the vortex theory with one electron per unit cell, the vortex condensates will break lattice symmetries. Some of these spatially-ordered states are superconductors and some are insulators. We will begin by describing the possible superconducting states within the dual theory (including, a *striped* superconductor), and then move on to a description of the insulating states. Ignoring charge fluctuations in the superconducting states (as we have in arriving at Eqn. 29) is not justified, and a good description of these states would require putting the charge fluctuations back in. Here, we content ourselves to characterizing the phases of our vortex system by their broken symmetries. We conclude with a summary of the possible transitions from superconductor to insulator within this mean field theory.
### A Superconductors
#### 1 Vortex Vacuums
The simplest superconducting phase is just the vortex vacuum. This is the standard BCS d-wave superconductor. Destruction of superconductivity occurs when single vortices proliferate out of the vacuum and condense, breaking the dual U(1) symmetry. The effective action for this transition is Eqn. 29.
#### 2 Paired Vortex Condensates
Condensation of single $`hc/2e`$ vortices necessarily breaks the dual U(1) symmetry (Eqn. 27), killing superconductivity. However, when *pairs* of vortices condense, the U(1) can be preserved. Consider the paired vortex condensate:
$$\varphi _2^{}\varphi _10,\varphi _1=\varphi _2=0.$$
(30)
We see that in this condensate, the dual U(1) is preserved, and the state is characterized by the phase of the condensate (setting the amplitude $`|\varphi _2^{}\varphi _1|=1`$ for simplicity),
$`\varphi _2^{}\varphi _1`$ $`=`$ $`e^{i\theta },`$ (31)
$`\theta `$ $``$ $`\theta _1(x)\theta _2(x).`$ (32)
Here, $`\theta _1`$ and $`\theta _2`$ are the phases of $`\varphi _1`$ and $`\varphi _2`$, respectively, and are still free to fluctuate. Only the combination $`\theta =\theta _1\theta _2`$ is uniform, reflecting the fact that the dual U(1) symmetry is preserved (i.e., $`\varphi _1`$ and $`\varphi _2`$ are uncondensed). The only term in the Landau-Ginzburg action which depends on $`\theta `$ is the $`v_8`$ term, giving:
$$S_v=v_8d^2x𝑑\tau \mathrm{cos}(4\theta ).$$
(33)
We see that the ground state depends on the sign of $`v_8`$:
$`v_8>0`$ $`:`$ $`\theta =n{\displaystyle \frac{\pi }{2}},`$ (34)
$`v_8<0`$ $`:`$ $`\theta ={\displaystyle \frac{\pi }{4}}+n{\displaystyle \frac{\pi }{2}},`$ (35)
with n an integer.
The spatial symmetries in Eqns.22-25, written in terms of the relative phase $`\theta `$, are:
$`T_{\widehat{x}}:\theta `$ $``$ $`\theta ,`$ (36)
$`T_{\widehat{y}}:\theta `$ $``$ $`\pi \theta ,`$ (37)
$`R_{\frac{\pi }{2}}:\theta `$ $``$ $`\theta +{\displaystyle \frac{\pi }{2}}.`$ (38)
From this we can see that the vortex condensate favored by $`v_8>0`$ breaks the lattice rotational symmetry and *one* of the two translational symmetries. We therefore associate this condensate with a stripe-type ordering: a *striped superconductor*. This state is particularly interesting given recent experimental results which suggest possible stripes in the superconducting state of $`La_{2x}Sr_xCuO_4`$ . The ground state for $`v_8<0`$ breaks all of the lattice symmetries; we identify this state with a “plaquette” order which will be made more explicit in upcoming sections when we discuss the insulating states of the vortex system. For now, we emphasize the possibility of spatially-ordered superconducting states which emerge quite naturally within our dual vortex description.
Still working in the dual description, these striped and plaquette superconductors are described by an effective theory of *one* vortex species, since the paired condensation has locked the two original vortices together: the vortex phases $`\theta _1(x)`$ and $`\theta _2(x)=\theta _1(x)\theta `$ still fluctuate within the superconducting phase, but not independently. When the remaining phase $`\theta _1`$ becomes constant over the sample, the dual U(1) is broken, and superconductivity is destroyed. Therefore, for these spatially ordered superconductors, $`S_{LG}`$ (Eqn. 29) reduces to:
$$S_v=d^2x𝑑\tau \left[|_\mu \varphi _1|^2+r|\varphi _1|^2+u(|\varphi _1|^2)^2\right].$$
(39)
It is worth noting that we have gone from a theory of a single fully-frustrated vortex to a theory of a single unfrustrated vortex via a theory of two vortices. This is possible because in a striped or plaquette superconductor, the unit cell is doubled. If one started from scratch in constructing a dual theory of these striped (plaquette) superconductors, the vortices would see $`2\pi `$ rather than $`\pi `$ flux per (doubled) unit cell and there would be only one low-energy mode.
### B Confined Insulators
When single vortices condense at half-filling, we move from the d-wave superconductor into a confined insulator. Within our dual formulation, these insulators are described by condensates which break the dual U(1) symmetry of Eqn. 27. In the case of superconductors which are vortex vacuums, because we have two vortex species, there are many ways to do this and therefore many possible single vortex condensates. We will see that these different vortex condensates correspond to different insulating states of electrons. We return to the case of the striped and plaquette superconductors after first enumerating the insulating states at the mean field level, using the action of Eqn. 29.
The most general U(1)-breaking vortex condensate is:
$`\varphi _1`$ $`=`$ $`|\varphi _1|e^{i\theta _1},`$ (40)
$`\varphi _2`$ $`=`$ $`|\varphi _2|e^{i\theta _2},`$ (41)
where $`|\varphi _1|`$, $`|\varphi _2|`$, $`\theta _1`$, and $`\theta _2`$ are all fixed real numbers. Within our dual Landau-Ginzburg model, condensing the vortices corresponds to setting $`r<0`$ and $`u_4>0`$. The signs of $`v_4`$ and $`v_8`$ then determine the ground state. For $`v_4<0`$, both vortex species acquire a non-zero amplitude $`|\varphi _1|=|\varphi _2|0`$ and their relative phase $`\theta _{12}=\theta _1\theta _2`$ is determined by the sign of $`v_8`$. On the other hand, if $`v_4>0`$, the ground states are condensates of either $`\varphi _1`$ or $`\varphi _2`$ and the sign of $`v_8`$ is irrelevant. Each of these condensates will correspond to a different insulating state of the electron system. We consider each case in turn.
#### 1 $`|\varphi _1|=|\varphi _2|0`$
These condensates are favored by $`v_4<0`$, and the relative phase ($`\theta _1\theta _2`$) is determined by the sign of $`v_8`$. Taking the magnitudes $`|\varphi _1|=|\varphi _2|=1`$, this term in the action can be rewritten as:
$$v_8[(\varphi _1^{}\varphi _2)^4+h.c.]=v_8\mathrm{cos}(4\theta _{12}).$$
(42)
In terms of this relative phase, the spatial symmetries are given by Eqns. 36-38 with the replacement $`\theta \theta _{12}`$.
####
This class of condensates is preferred by $`v_8>0`$. There are four general states, corresponding to each of the possible values of $`n`$. We see by the symmetry transformations in Eqns.36-38 that each of these states breaks the lattice rotational symmetry as well as breaking one of the two translational symmetries while leaving the other intact. On these grounds alone, we could guess that these states correspond to “stripe-like” phases. To be more concrete, we may go back to our real-space representation for the vortex field $`\mathrm{\Phi }(\stackrel{}{r})`$ in Eqn.20 and draw real-space pictures of these lattice states. The values of the fields at various points will be gauge dependent, but the location of frustrated bonds (which are places of higher energy density) is gauge-independent and therefore a good way to characterize the state of the system. This is shown for the case $`\theta _{12}=\frac{\pi }{2}`$, as an example, in Figure 3a. Investigations of this sort lead us to conclude that the four ground states of the system in this case are characterized by “stripes” of energy-density as shown in Figure 3b .
We now turn to a characterization of this system in terms of the electron degrees of freedom. Because we have broken the dual U(1) symmetry of the vortices and we are at half-filling, these states will be Mott insulators. The charge degrees of freedom live on the plaquettes of the dual lattice and are fixed at one charge of $`e`$ per dual plaquette. It has been suggested that the frustrated bonds of the dual lattice should correspond to singlet bonds of the electron system, since one expects regions of higher energy along the links where the electrons spend most of their time. This relationship between the frustrated bonds on the dual lattice and the singlet bonds on the original lattice is illustrated in Figure 4. These “striped” vortex phases then correspond to spin-Peierls (or “bond density wave”) order in the insulating electron system.
####
These condensates are favored by $`v_8<0`$. Here, however, each ground state breaks all of the discrete lattice symmetries. We proceed as above and obtain characterizations of these states in terms of the location of frustrated bonds. The result is a plaquette-like structure, as seen in Figure 5.
In terms of the electron degrees of freedom, we would like to again interpret the frustrated bonds of the dual lattice as regions where singlet-type bonds of the electron system reside. The plaquette-like structure of these vortex states may then correspond to a “plaquette RVB” state of the electron system, as shown in Figure 6.
#### 2 $`\varphi _10,\varphi _2=0`$ or $`\varphi _1=0,\varphi _20`$
These condensates are preferred in the case $`v_4>0`$. We may proceed as above in drawing real-space diagrams corresponding to these states. We find, as shown in Figure 7, that these states have vortex currents around each plaquette, of alternating sign.
In order to interpret this state, we will have to put back in the spinons which have been ignored in the previous discussion. The vortex-spinon action is:
$`S`$ $`=`$ $`S_v+S_s+S_{CS},`$ (43)
$`S_v`$ $`=`$ $`t_v{\displaystyle \underset{<i^{}j^{}>}{}}\mu _{i^{}j^{}}\mathrm{cos}(\theta _i^{}\theta _j^{}{\displaystyle \frac{a_{i^{}j^{}}}{2}}),`$ (44)
with $`\times \text{a}=2\pi ,`$ (45)
$`S_s`$ $`=`$ $`{\displaystyle \underset{<ij>}{}}\sigma _{ij}[t_{ij}^s\overline{f}_if_j+t_{ij}^\mathrm{\Delta }f_if_j]{\displaystyle \underset{i}{}}\overline{f}_if_i,`$ (46)
$`S_{CS}`$ $`=`$ $`{\displaystyle i\frac{\pi }{4}(1\underset{\mathrm{}^{}}{}\mu )(1\sigma _{ij})},`$ (47)
where $`i,j`$ label sites on the original lattice and $`i^{},j^{}`$ label sites on the dual lattice. Looking at $`S_v`$, we see that the alternating vortex currents would like to induce compensating fluctuations in either the $`a_{i^{}j^{}}`$ or $`\mu _{i^{}j^{}}`$ fields. Allowing fluctuations of the gauge field a (which describes charge fluctuations), and ignoring the coupling to the spinons, the alternating vortex currents would induce CDW order at wavevector $`(\pi ,\pi )`$. However, with a large on-site U, this state will be greatly suppressed. If we forbid charge fluctuations, we see that the alternating vortex currents will instead drive a mean field in the $`Z_2`$ gauge field:
$$\underset{\mathrm{}^{}}{}\mu _{i^{}j^{}}=(1)^{n_f}1,$$
(48)
(where $`n_f`$ is the number of spinons in the dual plaquette denoted by $`\mathrm{}`$), which corresponds to one spinon per unit cell. Unlike the previously-considered vortex condensates (with $`v_4<0`$), at the level of vortex mean field theory, this state has no broken translational symmetries. (However, we cannot rule out the breaking of symmetries by the charge and spin fluctuations.) We note that a possible candidate for this state which has one electron per until cell and uniform energy density is the antiferromagnet.
### C Summary of Vortex Theory
We have seen that our dual vortex theory describes both standard BCS and striped or plaquette d-wave superconductors as well as a host of confined insulating states. We now summarize the results of Landau theory for the transitions from the d-wave superconductor to the confined insulator at half-filling.
We consider first the transition from a vacuum of vortices (a superconductor) to a U(1)-breaking condensate of vortices (an insulator). Within mean field theory, the nature of the insulating state is determined by the signs of the coupling constants in Eqn 29, and we have the following possible *direct* transitions out of the symmetric d-wave superconductor:
$`v_4<0,v_8>0;`$ $`dSC\text{spin-Peierls,}`$ (49)
$`v_4<0,v_8<0;`$ $`dSC\text{plaquette RVB,}`$ (50)
$`v_4>0;`$ $`dSC\text{uniform state of electrons}.`$ (51)
One might hope to ascertain which of these insulating states is preferred close to a d-wave superconductor, including fluctuations beyond the mean field level, by considering the fixed points of the action in Eqn.29. In particular, we see that the sign of $`v_4`$ determines whether we enter one of the states of broken translational symmetry (spin-Peierls or plaquette RVB), or the state with uniform energy-density (possibly the antiferromagnet). The work of Blagoeva on the theory of two-component complex fields with these (and other) couplings gives a stable fixed point at $`v_4<0`$, to order $`ϵ^2`$ ($`ϵ=4D`$, $`D=d+1`$, in d spatial dimensions). This suggests that the transition dSC $``$ spin-Peierls would be preferred over dSC $``$ uniform state. This is tantalizing given the experimental evidence for intervening “stripey” phases between the superconducting and antiferromagnetic phases in the cuprates .
In the case of the striped and plaquette superconductors, when the single vortex species in Eqn. 39 condenses, superconductivity in these states is destroyed and we enter a confined insulating state. Because the relative phase $`\theta _1\theta _2`$ is already fixed within these superconductors, we see from our above analysis of the insulating phases that the insulating state is pre-determined. The striped superconductor (with $`\theta =n\pi /2`$) enters the spin-Peierls insulator, and the plaquette superconductor (with $`\theta =\pi /4+n\pi /2`$) enters the plaquette-RVB insulator. In other words, these spatially-ordered superconductors make transitions into insulating states with the same broken spatial symmetries:
$`v_8>0;\text{striped SC}`$ $``$ spin-Peierls, (52)
$`v_8<0;\text{plaquette SC}`$ $``$ plaquette-RVB. (53)
In the preceding section, we have considered states of electron systems at half-filling near a d-wave superconductor within a dual formulation in terms of vortices. Each phase is characterized by a dual (vortex) order parameter. At one electron per unit cell, the frustration of the vortex theory manifests itself in spontaneously broken spatial symmetries. Exploiting the fact that the vortex order parameters break spatial symmetries has helped us identify these vortex phases with more familiar phases of electrons (such as the spin-Peierls state), as well as phases like the striped superconductor. The dual formulation shows us the enhanced chance for striped superconductors near half-filling. In the next section, we will add back in the spinons (and along with them, their long-ranged statistical interaction with the vortices), and extract information about the critical properties of the confinement transition using field theory methods.
## IV confinement transition
Having explored the vortex sector of the theory with one electron per unit cell, we now wish to put the spinons back in and address the critical properties of the confinement transition. Because we will continue to work at half-filling, the confined states of electrons will be Mott insulators. While the theory of vortices and spinons coupled to $`Z_2`$ gauge fields may in principle be numerically accessible, the action suffers from the notorious fermion sign problem. Here, we discuss a special case which will turn out to be accessible to perturbative RG calculations.
Focusing on the spinon Hamiltonian (and dropping the $`Z_2`$ gauge field for the time being):
$$H_s=\underset{<rr^{}>}{}[t_{rr^{}}^sf_r^{}f_r^{}+t_{rr^{}}^\mathrm{\Delta }(f_rf_r^{}+c.c.)],$$
(54)
we choose the special case:
$$t^s=|t^\mathrm{\Delta }|t,$$
(55)
with nearest-neighbor d-wave pairing amplitude:
$$\begin{array}{cc}t_{\stackrel{}{r},\stackrel{}{r}\pm \widehat{x}}^\mathrm{\Delta }& =+t,\\ t_{\stackrel{}{r},\stackrel{}{r}\pm \widehat{y}}^\mathrm{\Delta }& =t.\end{array}$$
(56)
Following Affleck , we introduce the fields
$$\left(\begin{array}{c}d_r\\ d_r^{}\end{array}\right)=\{\begin{array}{cc}e^{i\frac{\pi }{8}\sigma _y}\left(\begin{array}{c}f_r\\ f_r^{}\end{array}\right),\hfill & \text{ for }y\text{ even, }\hfill \\ (i\sigma _y)e^{i\frac{\pi }{8}\sigma _y}\left(\begin{array}{c}f_r\\ f_r^{}\end{array}\right),\hfill & \text{for }y\text{ odd,}\hfill \end{array}$$
(57)
(where $`\sigma _y`$ is the usual Pauli matrix), the spinon Hamiltonian becomes:
$$H_s=\underset{<rr^{}>}{}t_{rr^{}}(d_{r\alpha }^{}d_{r^{}\alpha }+h.c.),$$
(58)
with
$$t_{rr^{}}=\{\begin{array}{cc}t\hfill & \text{for }y\text{ and }y^{}\text{ even,}\hfill \\ t\hfill & \text{else.}\hfill \end{array}$$
(59)
This is the Hamiltonian of fermions hopping in 2d in the presence of $`\pi `$ flux per plaquette. We have succeeded in finding a Hamiltonian for the spin sector which has a conserved fermion number. The original theory (Eqn.4) can now be written in terms of these $`d`$ fermion fields, the chargons and the $`Z_2`$ gauge field. Following a transformation which can get rid of the Berry’s phase term this theory can be modeled numerically with no fermion sign problem. Here, we instead proceed to a low-energy continuum Hamiltonian for the spin sector. To this end, we diagonalize $`H_s`$ to find two Dirac points. These are the usual d-wave quasiparticle nodes at $`(k_x,k_y)=(\pm \frac{\pi }{2},\pm \frac{\pi }{2})`$ except that due to the $`\pi `$ flux per plaquette, we have doubled the unit cell and halved the Brillouin zone; it now contains only two of these nodes, which we denote $`\stackrel{}{K}_1`$ and $`\stackrel{}{K}_2`$. In terms of long-wavelength fields living at these two nodes,
$$d_{j\alpha }(\stackrel{}{x})\psi _{1j\alpha }(\stackrel{}{x})e^{iK_1x}+\psi _{2j\alpha }(\stackrel{}{x})e^{iK_2x},$$
(60)
(where $`j=1,2`$ labels the sublattice), the continuum Hamiltonian is:
$`H_s`$ $`=`$ $`{\displaystyle d^2x\text{v}_s\psi _{1\alpha }^{}[\tau _1(i_x)+\tau _2(i_y)]\psi _{1\alpha }}`$ (62)
$`+\text{v}_s\psi _{2\alpha }^{}[\tau _2(i_x)+\tau _1(i_y)]\psi _{2\alpha },`$
where,
$`\tau _1`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\tau _x+\tau _z),`$ (63)
$`\tau _2`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\tau _x\tau _z).`$ (64)
Here, $`\stackrel{}{\tau }`$ acts in the sublattice space, and we have rotated the x and y axes at each node by $`45^{}`$. Just as the Hamiltonian for the $`d`$ fermions was diagonal in the spin label, so is this one, and we are left with a theory of four species of Dirac fermions. Note that the spinon characteristic velocity, $`\text{v}_s`$ is isotropic in space because of our choice $`t_s=|t_\mathrm{\Delta }|`$.
Defining Dirac matrices in 2+1 dimensions:
$$\begin{array}{cc}\text{at node}\stackrel{}{K}_1:& \text{at node}\stackrel{}{K}_2:\\ \gamma _0\tau _y,& \gamma _0\tau _y,\\ \gamma _1\tau _2,& \gamma _1\tau _1,\\ \gamma _2\tau _1,& \gamma _2\tau _2,\end{array}$$
(65)
$$(\gamma _\mu )^{}=\gamma _\mu ;\{\gamma _\mu ,\gamma _\nu \}=2\delta _{\mu \nu }\text{(at each node)},$$
(66)
we proceed to the Euclidean Lagrangian density:
$`_s`$ $`=`$ $`\overline{\psi }_a[\gamma _0_0+\text{v}_s\gamma _i_i]\psi _a,`$ (67)
$`\overline{\psi }`$ $``$ $`\psi ^{}\gamma _0.`$ (68)
The fields $`\psi _a`$ have two components (corresponding to the sublattice label), and summation conventions on the number of species $`a[1,4]`$ (one for each spin at each of the two nodes) and the spatial dimension $`i[1,2]`$ are in use.
We have succeeded in writing a low-energy effective theory for the spin sector which is just that of four species of two-component Dirac fermions. We may now write down a full low energy effective theory where, since the spinon and vortex sectors each display U(1) symmetries:
$`\psi _a`$ $``$ $`e^{i\alpha _\psi }\psi _a,`$ (69)
$`\varphi `$ $``$ $`e^{i\alpha _\varphi }\varphi ,`$ (70)
we may implement the statistical spinon-vortex interaction using U(1) (rather than $`Z_2`$) Chern-Simons gauge fields, $`\text{A}^\varphi `$ and $`\text{A}^\psi `$. We proceed to a field theory modeling the confinement transition between a spin-charge separated (d-wave) superconductor and a spin-charge confined Mott insulator. For simplicity, we consider the vortex theory with only one species (Eqn.39), and thereby consider transitions out of the striped/plaquette superconductor given in Eqns. 52 \- 53. The low-energy effective action is:
$`S`$ $`=`$ $`{\displaystyle d^2x𝑑\tau [_s+_v+_{CS}+_{vs}]},`$ (71)
$`_s`$ $`=`$ $`\overline{\psi }_a(\overline{)}ig\overline{)}A^\psi )\psi _a+\kappa \overline{\psi }_a(\gamma _i_iig\gamma _iA_i^\psi )\psi _a,`$ (72)
$`_v`$ $`=`$ $`|(_\mu igA_\mu ^\varphi |^2+m^2|\varphi |^2+u_0(|\varphi |^2)^2,`$ (73)
$`_{𝒞𝒮}`$ $`=`$ $`iϵ_{\mu \nu \lambda }A_\mu ^\varphi _\nu A_\lambda ^\psi ,`$ (74)
$`_{vs}`$ $`=`$ $`v_0\psi _a^{}\psi _a|\varphi |^2,`$ (75)
with $`a[1,N=4],\kappa ={\displaystyle \frac{\text{v}_s}{\text{v}_v}}1,`$ (76)
where $`\kappa `$ is a measure of the velocity anisotropy between vortices and spinons and will be treated as a perturbation. We have added the term $`_{vs}`$ in the interest of including all possible relevant interactions. The Chern-Simons term causes a vortex taken around a spinon to acquire a phase of:
$$\varphi e^{ig{\scriptscriptstyle \stackrel{}{A}^\varphi 𝑑\stackrel{}{l}}}\varphi =e^{ig^2}\varphi ,$$
(77)
and likewise for a spinon after encircling a vortex:
$$\psi e^{ig{\scriptscriptstyle \stackrel{}{A}^\psi 𝑑\stackrel{}{l}}}\psi =e^{ig^2}\psi ,$$
(78)
so that the full statistical interaction is achieved when:
$$g^2=\pi =2\pi \alpha ;\alpha =\frac{1}{2}$$
(79)
(where $`\alpha `$ is the so-called “statistics angle” and is equal to $`1/2`$ since the vortex and the spinon are relative semions). The theory as written neglects charge fluctuations, which is not justified within the superconducting phase. The full vortex theory would include an additional minimal coupling to a gauge field a . As seen in the dual XY model, this coupling causes runaway flows, and is probably best modelled numerically. At this point, we leave out the gauge field a and its attendant problems, but we will revisit this question shortly.
When the vortex Lagrangian is taken through criticality ($`m^2<0`$), the statistical interaction, mediated by the gauge fields $`A_\mu ^\varphi `$ and $`A_\mu ^\psi `$, will drive spin-charge confinement. Here, we seek the effect of these statistics on critical properties of the system. In particular, we wish to calculate $`\beta `$ functions for the couplings $`u_0`$, $`v_0`$, $`\kappa `$, and $`g`$, as well as the anomalous dimensions of the vortex and spinon fields.
We work in $`D=d+1=3`$ dimensions (indeed, our Chern-Simons flux attachment is not well-defined in higher dimensions), and define dimensionless couplings:
$`u`$ $`=`$ $`\mathrm{\Lambda }^1u_0K_{D=3},`$ (80)
$`v`$ $`=`$ $`v_0K_{D=3},`$ (81)
where factors of $`K_D=[2^{D1}\pi ^{D/2}\mathrm{\Gamma }(D/2)]^1`$ have been put in for later convenience. The bare propagators in the Landau gauge are:
fermions: $`G_0^\psi `$ $`={\displaystyle \frac{i\overline{)}k}{k^2}},`$ (83)
vortices: $`G_0^\varphi `$ $`={\displaystyle \frac{1}{k^2}},`$ (84)
gauge fields: $`S_0^{\mu \nu }`$ $`={\displaystyle \frac{ϵ^{\mu \nu \lambda }k^\lambda }{k^2}}=A_\mu ^\psi A_\nu ^\varphi ,`$ (86)
$`A^\varphi A^\varphi =A^\psi A^\psi =0.`$
(The fermion propagator is diagonal in the label $`a`$, so we have suppressed this index.)
For the $`\beta `$-functions we find, to lowest non-vanishing order (1 loop):
$`{\displaystyle \frac{du}{dl}}`$ $`=`$ $`u10u^2+({\displaystyle \frac{N}{3}}+C\kappa )v^2+\mathrm{},`$ (87)
$`{\displaystyle \frac{dv}{dl}}`$ $`=`$ $`4uv+\mathrm{},`$ (88)
$`{\displaystyle \frac{dg^2}{dl}}`$ $`=`$ $`0,`$ (89)
$`{\displaystyle \frac{d\kappa }{dl}}`$ $`=`$ $`0+\mathrm{}.`$ (90)
We expect that at higher orders, $`g`$ will enter into $`\frac{du}{dl}`$ and $`\frac{dv}{dl}`$ non-trivially, but that $`g`$ itself should not renormalize at any order, following the argument given by Semenoff et al. . The one-loop RG equations for $`u`$ and $`v`$ have a stable solution at $`v=0`$, so that the theory decouples into separate spinon and vortex theories. At this order, since the spinon and vortex sectors decouple, we may ignore the Chern-Simons gauge fields (effectively taking $`g=0`$) and include the effects of charge fluctuations by using the full dual XY model for the vortex sector:
$`_v`$ $`=`$ $`|(_\mu ie_0a^\mu )\varphi |^2+{\displaystyle \frac{1}{2}}|\stackrel{}{}\times \stackrel{}{a}|^2`$ (92)
$`+m^2|\varphi |^2+u_0(|\varphi |^2)^2.`$
Recently, much work has gone into tackling the critical properties of the $`e0`$ model , and we may use these results.
To first order, then, we find a *fixed line*, parameterized by values of the statistics angle $`\alpha `$ (or equivalently, the coupling $`g`$). At lowest non-vanishing order, this line is given by:
$`u^{}`$ $``$ $`u_{dual}^{},`$ (93)
$`v^{}`$ $``$ $`0,`$ (94)
$`(g^2)^{}`$ $`=`$ $`g^2=\pi ,`$ (95)
$`e^{}`$ $``$ $`e_{dual},`$ (96)
$`\kappa ^{}`$ $``$ $`\kappa ,`$ (97)
where, by $`u_{dual}^{}`$ and $`e_{dual}`$, we mean the values of these couplings at the fixed point of the dual XY model.
In order to see whether spinon-vortex velocity anisotropy grows, we need to take the $`\beta `$ function for $`\kappa `$ to its lowest non-vanishing order, which is two loops. The result is:
$$\frac{d\kappa }{dl}=\frac{31}{240}\frac{g^4}{\pi ^2}\kappa +\mathrm{}.$$
(98)
Since the system flows toward $`\kappa =0`$, it is legitimate to treat this term as a perturbation, and the theory becomes “relativistic” at the critical point.
We proceed by calculating the anomalous dimensions of the spinon and vortex fields, to lowest order, near the critical point. To that end, we consider the self-energies:
$`[G^\varphi (k)]^1`$ $`=`$ $`[G_0^\varphi (k)]^1+\mathrm{\Sigma }^\varphi (k),`$ (99)
$`[G^\psi (k)]^1`$ $`=`$ $`[G_0^\psi (k)]^1+\mathrm{\Sigma }^\psi (k).`$ (100)
Near the critical point, the anomalous dimensions are given by:
$`G^\varphi (k)`$ $``$ $`{\displaystyle \frac{1}{|k|^{2\eta _\varphi }}},`$ (102)
$`G^\psi (k)`$ $``$ $`{\displaystyle \frac{i\overline{)}k}{|k|^{2\eta _\psi }}}`$ (103)
(up to additive constants). Working at the fixed point $`(u,v,g^2)=(u^{}=u_{dual}^{},v^{}=0,g_{}^{2}{}_{}{}^{}=\pi )`$ and calculating the spinon and vortex self-energies to two loops in 3 dimensions, we find:
$`\eta _\varphi `$ $`=`$ $`\eta _{dual}{\displaystyle \frac{4}{3}}{\displaystyle \frac{(g^4)^{}}{16\pi ^2}}N+\mathrm{},`$ (104)
$`\eta _\psi `$ $`=`$ $`{\displaystyle \frac{1}{3}}{\displaystyle \frac{(g^4)^{}}{16\pi ^2}}+\mathrm{}.`$ (105)
Since we are in the case with one vortex species, we may take the numerical results of Hove and Sudbø for the anomalous dimension of the vortex field in the full dual XY model in $`D=3`$: $`\eta _\varphi 0.24`$. After plugging in $`N=4`$, and $`g_{}^{2}{}_{}{}^{}=\pi `$ in our result, we find:
$`\eta _\varphi 0.24{\displaystyle \frac{1}{3}}0.57,`$ (106)
$`\eta _\psi {\displaystyle \frac{1}{48}}0.02.`$ (107)
These critical exponents may reveal themselves in many quantities. In particular, the spectral function as probed by ARPES and the spin-spin correlations probed by NMR or neutron scattering. Within our theory, the low-energy electron correlator decouples into chargon and spinon pieces for $`g0`$:
$$c(x)c^{}(0)=b(x)b^{}(0)f(x)f^{}(0).$$
(108)
These correlators will exhibit anomalous dimensions $`\eta _b`$ and $`\eta _f`$, which can be expanded perturbatively around $`g^{}=0`$:
$`\eta _b`$ $`=`$ $`\eta _{XY}+C_b^{(2)}(g^{})^2+\mathrm{},`$ (109)
$`\eta _f`$ $`=`$ $`\eta _\psi =C_f^{(2)}(g^{})^2+\mathrm{},`$ (110)
where we have calculated $`C_f^{(2)}0.03`$. The anomalous dimension for the 3d XY model (appropriate for one vortex species) has been calculated by Hasenbusch and Török using Monte Carlo methods ; they find $`\eta _{XY}0.038`$. The anomalous dimension will also enter into the spin-spin correlation function. Within our model, it looks as though vertex correction diagrams will not contribute as much near the critical point as the direct $`[G_\psi ]^2`$ term.
## V conclusions
In this paper, we have used a gauge theory of strongly-interacting electrons to explore the regions near the superconducting state in the high-$`T_c`$ cuprates. This gauge theory exhibits spin-charge separated and spin-charge confined phases. We have seen that the presence of one electron per unit cell has profound implications for the regions near the superconducting state. Within a dual description, half-filling of electrons corresponds to fully-frustrated vortices, leading to a spontaneous breaking of translational symmetries in the electron system. From this, we have seen the possibility of striped superconductivity as well as a host of confined insulators descending from d-wave superconducting phases. We have then used Chern-Simons methods to calculate lowest-order critical properties of the confinement transition between these phases. Because we have worked at half-filling of electrons throughout, our results are of particular relevance to the undoped cuprate materials, which may be spin-charge confined. However, we also hope that the flavor of our results may be of interest in the heavily overdoped materials, where the confinement of spin and charge may result in a Fermi liquid phase.
We are grateful to Leon Balents, Patrick Lee, and Doug Scalapino for insightful discussions. This research was generously supported by the NSF under Grants DMR-97-04005, DMR95-28578 and PHY94-07194.
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# PHOTOMETRIC REDSHIFTS AND GRAVITATIONAL TELESCOPES
## 1 Introduction
Photometric redshifts (hereafter $`z_{phot}`$s) are a promising technique in deep universe studies. The interest for this technique has recently increased with the development of large field and deep field surveys, in particular the HDF. The relatively high number of objects accessible to photometry allows to enlarge the spectroscopic sample towards the faintest magnitudes. For this reason, $`z_{phot}`$s are also important when using lensing clusters as Gravitational Telescopes (GTs), for the study of both the background sources and the lenses.
One of the most widely used $`z_{phot}`$ techniques is the SED fitting procedure. The observed photometric spectral energy distributions (SEDs) are compared to those obtained from a set of template spectra, using the same photometric system. The method is based on the detection of strong spectral features, and it has been largely applied on HDF studies $`^\mathrm{?}`$ $`^\mathrm{?}`$ $`^\mathrm{?}`$ $`^\mathrm{?}`$ $`^\mathrm{?}`$ $`^\mathrm{?}`$ $`^\mathrm{?}`$ $`^\mathrm{?}`$. The examples presented in this paper have been produced using our public code called hyperz, which adopts a standard SED fitting method, but most results should be completely general when using other $`z_{phot}`$ tools. hyperz is presented in a recent paper by Bolzonella et al. (2000) $`^\mathrm{?}`$, and it is available on the web at http://webast.ast.obs-mip.fr/hyperz. When applying hyperz to the spectroscopic samples available on the HDF, the uncertainties are typically $`\delta z/(1+z)0.1`$, and this value gives an idea of the expected accuracy of $`z_{phot}`$s for the purposes of this paper.
## 2 Photometric redshifts and background sources
The idea is to take benefit from the amplification factor in GTs to study the properties of the distant background population of lensed galaxies. The typical amplification ranges between 1 to 3 magnitudes in the cluster core. In principle, GTs are useful to build-up and study an independent sample of high-$`z`$ galaxies, which complements the large samples obtained in standard field surveys. The advantage is that this sample is less biased in luminosity than the field.
### 2.1 Redshift distribution of faint galaxies
One of the main goals of the GT is to determine the $`z`$ distribution of a very faint subsample of high-$`z`$ lensed galaxies, invisible otherwise, and this can be achieved using $`z_{phot}`$s. In order to prevent the biases towards or against a particular type of galaxy or redshift domain, $`z_{phot}`$ shall be computed from broad-band photometry using a large wavelength interval, from B (U when possible) to K. This allows also to reduce the errors on $`z_{phot}`$ (see Bolzonella et al. 2000 for a detailed discussion).
We have obtained the (photometric) N($`z`$) distribution of arclets in several well known clusters (A2390, A370, Cl2244-02, AC114,…). Figure 1 displays a recent example, the $`z_{phot}`$ distribution for different source samples in MS1008-1224$`^\mathrm{?}`$, corresponding to different limits in magnitude. In this case, $`z_{phot}`$s were computed from VLT BVRI (FORS) and JK’ (ISAAC) public data obtained during the Science Verification phase. The sample includes 559 sources located on the central 2.5 arcminute field of ISAAC, excluding obvious cluster members. The typical number of high-$`z`$ sources found in the inner 1’ radius region of the cluster is $`30`$ to 50 at $`1z7`$, for a photometric survey performed with the HST and different 4m class telescopes. Clusters with well constrained mass distributions enable to recover precisely the $`N(z_{phot})`$ distribution of lensed galaxies, by correcting the relative impact parameter on each redshift bin.
### 2.2 Optimization of spectroscopic redshift surveys
An interesting issue for $`z_{phot}`$ when using GTs for the spectroscopic study of faint amplified sources is the optimization of the survey, that is selecting the best spectral domain in the visible or near-IR bands. This means in practice to produce a criterion based in $`z_{phot}`$ to discriminate between objects showing strong spectral features in the optical and in the near-IR.
An additional benefit of $`z_{phot}`$ is that this technique efficiently contributes to the identification of objects with ambiguous spectral features, such as single emission lines. An example of this is given in a recent paper by Campusano et al. (in preparation, see also Le Borgne et al., this conference).
### 2.3 Identification and study of very high-$`z`$ sources
The signal/noise ratio in spectra of amplified sources and the detection fluxes are improved beyond the limits of conventional techniques, whatever the wavelength used for this exercise. An example is the ultra-deep MIR survey of A2390$`^\mathrm{?}`$, and the SCUBA cluster lens survey$`^\mathrm{?}`$$`^\mathrm{?}`$. Number of $`z\stackrel{>}{}4`$ lensed galaxies have been found recently, and these findings strongly encourage our approach$`^\mathrm{?}`$, $`^\mathrm{?}`$, $`^\mathrm{?}`$, $`^\mathrm{?}`$. Cluster lenses are the natural way to search for primeval galaxies, in order to constraint the scenarios of galaxy formation.
High-$`z`$ lensed sources could be selected close to the appropriate critical lines, and identified using $`z_{phot}`$ criteria. $`z_{phot}`$s are computed from broad-band photometry on a large wavelength interval, from B (U when possible) to K.
Whatever the $`z_{phot}`$ method used, a crucial test is the comparison between the photometric and the spectroscopic redshifts obtained on a restricted subsample of objects. Thanks to the magnification factor, cluster lenses could be used to enlarge the training sample towards the faintest magnitudes.
### 2.4 Combining photometric and lensing redshifts
Lensing inversion and $`z_{phot}`$ techniques produce independent probability distributions for the redshift of amplified sources. Therefore, the combination of both methods provides an alternative way to determine the redshift distribution of the faintest high-$`z`$ sources. Figure 2 displays the results for 4 sources in A2390, with $`z_{phot}`$ computed from BgVrRiIJK’ photometry$`^\mathrm{?}`$. When comparing the $`z_{phot}`$ and lensing redshift ($`z_{lens}`$) values for a subsample of 98 arclets in the core of A2390, all selected according to morphological criteria (minimum elongation and right orientation are requested), we find that about $`60\%`$ of the sample have $`|z_{phot}z_{lens}|0.25`$. The discrepancy mostly corresponds to sources with $`z_{phot}\stackrel{>}{}2`$, but it must be noticed that the lensing inversion technique is characterized by a trend against the identification of high-$`z`$ images, whereas $`z_{phot}`$ does not show this trend. This behaviour is expected as a result of the relative low sensitivity to $`z`$ of the inversion method for high-$`z`$ values. Thus, in general, the $`z_{phot}`$ determination is more accurate than the $`z_{lens}`$. Nevertheless, the combination of both distributions is particularly useful when a degenerate solution appears using $`z_{phot}`$, as shown in Figure 2.
## 3 Photometric redshifts and cluster lenses
### 3.1 Identification of multiple images
Cluster lenses are useful only when their mass distribution is highly constrained by multiple images, revealed by HST and multicolor photometry. The $`z_{phot}`$ technique is useful to identify objects with similar SEDs within the errors, thus compatible with a multiple image configuration. $`z_{phot}`$s can be used advantageously when data span a large wavelength range, and particularly if near-IR photometry is available, because this allows to obtain accurate $`z_{phot}`$ in the sensitive region of $`0.8z2`$.
### 3.2 Scaling the mass in weak lensing analysis
$`z_{phot}`$s are particularly useful when deriving the mass from weak shear analysis: they are used to eliminate cluster and foreground galaxies from the analysis, and to scale the lensing distance modulus in order to compute the mass from the gravitational convergence. The average convergence $`\kappa \mathrm{\Sigma }/\mathrm{\Sigma }_{cr}`$, which corresponds to the ratio between the surface mass density and the critical value for lensing, may be obtained as a function of the radial distance $`\theta `$ using different methods (see Mellier 1999$`^\mathrm{?}`$ for a review). The mass within an aperture $`\theta `$ is given by
$$M(<\theta )=\kappa (<\theta )\mathrm{\Sigma }_{cr}\pi \left(\theta D_{ol}\right)^2=\kappa (<\theta )\theta ^2\frac{c^2}{4G}\frac{D_{ls}}{D_{os}D_{ol}}^1$$
where $`D_{ij}`$ are the angular size distances between the cluster ($`l`$), the observer ($`o`$) and the source ($`s`$), and $`\kappa (<\theta )`$ is the averaged convergence within the radius $`\theta `$. Using the $`N(z_{phot})`$ computed through a suitable filter set, the mean value of $`\frac{D_{ls}}{D_{os}D_{ol}}`$ can be computed, thus leading to a fair estimate of the mass. This method has been recently applied to the lensing clusters MS1008-1224$`^\mathrm{?}`$ and A2219$`^\mathrm{?}`$ (see also Abdel Salam and Gray, this conference). Because of the small projected surface across the redshift space, the effective surface which is “seen” through a cluster lens is relatively small, thus producing a strong variance from field to field. Obtaining the $`N(z_{phot})`$ distribution for each cluster could help to improve the mass determination. Nevertheless, the distortion on the $`N(z_{phot})`$ distribution itself depends on the mass distribution (sect. 2.1). An iterative process is needed to reduce the error bar on the mass, which is $`30\%`$ without this second order correction$`^\mathrm{?}`$.
### 3.3 Identification and Characterization of lensing structures
$`z_{phot}`$s are useful to improve the detection of clusters in wide-field surveys, and to identify the visible counterpart of complex lenses. It has been shown that including such a $`z_{phot}`$ technique in an automated identification algorithm allows to improve significantly the detection levels for clusters$`^\mathrm{?}`$, whatever the algorithm used$`^\mathrm{?}`$ $`^\mathrm{?}`$ $`^\mathrm{?}`$. In general, the $`S/N`$ is improved by a factor of 3 to 6 up to $`z1`$, depending on the redshift and magnitude limits. The detection efficiency in the $`0.8z2.2`$ region is improved only when using a $`z_{phot}`$ selection based on optical and near-IR filters. These comments also apply to multiple, complex and/or dark lenses, where $`z_{phot}`$s allow to identify the main lensing structures. Examples of composite lenses recently identified by $`z_{phot}`$ are MS1008-1224$`^\mathrm{?}`$, where a secondary lens exists at $`z0.9`$, and the multiple-quasar fields Q2345+007$`^\mathrm{?}`$ and the Cloverleaf$`^\mathrm{?}`$.
Cluster members could be also selected using $`z_{phot}`$ criteria. The present version of hyperz is also able to display the probability of each object to be at a fixed redshift. This is useful when looking for clusters of galaxies at a given (or guessed) redshift. In this way, the number density and luminosity density distributions of cluster galaxies can be computed, and $`M/L`$ ratios could be estimated.
## 4 Conclusions and Perspectives
Concerning the lensing clusters, $`z_{phot}`$s appear as an essential tool for mass calibration in weak shear studies (see the recent papers on MS1008 $`^\mathrm{?}`$ and A2219 $`^\mathrm{?}`$). $`z_{phot}`$s allow to improve the identification of the “visible” counterpart of lensing structures, to determine their redshifts and to measure $`M/L`$ ratios of clusters, groups, etc.
$`z_{phot}`$s allow to optimize the spectroscopic surveys of faint lensed sources (visible versus near-IR domains), and to identify ambiguous spectral features (emission lines). Before the recent VLT survey on AC114 (Le Borgne et al. this conference), all spectroscopically confirmed lensed sources were “bright” ($`M_B\stackrel{<}{}21`$). Because a redshift accuracy of $`\delta z0.1(1+z)`$ is enough for most applications, $`z_{phot}`$s allow to go further in magnitude when deriving the statistical properties of faint sources. Besides, combining $`z_{phot}`$ and lensing inversion techniques provides with an alternative way to determine the redshift distribution of sources. A future development should be the study of the systematics and biases introduced by GTs when they are used to access the distant population of faint sources. The typical uncertainty in the amplification factor is $`\mathrm{\Delta }m`$ 0.3 magnitudes, a value which is similar to model uncertainties when deriving intrinsic luminosities and SFRs of background sources with relatively well constrained SED ($`30\%`$ accuracy). Only well known lenses are actually useful as GTs.
Conversely, lensing clusters could be considered as a tool to calibrate $`z_{phot}`$s beyond the reach of standard spectroscopy, up to the faintest limits in magnitude, thus allowing to extend the training set for $`z_{phot}`$s. The sample of lensing clusters available has to be enlarged, in order to obtain a weakely biased image of the averaged properties of sources along different lines-of-sight. In order to sample a statistically significant field at $`z2`$ in the strong amplification domain (close to the corresponding caustic lines), we need to study about 10 different and well-known cluster lenses. The Ultra-Deep Photometric Survey of cluster lenses and the subsequent spectroscopic follow up of sources constitute a well defined program for 10m ground-based telescopes and the future NGST.
## Acknowledgments
This work was supported by the TMR Lensnet ERBFMRXCT97 - 0172 (http : // www.ast.cam. ac.uk /IoA/lensnet), the ECOS SUD Program, the French Centre National de la Recherche Scientifique, and the French Programme National de Cosmologie (PNC).
## References
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# 1 Introduction
## 1 Introduction
Five-dimensional (5D) $`𝒩`$=1 theories with global exceptional symmetries are non-trivial interacting superconformal theories . They appear in the strong-coupling limit of 5D $`𝒩`$=1 supersymmetric $`SU(2)`$ gauge theory with $`N_f`$ quark hypermultiplets. For $`N_f8`$ the microscopic global symmetry $`SO(2N_f)\times U(1)`$ is enhanced to $`E_{N_f+1}`$. When $`N_f=0`$ we have two theories with global $`E_1`$ and $`\stackrel{~}{E}_1`$ symmetry. The $`\stackrel{~}{E}_1`$ theory is shown to further flow down to the $`\stackrel{~}{E}_0`$ theory which has no global symmetry.
It is well-known that these 5D $`E_N`$ theories are obtained by compactifying M-theory on a Calabi–Yau threefold with a shrinking four-cycle realized by del Pezzo surfaces . Then the geometrical meaning of the sequence of $`E`$-type symmetries is most naturally understood in terms of the blowing-up process of $`𝑷^2`$ (or of $`𝑷^1\times 𝑷^1`$) which yields del Pezzo surfaces. The correspondence is summarized in Table 1.
Another interesting realization of 5D $`E`$-type theories is provided by the IIB 5-brane web construction including background 7-branes . The advantageous point in this is that the affine property of $`E_N`$ symmetry, which has been known to occur naturally as the action of the Weyl group of the affine algebra $`\widehat{E}_N`$ in the study of the del Pezzo surfaces , can be captured explicitly by the 7-brane configurations, thanks to the recent advances in 7-brane technology .
Furthermore it is shown that the Coulomb branch of 5D $`E_N`$ theories compactified on a circle is described in terms of a D3-brane probing the affine 7-brane backgrounds . This provides us with an intuitive physical picture behind important calculations performed in . Then, in the light of the analysis of and the 7-brane picture, we see that the $`E_N`$ theories on $`𝑹^4\times S^1`$ is considered as either the D3–7-brane system or the local system realized in singular Calabi–Yau space with a shrinking del Pezzo four-cycle; these two systems are mirror dual to each other. In this paper, thus, we revisit local mirror symmetry associated with del Pezzo surfaces in Calabi–Yau threefolds and elucidate on the relation between IIB string junctions in the affine 7-brane backgrounds and IIA D-branes wrapped on del Pezzo surfaces. A map between these two objects has been worked out recently in in which an isomorphism between the junction lattices and the even homology lattice of del Pezzo surfaces, which is identified with the Ramond-Ramond (RR) charge lattice of IIA D-branes, is shown by comparing the intersection forms. In our approach, on the other hand, we evaluate explicitly the mirror map and the BPS central charge so as to convert junction charges to the central charges of D-branes.
The paper is organized as follows: In the next section, we calculate in detail the Seiberg–Witten (SW) period integrals which describe the Coulomb branch of 5D $`E_N`$ theories on $`𝑹^4\times S^1`$. The monodromy matrices and the prepotentials are obtained. In section 3, the mirror map is constructed. In section 4, we analyze D-branes localized on a surface. Several important invariants of the BPS D-branes such as the RR charge, the central charge and intersection pairings are given in algebro-geometrical terms. In section 5, the results of the preceding section are utilized to verify the map between string junctions and D-branes on a del Pezzo four-cycle.
## 2 Calculation of the periods and monodromies
The elliptic curves for the $`\widehat{E}_N`$ 7-brane configurations are obtained in . These curves describe the Coulomb branch of 5D $`E_N`$ theories compactified on $`S^1`$ . For $`E_{N=8,7,6}`$ theories, they are given by
$`\widehat{E}_8`$ $`:y^2=x^3+R^2u^2x^22u^5,`$
$`\widehat{E}_7`$ $`:y^2=x^3+R^2u^2x^2+2u^3x,`$ (2.1)
$`\widehat{E}_6`$ $`:y^2=x^3+R^2u^2x^22Riu^3xu^4,`$
where $`u`$ is a complex moduli parameter with mass dimension 6, 4, 3 for $`\widehat{E}_8`$, $`\widehat{E}_7`$, $`\widehat{E}_6`$ and $`R`$ is the radius of $`S^1`$. These curves have the discriminant
$`\widehat{E}_8`$ $`:\mathrm{\Delta }=4u^{10}(2R^6u27),`$
$`\widehat{E}_7`$ $`:\mathrm{\Delta }=4u^9(R^4u8),`$ (2.2)
$`\widehat{E}_6`$ $`:\mathrm{\Delta }=iu^8(4R^3u+27i),`$
whose zeroes at $`u=0`$ represent coalescing $`(N+2)`$ 7-branes of $`E_{N=8,7,6}`$ configurations and a single zero at $`u0`$ represents a 7-brane which is responsible for extending $`E_N`$ to the affine system. The SW differentials $`\lambda _{\text{SW}}`$ associated with the $`\widehat{E}_N`$ curves are defined by
$$\frac{d\lambda _{\text{SW}}}{du}=\frac{dx}{y}+d(),$$
(2.3)
which are known to take the logarithmic form . For the curves (2.1) we find
$`\widehat{E}_8`$ $`:\lambda _{\text{SW}}={\displaystyle \frac{\kappa }{2\pi iR}}\mathrm{log}\left({\displaystyle \frac{yRux}{y+Rux}}\right){\displaystyle \frac{dx}{x}},`$
$`\widehat{E}_7`$ $`:\lambda _{\text{SW}}={\displaystyle \frac{\kappa }{2\pi iR}}\mathrm{log}\left({\displaystyle \frac{yRux}{y+Rux}}\right){\displaystyle \frac{dx}{x}},`$ (2.4)
$`\widehat{E}_6`$ $`:\lambda _{\text{SW}}={\displaystyle \frac{\kappa }{2\pi iR}}\mathrm{log}\left({\displaystyle \frac{yRux+iu^2}{y+Ruxiu^2}}\right){\displaystyle \frac{dx}{xiu/R}},`$
where $`\kappa `$ is a normalization constant and the factor $`1/R`$ ensures that $`\lambda _{\text{SW}}`$ has mass dimension unity. Notice that $`\lambda _{\text{SW}}`$ possesses the pole at $`x=0`$ for $`\widehat{E}_8`$, $`\widehat{E}_7`$ and $`x=iu/R`$ for $`\widehat{E}_6`$ because of the multivaluedness of the logarithm. Hence the set of period integrals $`\mathrm{\Pi }=(s,a,a_D)`$ consists of
$`s`$ $`={\displaystyle _C}\lambda _{\text{SW}}={\displaystyle \frac{2\pi i\kappa }{R}},`$ (2.5)
$`a(u)`$ $`={\displaystyle _A}\lambda _{\text{SW}}={\displaystyle ^u}𝑑u^{}\varpi (u^{}),`$ (2.6)
$`a_D(u)`$ $`={\displaystyle _B}\lambda _{\text{SW}}={\displaystyle ^u}𝑑u^{}\varpi _D(u^{}),`$ (2.7)
where the contour $`C`$ surrounds the pole of $`\lambda _{\text{SW}}`$ and $`A`$, $`B`$ are the homology cycles on $`\widehat{E}_N`$ torus along which the torus periods $`\varpi (u)`$, $`\varpi _D(u)`$ are defined as usual,
$$\varpi (u)=_A\frac{dx}{y},\varpi _D(u)=_B\frac{dx}{y}.$$
(2.8)
### 2.1 Picard–Fuchs equations
The period $`\mathrm{\Pi }`$ is evaluated with the use of the Picard–Fuchs equations. It is convenient to introduce a dimensionless variable
$$z=\frac{2}{27}R^6u,\frac{1}{8}R^4u,\frac{4i}{27}R^3u$$
(2.9)
for $`\widehat{E}_8`$, $`\widehat{E}_7`$, $`\widehat{E}_6`$, respectively, to write down the Picard–Fuchs equations. We obtain
$$_{\widehat{E}_N}\mathrm{\Pi }(z)=0,$$
(2.10)
where
$`_{\widehat{E}_8}`$ $`=\left[36z^2(z1){\displaystyle \frac{d^2}{dz^2}}+4z(27z18){\displaystyle \frac{d}{dz}}+(36z5)\right]{\displaystyle \frac{d}{dz}},`$
$`_{\widehat{E}_7}`$ $`=\left[16z^2(z1){\displaystyle \frac{d^2}{dz^2}}+16z(3z2){\displaystyle \frac{d}{dz}}+(16z3)\right]{\displaystyle \frac{d}{dz}},`$ (2.11)
$`_{\widehat{E}_6}`$ $`=\left[9z^2(z1){\displaystyle \frac{d^2}{dz^2}}+z(27z18){\displaystyle \frac{d}{dz}}+(9z2)\right]{\displaystyle \frac{d}{dz}}.`$
In order to derive the solution we first solve the second-order equation for $`d\mathrm{\Pi }/dz`$, and then integrate over $`z`$ to get $`\mathrm{\Pi }(z)`$. Substituting the form $`d\mathrm{\Pi }/dz=z^\rho F(z)`$ it is seen that $`\rho =\frac{5}{6}`$, $`\frac{3}{4}`$, $`\frac{2}{3}`$ for $`\widehat{E}_8`$, $`\widehat{E}_7`$, $`\widehat{E}_6`$, and $`F(z)`$ obeys the standard hypergeometric equation
$$\left[z(1z)\frac{d^2}{dz^2}+\left(\gamma (\alpha +\beta +1)z\right)\frac{d}{dz}\alpha \beta \right]F(z)=0,$$
(2.12)
with $`\alpha =\beta =\rho +1`$ and $`\gamma =2(\rho +1)`$, that is,
$$\alpha =\frac{1}{6},\frac{1}{4},\frac{1}{3}$$
(2.13)
for $`\widehat{E}_8,\widehat{E}_7,\widehat{E}_6`$.<sup>2</sup><sup>2</sup>2We shall use the notation $`\alpha `$ specifically to denote these numbers throughout this paper. The torus periods around the regular singular points at $`z=0,1,\mathrm{}`$ are then expressed as
$$\left(\begin{array}{c}\varpi _D^{()}(z)\\ \varpi ^{()}(z)\end{array}\right)=C_{}\left(\begin{array}{c}\phi _1^{()}(z)\\ \phi _2^{()}(z)\end{array}\right),$$
(2.14)
where $``$ stands for $`0,1,\mathrm{}`$ corresponding to $`z=0,1,\mathrm{}`$, $`C_{}`$ are $`2\times 2`$ coefficient matrices and the set of fundamental solutions has been taken as follows:
$$\{\begin{array}{cc}\phi _1^{(0)}(z)=z^{(1\alpha )}{}_{2}{}^{}F_{1}^{}(\alpha ,\alpha ;2\alpha ;z),\hfill & \\ \phi _2^{(0)}(z)=z^\alpha {}_{2}{}^{}F_{1}^{}(1\alpha ,1\alpha ;2(1\alpha );z),\hfill & \end{array}$$
(2.15)
$$\{\begin{array}{cc}\phi _1^{(1)}(z)=z^{(1\alpha )}{}_{2}{}^{}F_{1}^{}(\alpha ,\alpha ;1;1z),\hfill & \\ \phi _2^{(1)}(z)=z^{(1\alpha )}({}_{2}{}^{}F_{1}^{}(\alpha ,\alpha ,;1;1z)\mathrm{log}(1z)+{}_{2}{}^{}F_{1}^{}{}_{}{}^{}(\alpha ,\alpha ;1;1z)),\hfill & \end{array}$$
(2.16)
and
$$\{\begin{array}{cc}\phi _1^{(\mathrm{})}(z)=\frac{1}{z}{}_{2}{}^{}F_{1}^{}(\alpha ,1\alpha ;1;\frac{1}{z}),\hfill & \\ \phi _2^{(\mathrm{})}(z)=\frac{1}{z}\left({}_{2}{}^{}F_{1}^{}(\alpha ,1\alpha ;1;\frac{1}{z})\mathrm{log}(z)+{}_{2}{}^{}F_{1}^{}{}_{}{}^{}(\alpha ,1\alpha ;1;\frac{1}{z})\right).\hfill & \end{array}$$
(2.17)
Here $`{}_{2}{}^{}F_{1}^{}(\alpha ,\beta ;\gamma ;z)`$ is the hypergeometric function
$${}_{2}{}^{}F_{1}^{}(\alpha ,\beta ;\gamma ;z)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(\alpha )_n(\beta )_n}{(\gamma )_n}\frac{z^n}{n!}$$
(2.18)
and
$${}_{2}{}^{}F_{1}^{}{}_{}{}^{}(\alpha ,\beta ;\gamma ;z)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{(\alpha )_n(\beta )_n}{(\gamma )_n}\left[\underset{k=0}{\overset{n1}{}}\left(\frac{1}{\alpha +k}+\frac{1}{\beta +k}\frac{2}{\gamma +k}\right)\right]\frac{z^n}{n!},$$
(2.19)
where $`(\alpha )_n=\mathrm{\Gamma }(\alpha +n)/\mathrm{\Gamma }(\alpha )`$ with $`\mathrm{\Gamma }(x)`$ being the gamma function.
The coefficient matrix $`C_1`$ is determined by directly evaluating the elliptic integrals. The result reads
$$C_1=\stackrel{~}{c}R^{\frac{1\alpha }{\alpha }}\left(\begin{array}{cc}\pi & 0\\ \frac{i}{2}\left(\mathrm{log}vi\pi \right)& \frac{i}{2}\end{array}\right)$$
(2.20)
where $`\stackrel{~}{c}=2^2/3^3`$, $`2^{3/2}`$, $`2^{5/2}/3^2`$<sup>3</sup><sup>3</sup>3In what follows we will set $`\stackrel{~}{c}=1`$ for simplicity. and
$$v=432,64,27$$
(2.21)
for $`\widehat{E}_8,\widehat{E}_7,\widehat{E}_6`$. Then, performing the analytic continuation we calculate $`C_0`$ and $`C_{\mathrm{}}`$. The connection matrices $`X,Y`$ defined by
$$\left(\begin{array}{c}\phi _1^{(0)}\\ \phi _2^{(0)}\end{array}\right)=X\left(\begin{array}{c}\phi _1^{(1)}\\ \phi _2^{(1)}\end{array}\right),\left(\begin{array}{c}\phi _1^{(0)}\\ \phi _2^{(0)}\end{array}\right)=Y\left(\begin{array}{c}\phi _1^{(\mathrm{})}\\ \phi _2^{(\mathrm{})}\end{array}\right)$$
(2.22)
are found with the aid of the Barnes’ integral representation of the hypergeometric function \[13, p. 286\]. It turns out that
$$X=\left(\begin{array}{cc}\xi _1\eta _1& \xi _1\\ \xi _2\eta _2& \xi _2\end{array}\right),Y=\left(\begin{array}{cc}\xi _1\eta _1e^{i\pi \alpha }& \xi _1e^{i\pi \alpha }\\ \xi _2\eta _2e^{i\pi (1\alpha )}& \xi _2e^{i\pi (1\alpha )}\end{array}\right),$$
(2.23)
where
$`\xi _1={\displaystyle \frac{\mathrm{\Gamma }(2\alpha )}{\mathrm{\Gamma }^2(\alpha )}},`$ $`\xi _2={\displaystyle \frac{\mathrm{\Gamma }(22\alpha )}{\mathrm{\Gamma }^2(1\alpha )}},`$ (2.24)
$`\eta _1=2(\psi (1)\psi (\alpha )),`$ $`\eta _2=2(\psi (1)\psi (1\alpha )),`$ (2.25)
and $`\psi (x)=\frac{d}{dx}\mathrm{log}\mathrm{\Gamma }(x)`$ is the digamma function. Thus we obtain
$`C_0`$ $`=C_1X^1={\displaystyle \frac{\pi R^{\frac{1\alpha }{\alpha }}}{2(2\alpha 1)}}\left(\begin{array}{cc}2\xi _2& 2\xi _1\\ \frac{\omega }{\mathrm{sin}\pi \alpha }\xi _2& \frac{\overline{\omega }}{\mathrm{sin}\pi \alpha }\xi _1\end{array}\right),\omega =e^{i\pi \left(\frac{1}{2}\alpha \right)},`$ (2.26)
$`C_{\mathrm{}}`$ $`=C_1X^1Y=R^{\frac{1\alpha }{\alpha }}\left(\begin{array}{cc}(\mathrm{log}v+i\pi )\mathrm{sin}\pi \alpha & \mathrm{sin}\pi \alpha \\ \frac{i\pi }{2\mathrm{sin}\pi \alpha }& 0\end{array}\right).`$ (2.27)
When the periods go around each regular singular point counterclockwise they undergo the monodromy. If we denote as $`T_{}`$ the monodromy matrix of $`\left(\begin{array}{c}\phi _1^{()}\\ \phi _2^{()}\end{array}\right)`$, the monodromy matrix $`\widehat{M}_{}`$ acting on $`\left(\begin{array}{c}\varpi _D^{()}\\ \varpi ^{()}\end{array}\right)`$ is given by $`\widehat{M}_{}=C_{}T_{}C_{}^1`$. For $`\widehat{E}_{N=8,7,6}`$ one computes
$$\widehat{M}_0=\left(\begin{array}{cc}N8& N9\\ 1& 1\end{array}\right),\widehat{M}_1=\left(\begin{array}{cc}1& 0\\ 1& 1\end{array}\right),\widehat{M}_{\mathrm{}}=\left(\begin{array}{cc}1& N9\\ 0& 1\end{array}\right)$$
(2.28)
which obey $`\widehat{M}_{\mathrm{}}=\widehat{M}_0\widehat{M}_1`$, $`\widehat{M}_0^6=1`$ ($`\widehat{E}_8`$), $`\widehat{M}_0^4=1`$ ($`\widehat{E}_7`$), $`\widehat{M}_0^3=1`$ ($`\widehat{E}_6`$). If we follow the convention in , $`K_{}=\widehat{M}_{}^1`$ yields the monodromy around the 7-brane configurations. For the $`[p,q]`$ 7-brane $`𝑿_{[p,q]}`$, the monodromy matrix reads
$$K_{[p,q]}=\left(\begin{array}{cc}1+pq& p^2\\ q^2& 1pq\end{array}\right),$$
(2.29)
whereas for the $`\widehat{E}_N`$ 7-branes we have
$$K(\widehat{𝑬}_N)=\left(\begin{array}{cc}1& 9N\\ 0& 1\end{array}\right).$$
(2.30)
In (2.28) we observe $`\widehat{M}_{\mathrm{}}^1=K(\widehat{𝑬}_N)`$, $`\widehat{M}_1^1=K_{[0,1]}`$ and $`\widehat{M}_0^1=K_{[3,2]}K_{[3,1]}K_{[1,0]}^N`$. Accordingly our 7-brane configuration is identified as
$$\widehat{𝑬}_N=𝑨^N𝑿_{[3,1]}𝑿_{[3,2]}𝑿_{[0,1]},$$
(2.31)
where we have used the notation in and $`𝑨=𝑿_{[1,0]}`$ is the D7-brane. This is shown to be equivalent to the canonical one $`\widehat{𝑬}_N=𝑨^{N1}𝑩𝑪^2𝑿_{[3,1]}`$ by making use of the brane move .
### 2.2 Seiberg–Witten periods
Now our task is to calculate the Seiberg–Witten periods $`a(z)`$, $`a_D(z)`$ from the torus periods through (2.5)–(2.7). The important subtlety arises in evaluating the integration constants. First of all, since $`\lambda _{\text{SW}}`$ vanishes at $`u=0`$, $`a(z)`$ and $`a_D(z)`$ must vanish as well at $`z=0`$ . In the patch $`|z|<1`$, the SW periods are thus given by<sup>4</sup><sup>4</sup>4In the following we will ignore the irrelevant overall constants and put $`R=1`$ for simplicity. There is no difficulty in recovering them.
$`a_D^{(0)}(z)`$ $`={\displaystyle _0^z}𝑑z^{}\varpi _D^{(0)}(z^{}),`$ (2.32)
$`a^{(0)}(z)`$ $`={\displaystyle _0^z}𝑑z^{}\varpi ^{(0)}(z^{}).`$ (2.33)
We note that these are succinctly expressed in terms of the generalized hypergeometric function
$${}_{3}{}^{}F_{2}^{}(\alpha _1,\alpha _2,\alpha _3;\beta _1,\beta _2;z)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(\alpha _1)_n(\alpha _2)_n(\alpha _3)_n}{(\beta _1)_n(\beta _2)_n}\frac{z^n}{n!}$$
(2.34)
in such a way that
$$\left(\begin{array}{c}a_D^{(0)}(z)\\ a^{(0)}(z)\end{array}\right)=C_0\left(\begin{array}{c}V_1^{(0)}(z)\\ V_2^{(0)}(z)\end{array}\right),$$
(2.35)
where
$`V_1^{(0)}(z)`$ $`={\displaystyle \frac{1}{\alpha }}z^\alpha {}_{3}{}^{}F_{2}^{}(\alpha ,\alpha ,\alpha ;2\alpha ,1+\alpha ;z),`$ (2.36)
$`V_2^{(0)}(z)`$ $`={\displaystyle \frac{1}{1\alpha }}z^{1\alpha }{}_{3}{}^{}F_{2}^{}(1\alpha ,1\alpha ,1\alpha ;2(1\alpha ),2\alpha ;z).`$ (2.37)
Next the SW periods in the patch $`|1z|1`$ are
$`a_D^{(1)}(z)`$ $`={\displaystyle _1^z}𝑑z^{}\varpi _D^{(1)}(z^{})+c_D^{(1)},`$ (2.38)
$`a^{(1)}(z)`$ $`={\displaystyle _1^z}𝑑z^{}\varpi ^{(1)}(z^{})+c^{(1)},`$ (2.39)
where $`c_D^{(1)}`$ and $`c^{(1)}`$ are integration constants. Notice that the analytic continuation allows us to write
$$a_D^{(1)}(z)=_0^z𝑑z^{}\varpi _D^{(1)}(z^{})$$
(2.40)
which obeys $`a_D^{(1)}(0)=a_D^{(0)}(0)=0`$. Therefore,
$`c_D^{(1)}`$ $`=a_D^{(1)}(1)`$
$`=\pi {\displaystyle _0^1}𝑑xx^{(1\alpha )}{}_{2}{}^{}F_{1}^{}(\alpha ,\alpha ;1;1x)`$
$`=\pi {\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\alpha )_n^2}{n!}}{\displaystyle \frac{\mathrm{\Gamma }(\alpha )}{\mathrm{\Gamma }(\alpha +n+1)}}`$
$`=\pi {\displaystyle \frac{\mathrm{\Gamma }(\alpha )}{\mathrm{\Gamma }(\alpha +1)}}{}_{2}{}^{}F_{1}^{}(\alpha ,\alpha ;1+\alpha ;1).`$ (2.41)
Here the well-known formula
$${}_{2}{}^{}F_{1}^{}(\alpha ,\beta ;\gamma ;1)=\frac{\mathrm{\Gamma }(\gamma )\mathrm{\Gamma }(\gamma \alpha \beta )}{\mathrm{\Gamma }(\gamma \alpha )\mathrm{\Gamma }(\gamma \beta )},\text{Re}\gamma >0,\text{Re}(\gamma \alpha \beta )>0$$
(2.42)
helps us to find
$$c_D^{(1)}=\pi \mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(1\alpha )=\frac{\pi ^2}{\mathrm{sin}\pi \alpha }.$$
(2.43)
Since $`c_D^{(1)}`$ is also obtained from $`c_D^{(1)}=a_D^{(0)}(1)`$, (2.43) yields the nontrivial identity for the generalized hypergeometric function $`{}_{3}{}^{}F_{2}^{}`$ at $`z=1`$ via (2.35). For $`\alpha =\frac{1}{3}`$, this identity was first found numerically in , for which the above calculation affords analytic proof.
On the other hand, the similar integral for $`a^{(1)}(1)`$ is not helpful to determine $`c^{(1)}`$ analytically. Hence, evaluating $`c^{(1)}=a^{(0)}(1)`$ numerically we determine
$$c^{(1)}=\frac{\pi ^2}{\mathrm{sin}\pi \alpha }\times \{\begin{array}{cc}0.5000+0.9281i\text{for}\widehat{E}_8\hfill & \\ 0.5000+0.6103i\text{for}\widehat{E}_7\hfill & \\ 0.5000+0.4628i\text{for}\widehat{E}_6.\hfill & \end{array}$$
(2.44)
These constants were first evaluated in .<sup>2</sup><sup>2</sup>2We should note that our choice of the branch is $`z=e^{i\pi }z`$. If the other branch $`z=e^{i\pi }z`$ was chosen one would have $`0.5000`$ instead of $`0.5000`$ in (2.44) in agreement with .
Let us now turn to the analysis of the SW periods around $`z=\mathrm{}`$. In order to fix the integration constants, we adopt the idea in and employ the Barnes’ integral representation of the hypergeometric function
$${}_{2}{}^{}F_{1}^{}(\alpha ,\beta ;\gamma ;z)=\frac{\mathrm{\Gamma }(\gamma )}{\mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(\beta )}_i\mathrm{}^i\mathrm{}\frac{dt}{2\pi i}\frac{\mathrm{\Gamma }(t+\alpha )\mathrm{\Gamma }(t+\beta )\mathrm{\Gamma }(t)}{\mathrm{\Gamma }(t+\gamma )}(z)^t,$$
(2.45)
where $`|\mathrm{arg}(z)|<\pi `$. When the integration contour is closed on the right we have the power series as presented in (2.15) which converges for $`|z|<1`$. For our purpose, we first check how (2.45) can be used to reproduce the SW periods in the patch $`|z|<1`$. The naive idea is that expressing the torus periods in the form (2.45), we first make the $`z`$-integral and then perform the contour integral with respect to $`t`$. We are thus led to consider the integral
$$I_\alpha (z)=\frac{\mathrm{\Gamma }(2\alpha )}{\mathrm{\Gamma }^2(\alpha )}_i\mathrm{}^i\mathrm{}\frac{dt}{2\pi i}\frac{\mathrm{\Gamma }^2(t+\alpha )\mathrm{\Gamma }(t)}{\mathrm{\Gamma }(t+1\alpha )(t+\alpha )}(z)^tz^\alpha .$$
(2.46)
Closing the contour on the right we immediately obtain
$$I_\alpha (z)_{\text{right}}=\frac{1}{\alpha }z^\alpha {}_{3}{}^{}F_{2}^{}(\alpha ,\alpha ,\alpha ;2\alpha ,1+\alpha ;z),$$
(2.47)
and hence (2.35) is reproduced.
If we close the contour on the left, we obtain the expression which is valid for $`|z|>1`$. Thus
$$\left(\begin{array}{c}a_D^{(\mathrm{})}(z)\\ a^{(\mathrm{})}(z)\end{array}\right)=C_0\left(\begin{array}{c}I_\alpha (z)_{\text{left}}\\ I_{1\alpha }(z)_{\text{left}}\end{array}\right).$$
(2.48)
Upon doing the contour integral on the left, one has to pick up the triple pole at $`t=\alpha `$ in addition to the standard poles in the Barnes’ integral. This in fact gives rise to the integration constant. After some algebra we arrive at
$$\left(\begin{array}{c}a_D^{(\mathrm{})}(z)\\ a^{(\mathrm{})}(z)\end{array}\right)=C_{\mathrm{}}\left(\begin{array}{c}V_1^{(\mathrm{})}(z)\\ V_2^{(\mathrm{})}(z)\end{array}\right)+C_0\left(\begin{array}{c}b_1\\ b_2\end{array}\right),$$
(2.49)
where
$`V_1^{(\mathrm{})}(z)`$ $`=\mathrm{log}(z)g(z),`$ (2.50)
$`V_2^{(\mathrm{})}(z)`$ $`={\displaystyle \frac{1}{2}}\mathrm{log}^2(z)\mathrm{log}(z)g(z)+h(z),`$ (2.51)
with
$`g(z)`$ $`={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\alpha )_n(1\alpha )_n}{(n!)^2n}}{\displaystyle \frac{1}{z^n}},`$ (2.52)
$`h(z)`$ $`={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\alpha )_n(1\alpha )_n}{(n!)^2n}}\left[{\displaystyle \underset{k=0}{\overset{n1}{}}}\left({\displaystyle \frac{1}{\alpha +k}}+{\displaystyle \frac{1}{1\alpha +k}}{\displaystyle \frac{2}{1+k}}\right){\displaystyle \frac{1}{n}}\right]{\displaystyle \frac{1}{z^n}}.`$ (2.53)
In (2.49) the integration constants read
$$b_1=\frac{1}{2}e^{i\pi \alpha }\xi _1\left(\eta _1^2+\frac{\pi ^2}{3}\right),b_2=\frac{1}{2}e^{i\pi (1\alpha )}\xi _2\left(\eta _2^2+\frac{\pi ^2}{3}\right).$$
(2.54)
Further manipulations yield
$$a_D^{(\mathrm{})}(z)=4\pi ^2\mathrm{sin}\pi \alpha \frac{H_2(z)}{(2\pi i)^2},a^{(\mathrm{})}(z)=\frac{\pi ^2}{\mathrm{sin}\pi \alpha }\frac{H_1(z)}{2\pi i},$$
(2.55)
where
$`H_1(z)`$ $`=\mathrm{log}(vz)+g(z),`$ (2.56)
$`H_2(z)`$ $`={\displaystyle \frac{1}{2}}\mathrm{log}^2(vz)+\mathrm{log}(vz)g(z)h(z){\displaystyle \frac{\pi ^2}{2}}\left({\displaystyle \frac{1}{\mathrm{sin}^2\pi \alpha }}+{\displaystyle \frac{1}{3}}\right).`$ (2.57)
Now that we have fixed all the integration constants, it is seen that the monodromy matrices with integral entries are obtained by setting the constant period
$$s=\frac{\pi ^2}{\mathrm{sin}\pi \alpha }.$$
(2.58)
As a result, the monodromy matrices acting on the period vector $`{}_{}{}^{t}\mathrm{\Pi }={}_{}{}^{t}(s,a(z),a_D(z))`$ for $`\widehat{E}_{N=8,7,6}`$ theories turn out to be
$`M_0=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 1\\ 0& N9& N8\end{array}\right),M_1=\left(\begin{array}{ccc}1& 0& 0\\ 1& 1& 1\\ 0& 0& 1\end{array}\right),M_{\mathrm{}}=\left(\begin{array}{ccc}1& 0& 0\\ 1& 1& 0\\ 9N& N9& 1\end{array}\right)`$
which obey $`M_{\mathrm{}}=M_0M_1`$, $`M_0^6=1`$ ($`\widehat{E}_8`$), $`M_0^4=1`$ ($`\widehat{E}_7`$) and $`M_0^3=1`$ ($`\widehat{E}_6`$).
Finally the BPS central charge is expressed in terms of the period integrals
$$Z=pa(z)qa_D(z)+ns,$$
(2.60)
where $`(p,q)𝒁^2`$ are electric and magnetic charges and $`n𝒁`$ is global conserved charge of the BPS states. In view of the D3-probe picture, (2.60) is the central charge corresponding to the string junction,
$$𝑱=\underset{i=1}{\overset{N}{}}\lambda _i𝝎^i+P𝝎^p+Q𝝎^qn𝒛_{[0,1]},$$
(2.61)
where $`P=p`$, $`Q=q+n`$ and our notation is as follows; $`\{\lambda _i\}`$ is the Dynkin label of $`E_N`$ weights and the $`𝝎^i`$ are the junctions of zero asymptotic charges representing the fundamental weights of the Lie algebra $`E_N`$, $`𝝎^p`$, $`𝝎^q`$ are the $`E_N`$ singlet junctions with asymptotic charges $`(1,0)`$, $`(0,1)`$, respectively, and $`𝒛_{[0,1]}`$ stands for the outgoing (0,1)-string emanating from the 7-brane $`𝑿_{[0,1]}`$ at $`z=1`$. See Fig. 1a for the $`\widehat{E}_N`$ brane-junction configurations probed by a single D3-brane. To justify the above correspondence, for instance, take a junction with charges $`(p,q;n)=(0,1;1)`$. This is a single BPS string stretched between the D3-brane and the 7-brane $`𝑿_{[0,1]}`$. Clearly its mass vanishes when the D3-brane is located at the point $`z=1`$. This is indeed verified from (2.60) since $`Z(z)=a_D(z)+s`$ for $`(0,1;1)`$ charges and $`Z(1)=a_D(1)+s=0`$ at $`z=1`$ by virtue of (2.43).
We note that there are no terms in (2.60) which reflect the presence of the $`E_N`$ non-singlet junctions in (2.61) explicitly. To have such terms one has to turn on mass deformation parameters in the $`\widehat{E}_N`$ curves so that the SW differentials have poles with non-vanishing residues other than the one at $`x=0`$ ($`x=iu/R`$ for $`\widehat{E}_6`$). This will produce additional terms in (2.60) which may depend directly on the $`E_N`$ representations.
Notice that the string junction (2.61) is equivalently expressed as
$$𝑱=\underset{i=1}{\overset{N}{}}\lambda _i𝝎^i+p𝝎^p+q𝝎^q+n𝜹^{(1,0)},$$
(2.62)
where $`𝜹^{(1,0)}`$ is a loop junction which represents the imaginary root of the affine algebra $`\widehat{E}_N`$, see Fig. 1b. In the junction realization of $`\widehat{E}_N`$, the level $`k`$ of representations is given by $`k(𝑱)=(𝑱,𝜹^{(1,0)})=q`$, while the grade $`\overline{n}(\lambda )`$ for a weight vector $`\{\lambda _i\}`$ is equal to $`n`$ in $`𝑱`$ up to a constant shift .
On the D3-brane probing the region $`|z|<1`$, in the limit $`R0`$, the theory reduces to the 4D superconformal theory with global exceptional symmetries . Recovering the $`R`$-dependence we have the central charge in the form
$$Z=pa(u)qa_D(u)+\frac{n}{R}s,$$
(2.63)
where $`a(u)`$, $`a_D(u)\text{const}u^\alpha `$ and the last term represents the KK modes associated with the $`S^1`$ compactification. We thus observe that the grading of KK modes is regarded as the the grade in the affine symmetry $`\widehat{E}_N`$. As will be seen in section 5.3 these KK modes are identified with the D0-branes in M-theory. The decoupling of the KK modes as $`R0`$ is ensured since the 7-brane $`𝑿_{[0,1]}`$ moves away to infinity on the $`u`$-plane, and hence the loop junction $`𝜹^{(1,0)}`$ decouples, leaving the finite symmetry $`E_N`$ for 4D theories.
## 3 Mirror map
5D $`𝒩`$=1 supersymmetric $`SU(2)`$ gauge theories with exceptional global symmetries are realized by compactifying M-theory on a Calabi–Yau threefold with a vanishing del Pezzo four-cycle . In the local mirror geometry of the singularity associated with a shrinking del Pezzo four-cycle of the type $`E_{N=8,7,6}`$ is modeled by the Landau-Ginzburg potential
$`W_{E_8}`$ $`={\displaystyle \frac{1}{x_0^6}}+x_1^2+x_2^3+x_3^6+x_4^6\psi x_0x_1x_2x_3x_4,`$
$`W_{E_7}`$ $`={\displaystyle \frac{1}{x_0^4}}+x_1^2+x_2^4+x_3^4+x_4^4\psi x_0x_1x_2x_3x_4,`$ (3.1)
$`W_{E_6}`$ $`={\displaystyle \frac{1}{x_0^3}}+x_1^3+x_2^3+x_3^3+x_4^3\psi x_0x_1x_2x_3x_4.`$
The equations $`W_{E_N}=0`$ describe non-compact Calabi–Yau manifolds and $`\psi `$ is a complex moduli parameter. The point $`\psi =0`$ is the $`E_N`$ symmetric point (the Landau-Ginzburg point) at which the del Pezzo four-cycle collapses. The large complex structure limit is taken by letting $`\psi =\mathrm{}`$.
It turns out that the natural complex modulus is $`\psi ^{\mathrm{}}`$ with $`\mathrm{}=6`$ ($`E_8`$), 4 ($`E_7`$) and 3 ($`E_6`$) and, as will be seen momentarily, $`\psi ^{\mathrm{}}`$ and $`z`$ in the previous section is related through
$$z=\frac{\psi ^6}{432}(E_8),\frac{\psi ^4}{64}(E_7),\frac{\psi ^3}{27}(E_6).$$
(3.2)
The period integral $`\stackrel{~}{\mathrm{\Pi }}`$ of the holomorphic three-form $`\mathrm{\Omega }`$ associated with (3.1) are defined over a basis of three-cycles on the mirror Calabi–Yau . They obey the Picard–Fuchs equations
$$_{\text{ell}}^{E_N}\vartheta \stackrel{~}{\mathrm{\Pi }}=0,$$
(3.3)
where $`\vartheta =c\frac{}{c}`$ with $`c=\psi ^{\mathrm{}}`$ and $`_{\text{ell}}^{E_N}`$ are the Picard–Fuchs operators corresponding to the elliptic singularities of type $`\widehat{E}_N`$
$`_{\text{ell}}^{E_8}`$ $`=\vartheta ^212c(6\vartheta +5)(6\vartheta +1),`$
$`_{\text{ell}}^{E_7}`$ $`=\vartheta ^24c(4\vartheta +3)(4\vartheta +1),`$ (3.4)
$`_{\text{ell}}^{E_6}`$ $`=\vartheta ^23c(3\vartheta +2)(3\vartheta +1).`$
Making a change of variables with (3.2), one can check that (3.3) is equivalent to (2.11). In fact, it is shown in that the Calabi–Yau periods $`\stackrel{~}{\mathrm{\Pi }}(c)`$ over three-cycles reduce to the SW periods $`\mathrm{\Pi }`$ after performing the integrals over appropriate two-cycles. Thus we have $`\stackrel{~}{\mathrm{\Pi }}=\mathrm{\Pi }=(s,a(z),a_D(z))`$.
Under mirror symmetry three-cycles on the IIB side is mapped to the zero, two, four-cycles on the IIA side where the four-cycle is the del Pezzo surface. Following the standard machinery we find that the complex Kähler modulus $`t`$ is given by
$$t(z)=\frac{H_1(z)}{2\pi i}$$
(3.5)
and its dual $`t_d`$ becomes
$$t_d(z)=\frac{}{t}=\frac{H_2(z)}{(2\pi i)^2}.$$
(3.6)
Thus the SW periods and the flat coordinate system are related via
$$\left(\begin{array}{c}s\\ a(z)\\ a_D(z)\end{array}\right)=\frac{\pi ^2}{\mathrm{sin}\pi \alpha }\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 4\mathrm{sin}^2\pi \alpha \end{array}\right)\left(\begin{array}{c}1\\ t(z)\\ t_d(z)\end{array}\right).$$
(3.7)
The central charge (2.60) is then rewritten as
$$Z=\frac{\pi ^2}{\mathrm{sin}\pi \alpha }\left(4\mathrm{sin}^2\pi \alpha qt_d(z)+pt(z)+n\right).$$
(3.8)
In the large radius region $`|z|>1`$, it is shown from (2.55) that
$$t_d=\frac{t^2}{2}\frac{t}{2}\frac{1}{24}\left(\frac{3}{\mathrm{sin}^2\pi \alpha }2\right)+O(e^{2\pi it}),$$
(3.9)
from which we see that $`t_d`$ is the central charge of the D4-brane. Integrating this over $`t`$ we obtain the prepotential
$$=\frac{t^3}{6}\frac{t^2}{4}\frac{1}{24}\left(\frac{3}{\mathrm{sin}^2\pi \alpha }2\right)t+\text{const.}\frac{1}{(2\pi i)^3}\underset{k=1}{\overset{\mathrm{}}{}}n_kLi_3(e^{2\pi ikt}),$$
(3.10)
where $`Li_3(x)`$ is the trilogarithm and the instanton coefficients $`n_k`$ are presented in Table 2. We note that the $`n_k`$’s multiplied by $`4\mathrm{sin}^2\pi \alpha `$ for $`\widehat{E}_N`$ coincide with the Gromov–Witten invariants obtained in .
## 4 D-branes on a surface
The aim of this section is to obtain several basic invariants of BPS D-branes bounded on a surface, which are expressed as coherent sheaves on the surface, such as the RR charge, the central charge, and some intersection pairings with the aid of some algebraic geometry techniques. They play a crucial role in the duality map between D-branes on del Pezzo surfaces and the string junctions with $`\widehat{E}_N`$ symmetry, which we will discuss in the next section. An important note is that though we identify a BPS D-brane on a surface with a coherent sheaf, we do not enter into the conditions, such as the stability, to be satisfied by the sheaf to represent a true BPS brane, with few exceptions; in effect, what we really need is not a sheaf itself but its invariants such as the Chern character in this article.
### 4.1 D-branes on a Calabi–Yau threefold
We represent a BPS D-brane on a Calabi–Yau threefold $`X`$ by a coherent $`𝒪_X`$-module $`𝒢`$. The RR charge of $`𝒢`$ is given by the Mukai vector
$$v_X\left(𝒢\right)=\text{ch}\left(𝒢\right)\sqrt{\text{Todd}\left(T_X\right)}H_2(X;𝑸):=\underset{i=0}{\overset{3}{}}H_{2i}(X;𝑸),$$
(4.1)
where $`\text{ch}(𝒢)=_{i=0}^3\text{ch}_i(𝒢)`$ is the Chern character with $`\text{ch}_i(𝒢)H_{62i}(X;𝑸)`$, which can be computed as follows: there always exits a resolution of $`𝒢`$ by locally free sheaves $`(V_{})`$, that is, sheaves of sections of holomorphic vector bundles: $`0V_3V_2V_1V_0𝒢0`$ (exact), thus we can set $`\text{ch}(𝒢):=_{i=0}^3(1)^i\text{ch}(V_i)`$, which does not depend on the choice of the resolution, and $`\sqrt{\text{Todd}(T_X)}=[X]+c_2(X)/24`$, the effect of which on the RR charges has been called a geometric version of the Witten effect in .
The intersection form on D-branes on $`X`$, which is of great importance in investigating topological aspects of D-branes , is given by
$$I_X(𝒢_1,𝒢_2)=v_X(𝒢_1)^{}v_X(𝒢_2)_X=\text{ch}\left(𝒢_1\right)^{}\text{ch}\left(𝒢_2\right)\text{Todd}\left(T_X\right)_X,$$
(4.2)
where $`\mathrm{}_X`$ evaluates the degree of $`H_0(X;𝑸)𝑸`$ component, and $`vv^{}`$ flips the sign of $`H_0(X)H_4(X)`$ part. In particular, if $`𝒢`$ itself is locally free, then $`\text{ch}(𝒢)^{}=\text{ch}(𝒢^{})`$, where $`𝒢^{}=om_X(𝒢,𝒪_X)`$ is the dual sheaf.
It is easy to check $`I_X(𝒢_2,𝒢_1)=I_X(𝒢_1,𝒢_2)`$. On the other hand, the Hirzebruch–Riemann–Roch formula tells us
$$I_X(𝒢_1,𝒢_2)=\underset{i=0}{\overset{3}{}}(1)^i\text{dim}\text{Ext}_X^i(𝒢_1,𝒢_2),$$
(4.3)
according to which the skew-symmetric property of the intersection form $`I_X`$ may be attributed to the Serre duality: $`\text{Ext}_X^i(𝒢_2,𝒢_1)\text{Ext}_X^{3i}(𝒢_1,𝒢_2)^{}`$ . Incidentally, the H.R.R. formula (4.3) also assures that $`I_X`$ takes values in $`𝒁`$.
Let $`J_XH_4(X;𝑹)`$ be a Kähler form on $`X`$, which is identified with an $`𝑹`$-extended ample divisor here. The classical expression of the central charge of the D-brane $`𝒢`$ is then given by
$$Z_{J_X}^{\text{cl}}(𝒢)=e^{J_X}v_X(𝒢)_X=\underset{k=0}{\overset{3}{}}\frac{(1)^k}{k!}J_X^kv_{X,k}(𝒢)_X,$$
(4.4)
where $`v_{X,k}`$ is the $`H_{2k}(X)`$ component of $`v_XH_2(X)`$. On the other hand, the quantum central charge $`Z_{J_X}(𝒢)`$ differs from its classical counterpart (4.4) by the terms of order $`O(e^{2\pi J_X\beta _X})`$ where two-cycle $`\beta `$ is in the Mori cone of $`X`$, which is dual to the Kähler cone; the exact Kähler moduli dependence of $`Z_{J_X}`$ can be determined in principle by the Picard–Fuchs equations for the periods of the mirror Calabi–Yau threefold $`X^{}`$ .
### 4.2 D-branes localized on a surface
Let $`f:SX`$ be an embedding of a projective surface $`S`$ in a Calabi–Yau threefold $`X`$. If $`c_1(S)H_2(S)`$ is nef, which means that its intersection with any effective curve $`C`$ on $`S`$ are non-negative: $`c_1(S)C0`$, there is a smooth elliptic Calabi–Yau threefold over $`S`$ with $`S`$ its cross section, which we can take as a model of embedding . Other examples of embedding can be found in .
Now let us take the limit of infinite elliptic fiber, so that the D-branes the central charge of which remains finite are those which are confined to the surface $`S`$, where we should note that some D-branes on $`S`$, a D0-brane for example, can move along elliptic fibers so as to leave $`S`$ even if $`S`$ itself is rigid in $`X`$. The properties of the D-branes localized on $`S`$ then depend not on the details of the global model $`X`$, but only on the intrinsic geometry of $`S`$ and its normal bundle $`N_S=N_{S|X}`$, which is isomorphic to the canonical line bundle $`K_S`$. In particular, this means that we can compute the central charges of BPS D-branes using local mirror symmetry principle on $`S`$ .
A D-brane sticking to $`S`$ can be described by a coherent $`𝒪_S`$-module $``$. The Euler number of it defined by $`\chi ()=_{i=0}^2(1)^ih^i(S,)`$, where $`h^i(S,)=\text{dim}H^i(S,)`$, is an important invariant, which can be obtained as follows: first we need the Todd class of $`S`$
$`\text{Todd}(T_S)`$ $`=[S]+{\displaystyle \frac{1}{2}}c_1(S)+\chi (𝒪_S)[\text{pt}],`$ (4.5)
$`\chi (𝒪_S)`$ $`={\displaystyle \frac{1}{12}}c_2(S)+c_1(S)^2_S,`$ (4.6)
second, by the H.R.R. formula, we have
$$\chi ()=\text{ch}()\text{Todd}(T_S)_S=r\chi (𝒪_S)+\text{ch}_2()+\frac{1}{2}c_1(S)c_1()_S.$$
(4.7)
There is a canonical push-forward homomorphism $`f_{}:H_2(S;𝑸)H_2(X;𝑸)`$, which maps a cycle on $`S`$ to that on $`X`$. Similarly, we can define the coherent sheaf $`f_!`$ on $`X`$ by extending $``$ by zero to $`X\backslash S`$.<sup>3</sup><sup>3</sup>3The symbol $`f_!`$ is originally defined to be $`_i(1)^iR^if_{}`$, an element of the K group of coherent $`𝒪_X`$-modules ; it reduces to the single direct image sheaf $`f_{}`$ on $`X`$, because all the higher direct images $`R^if_{}`$ vanish for embedding $`f`$ \[35, p. 102\]. The celebrated Grothendieck–Riemann–Roch formula for embedding $`f:SX`$ relates the Chern characters of $``$ and $`f_!`$ as follows:
$$\text{ch}(f_!)=f_{}\left(\text{ch}()\frac{1}{\text{Todd}(N_S)}\right).$$
(4.8)
Multiplying the both hand side of (4.8) by the square root of $`\text{Todd}(T_X)`$, we have
$$\text{ch}(f_!)\sqrt{\text{Todd}(T_X)}=f_{}\left(\text{ch}()\sqrt{\frac{\text{Todd}(T_S)}{\text{Todd}(N_S)}}\right),$$
(4.9)
where we have used the projection formula \[36, p. 273\], \[37, p. 426\]:
$$f_{}\left(\alpha f^{}\beta \right)=f_{}\alpha \beta ,\alpha H_2(S;𝑸),\beta H_2(X;𝑸),$$
(4.10)
and $`f^{}\text{Todd}(T_X)=\text{Todd}(T_S)\text{Todd}(N_S)`$, which follows from the short exact sequence of bundles on $`S`$: $`0T_Sf^{}T_XN_S0,`$ combined with the multiplicative property of the Todd class. As the left hand side of (4.9) is the D-brane charge of $``$ regarded as a brane on $`X`$, we arrive at the intrinsic description of the RR charge on $`S`$:
$$v_S()=\text{ch}()\sqrt{\frac{\text{Todd}(T_S)}{\text{Todd}(K_S)}}=\text{ch}()e^{\frac{1}{2}c_1(S)}\sqrt{\frac{\widehat{A}(T_S)}{\widehat{A}(K_S)}}H_2(S;𝑸),$$
(4.11)
which is a complex-analytic (or algebraic) derivation of the RR charge which has originally been obtained in $`C^{(\mathrm{})}`$ category . The gravitational correction factor for $`S`$ admits the following expansion:
$$\sqrt{\frac{\text{Todd}(T_S)}{\text{Todd}(K_S)}}=[S]+\frac{1}{2}c_1(S)+\frac{1}{24}\left(c_2(S)+3c_1(S)^2\right)H_2(S;𝑸).$$
(4.12)
As a simple exercise let us compute the RR charge of a sheaf on $`S`$. To this end, let $`\iota :CS`$ be an embedding of a smooth genus $`g`$ curve in $`S`$ with normal bundle $`N_C=N_{C|S}`$. Then from a line bundle $`L_C`$ on $`C`$, we obtain a torsion sheaf $`\iota _!L_C`$ on $`S`$. $`\text{ch}(\iota _!L_C)`$ can be computed again from the G.R.R. formula:
$$\text{ch}(\iota _!L_C)=\iota _{}\left(\text{ch}(L_C)\frac{1}{\text{Todd}(N_C)}\right)=\iota _{}[C]+\left(\mathrm{deg}(L_C)\frac{1}{2}\mathrm{deg}(N_C)\right)[\text{pt}],$$
(4.13)
where $`\mathrm{deg}(L):=c_1(L)_C`$ for a line bundle $`L`$ on $`C`$. The RR charge of the torsion $`𝒪_S`$-module $`\iota _!L_C`$ can then be computed as
$$v_S(\iota _!L_C)=\iota _{}[C]+\chi (L_C)[\text{pt}]H_2(S)H_0(S),$$
(4.14)
where we have used the classical adjunction and the self-intersection formulae on $`S`$:
$`2g(C)2`$ $`=[C][C][C]c_1(S)_S,`$ (4.15)
$`\mathrm{deg}(N_C)`$ $`=[C][C]_S,`$ (4.16)
as well as the classical Riemann–Roch formula on $`C`$:
$$\chi (L_C)=h^0(C,L_C)h^1(C,L_C)=\mathrm{deg}(L_C)+1g.$$
(4.17)
Let us next turn to intersection pairings on D-branes. It seems that the most appropriate intersection form on D-branes on $`S`$ could depend on one’s purpose. Below, we will describe three candidates, each of which we think has its own reason to be chosen as an intersection form.
The first uses the Mukai vector $`v_S`$ (4.11) and defines a symmetric form:
$`I_S(_1,_2)`$ $`=v_S(_1)^{}v_S(_2)_S`$
$`={\displaystyle \frac{1}{12}}r_1r_2\chi (S)+r_1\text{ch}_2(_2)+r_2\text{ch}_2(_1)c_1(_1)c_1(_2)_S,`$ (4.18)
where $`\text{ch}()=r[S]+c_1()+\text{ch}_2()`$, $`\chi (S)=c_2(S)_S`$ the Euler number of $`S`$, and $`v^{}=v_0+v_1v_2`$ , with $`v_i`$ being the $`H_{2i}(S)`$ component of $`v`$. It should also be noted that $`I_S`$ is not $`𝒁`$-valued in general.
The second is the skew-symmetric form $`I_Xf_!`$ induced from the one on the ambient Calabi–Yau threefold $`X`$ (4.2). As shown below, however, this form has an description intrinsic to $`S`$ independent of the details of embedding:
$`I_X(f_!_1,f_!_2)`$ $`=f_{}(v_S(_1))^{}f_{}(v_S(_2))_X`$
$`=r_1c_1(_2)c_1(S)r_2c_1(_1)c_1(S)_S,`$ (4.19)
where we have used the self-intersection formula: $`f^{}f_{}[S]=c_1(N_S)`$ \[37, p. 431\], as well as the projection formula (4.10) to show that for $`[S]H_4(S)`$, $`[C]H_2(S;𝑸)`$
$$f_{}[S]f_{}[C]=f_{}\left([C]f^{}f_{}[S]\right)=f_{}\left([C]c_1(S)\right).$$
(4.20)
To be explicit, consider $`S=𝑷^2`$. $`H_2(𝑷^2)`$ is isomorphic to $`𝒁`$, the ample generator of which we denote by $`l`$. Then $`c_1(𝑷^2)=3l`$, and $`ll_{𝑷^2}=1`$. Following Diaconescu and Gomis , we express the Chern character of a coherent sheaf $``$ on $`𝑷^2`$ as
$$r()=n_2,c_1()=n_1l,\text{ch}_2()=n_0[\text{pt}],$$
(4.21)
where $`n_1,n_2𝒁`$, and $`n_0\frac{1}{2}n_1+𝒁`$. Our second intersection form can be written in these variables as follows:
$$I_X(f_!,f_!^{})=3(n_1^{}n_2n_1n_2^{}),$$
(4.22)
which is precisely the intersection form in up to sign. An interesting remark is that the intersection form (4.22) introduced in is based on that of one-cycles on an $`\widehat{E}_6`$ torus contained in $`𝑷^2`$, while our $`I_Xf_!`$ is induced from that on a Calabi–Yau threefold $`X`$ which contains $`𝑷^2`$.
The third, which has been used in to identify the RR charge lattice $`H_2(S)`$ with $`S`$ a del Pezzo surface and the string junction charge lattices, generalizing the result in , would be the most natural one also from the mathematical point of view :
$`\chi _S(_1,_2)`$ $`={\displaystyle \underset{i=0}{\overset{2}{}}}(1)^idim\text{Ext}_S^i(_1,_2)`$
$`=\text{ch}(_1)^{}\text{ch}(_2)\text{Todd}(T_S)_S`$
$`=r_1r_2\chi (𝒪_S)+r_1\text{ch}_2(_2)+r_2\text{ch}_2(_1)c_1(_1)c_1(_2)_S`$
$`+{\displaystyle \frac{1}{2}}r_1c_1(_2)c_1(S)r_2c_1(_1)c_1(S)_S.`$ (4.23)
The skew-symmetric part of the third form $`\chi _S`$ coincides with $`I_Xf_!`$:
$$\chi _S(_1,_2)\chi _S(_2,_1)=I_X(f_!_1,f_!_2),$$
(4.24)
while the relation between the symmetric part of $`\chi _S`$ and $`I_S`$ becomes
$$\frac{1}{2}\left(\chi _S(_1,_2)+\chi _S(_2,_1)\right)=\frac{1}{12}r_1r_2c_1(S)^2_S+I_S(_1,_2).$$
(4.25)
In view of the Serre duality: $`\text{Ext}_S^i(_1,_2)\text{Ext}_S^{2i}(_2,_1K_S)^{}`$, the skew-symmetric part of $`\chi _S`$ (4.24) comes from the non-triviality of the canonical line bundle $`K_S`$.
According to the Bogomolov inequality, the discriminant of a sheaf $``$, defined by
$$\mathrm{\Delta }():=2r\text{ch}_2()+c_1()^2_S,$$
must be non-negative if $``$ is torsion-free<sup>4</sup><sup>4</sup>4 Roughly speaking, a torsion-free sheaf on surface $`S`$ is a sheaf of sections of a vector bundle with at worst point-like singularities.and semi-stable , which puts the following constraint on the self-intersection number of a torsion-free $`𝒪_S`$-module $``$ corresponding to a true BPS D-brane:
$$\chi _S(,)=r^2\chi (𝒪_S)\mathrm{\Delta }()r^2\chi (𝒪_S).$$
(4.26)
Let $`J_SH_2(S;𝑹)`$ be a Kähler class on $`S`$. The classical central charge of $``$ measured by $`J_S`$ then admits an expression intrinsic to $`S`$:
$$Z_{J_S}^{\text{cl}}()=e^{J_S}v_S()_S.$$
(4.27)
In particular, if $`J_S`$ is obtained as a restriction of a Kähler class $`J_X`$ on $`X`$, then (4.27) coincides with $`Z_{J_X}^{\text{cl}}(f_!)`$: the central charge measured on $`X`$ by $`J_X`$.
## 5 String junctions versus del Pezzo surfaces
### 5.1 del Pezzo surfaces
A del Pezzo surface is a surface the first Chern class of which is ample. Apart from $`\stackrel{~}{B}_1=𝑷^1\times 𝑷^1`$, they are obtained by blowing up generic $`N`$ points on $`𝑷^2`$ for $`0N8`$, which we call $`B_N`$ in this article. Our main interest is of course in the three cases $`N=8,7,6`$. The homology groups of $`B_N`$ are
$$H_2\left(B_N\right)=𝒁[B_N]𝒁l𝒁e_1\mathrm{}𝒁e_N𝒁[\text{pt}],$$
(5.1)
where $`l`$ represents the pull-back of a line of $`𝑷^2`$, and $`e_1,\mathrm{},e_N`$ the exceptional divisors, by which the first Chern class is written as $`c_1(B_N)=3l_{i=1}^Ne_i`$. We also note that the Picard group of $`B_N`$ is isomorphic to $`H_2(B_N)`$, which means that each element of $`H_2(B_N)`$ is realized as the first Chern class of a holomorphic line bundle on $`B_N`$ which is unique up to isomorphism. Intersection pairings on $`H_2(B_N)`$ are given by
$$ll_{B_N}=1,le_i_{B_N}=0,e_ie_j_{B_N}=\delta _{i,j}.$$
(5.2)
We list here some topological invariants:
$$c_1(B_N)^2_{B_N}=9N,c_2(B_N)_{B_N}=3+N,\chi (𝒪_{B_N})=1.$$
(5.3)
The gravitational correction factor in the Mukai vector (4.12) is given by
$$\sqrt{\frac{\text{Todd}(T_{B_N})}{\text{Todd}(K_{B_N})}}=[B_N]+\frac{1}{2}c_1(B_N)+\frac{1}{12}(15N)[\text{pt}]H_2(B_N;𝑸).$$
(5.4)
The degree of $`[C]H_2(B_N;𝑸)`$ is defined by $`d([C])=[C]c_1(B_N)_{B_N}`$. If we expand it as $`[C]=a_0l+_{i=1}^Na_ie_i`$, then we see $`d([C])=3a_0+_{i=1}^Na_i`$. It is also convenient to associate the following two quantities to a coherent $`𝒪_{B_N}`$-module $``$: let $`d()`$ be the degree of $`c_1()`$, and $`k()=\text{ch}_2()_{B_N}`$, in terms of which the Euler number of $``$ (4.7) can be expressed as
$$\chi ()=h^0(S,)h^1(S,)+h^2(S,)=r()+\frac{1}{2}d()+k().$$
(5.5)
It is known that the degree zero sublattice of $`H_2(B_N)`$ is isomorphic to the $`E_N`$ root lattice, the root system of which is composed of the self-intersection $`2`$ elements. For $`3N8`$, the $`N`$ simple roots can be chosen as
$$𝜶_i=e_ie_{i+1},1i<N,𝜶_N=le_1e_2e_3.$$
(5.6)
The fundamental weight $`𝒘^iH_2(B_N;𝑸)`$ is uniquely determined by $`𝒘^i𝜶_j_{B_N}=\delta _j^i`$, and $`d(𝒘^i)=0`$. Any $`[C]H_2(B_N)`$ then admits the following orthogonal decomposition into the degree and the $`E_N`$ weight:
$$[C]=\frac{d([C])}{9N}c_1(B_N)+\underset{i=1}{\overset{N}{}}\lambda _i([C])𝒘^i,$$
(5.7)
where the Dynkin labels are determined by $`\lambda _i([C])=[C]𝜶_i_{B_N}`$; for a coherent sheaf $``$, its Dynkin labels $`\lambda _i()`$ can be defined in the same way using $`c_1()H_2(B_N)`$, that is,
$$c_1()=\frac{d()}{9N}c_1(B_N)+\underset{i=1}{\overset{N}{}}\lambda _i()𝒘^i.$$
(5.8)
The third intersection form $`\chi _N:=\chi _{B_N}`$ (4.23) can then be expressed as
$$\chi _N(_1,_2)=r_1r_2+r_1k_2+r_2k_1+\lambda _1\lambda _2\frac{d_1d_2}{9N}+\frac{1}{2}\left(r_1d_2r_2d_1\right).$$
(5.9)
For more information on del Pezzo surfaces see, for example, .
$`B_N`$ has a natural one-parameter family of complexified Kähler classes $`J_S=tc_1(B_N)`$, with $`\text{Im}(t)>0`$, because $`c_1(B_N)`$ is an ample divisor. The degree of a curve defined above is nothing but the volume of it measured by the normalized Kähler class $`J_S=c_1(B_N)`$.
Using (4.11) combined with (5.4), it is now straightforward to compute the classical central charge of a D-brane $``$ measured by the Kähler class $`J_S=tc_1(B_N)`$:
$`Z_t^{\text{cl}}()`$ $`=e^{tc_1(B_N)}v_{B_N}()_{B_N}`$
$`=r()(9N)\left({\displaystyle \frac{t^2}{2}}{\displaystyle \frac{t}{2}}+{\displaystyle \frac{1}{12}}{\displaystyle \frac{3N}{9N}}\right)+d()t\chi ().`$ (5.10)
Recall that the exact quantum central charge $`Z()`$ yields the classical one evaluated above (5.10) modulo the instanton correction terms $`O(e^{2\pi it})`$. In particular, for the cases $`N=8,7,6`$, the instanton expansion of the central charge (3.9) of the D4-brane obtained in the previous section takes the form:
$$t_d=\left(\frac{t^2}{2}\frac{t}{2}+\frac{1}{12}\frac{3N}{9N}\right)+O(e^{2\pi it}),$$
(5.11)
in writing which we have noticed that $`4\mathrm{sin}^2\pi \alpha =9N`$ for $`\widehat{E}_{N=8,7,6}`$. The classical part of $`t_d`$ coincides with the one in (5.10). Therefore we can make the following identification of the quantum central charge of the D-brane $``$ measured by the Kähler form $`tc_1(B_N)`$ by
$$Z_t()=r()(9N)t_d+d()t\chi ().$$
(5.12)
We are now in a position to compare the two central charges of the one and the same $`\widehat{E}_N`$ theory; one obtained by the analysis of the SW periods (2.60), and the other by the geometric method (5.12). They are related under mirror symmetry through (3.8) so that we have the following dictionary between the charges, after a trivial rescaling of the former:
$`p`$ $`=d(),`$
$`q`$ $`=r(),`$ (5.13)
$`n`$ $`=\chi ()=r()+{\displaystyle \frac{1}{2}}d()+k().`$
In passing, we give a comment on the $`N=0`$ case. Upon a change of variables: $`t=\frac{1}{3}t_b+\frac{1}{2}`$, where the constant shift implies the existence of the NS B-field flux on $`𝑷^2`$ , (5.10) becomes
$$Z^{\text{cl}}()=e^{t_bl}\text{ch}()\sqrt{\frac{\widehat{A}(T_{𝑷^2})}{\widehat{A}(N_{𝑷^2})}}_{𝑷^2}=r()\left(\frac{t_b^2}{2}+\frac{1}{8}\right)+\frac{1}{3}d()t_bk(),$$
(5.14)
where $`l`$ is the ample generator of divisors. This is precisely the classical central charge on $`𝑷^2`$ treated as the $`𝒁_3`$ orbifold in .
### 5.2 $`𝑬_\mathrm{𝟗}`$ almost del Pezzo surface
An $`E_9`$ almost del Pezzo surface $`B_9`$ is a surface obtained by blowing up nine points of $`𝑷^2`$ which are the complete intersection of two cubics on $`𝑷^2`$. It has the structure of elliptic fibration $`\pi :B_9𝑷^1`$, which has twelve degenerate fibers leaving the total space non-singular for a generic choice of parameters, which we assume throughout the paper. As its name stands for, $`B_9`$ shares many properties with the del Pezzo surfaces $`B_N`$; to be explicit, among the formulae in the preceding subsection, (5.1)–(5.5) remain valid for $`B_9`$ if one simply sets $`N=9`$ there, as well as the definition of the first Chern class $`c_1(B_9)`$. It is then clear that the elements of $`H_2(B_9)`$ orthogonal to both $`c_1(B_9)`$ and $`e_9`$ generate the $`E_8`$ root lattice. Let $`[\text{F}]`$ and $`[\text{B}]`$ be the class of $`H_2(B_9)`$ defined by the fiber and a cross section of the fibration respectively. We can then make the following identification:
$$[\text{F}]=c_1(B_9)=3l\underset{i=1}{\overset{9}{}}e_i,[\text{B}]=e_9,$$
(5.15)
the intersection pairings of which we give here for convenience:
$$[\text{F}][\text{F}]_{B_9}=0,[\text{F}][\text{B}]_{B_9}=1,[\text{B}][\text{B}]_{B_9}=1.$$
(5.16)
As opposed to the case of del Pezzo surfaces, $`c_1(B_9)`$ is no longer an ample divisor; it is only a nef divisor with self-intersection zero, so that $`c_1(B_9)`$ alone cannot define a Kähler class on $`B_9`$. However, the structure of elliptic fibration suggests the following natural two-parameter family of the complexified Kähler classes on $`B_9`$:
$$J=t_1[\text{F}]+t_2\left([\text{F}]+[\text{B}]\right)=(t_1+t_2)c_1(B_9)+t_2e_9,$$
(5.17)
where the imaginary parts of $`t_1`$ and $`t_2`$ parametrize the volume of the curve $`[\text{B}]`$ and $`[\text{F}]`$ measured by $`\text{Im}(J)`$ respectively, which span the Kähler sub-cone of our model. A serious treatment of the stability condition of coherent $`𝒪_{B_9}`$-modules would face with the problem of subdivision of the Kähler cone, because the stability condition depends on the choice of Kähler class $`\text{Im}(J)`$ , which we will not discuss further in this article. By the way, another two-parameter family of Kähler classes $`\stackrel{~}{J}`$ treated in can be written in our notation as $`\stackrel{~}{J}=\stackrel{~}{t}_1[\text{F}]+\stackrel{~}{t}_2l`$.
The classical central charge of the D-brane represented by a coherent $`𝒪_{B_9}`$-module $``$ can be computed as:
$`Z^{\text{cl}}()`$ $`=e^{t_1[\text{F}]t_2([\text{F}]+[\text{B}])}\text{ch}()\left([B_9]+{\displaystyle \frac{1}{2}}[\text{F}]+{\displaystyle \frac{1}{2}}[\text{pt}]\right)_{B_9}`$
$`=r()\left({\displaystyle \frac{t_2^2}{2}}+t_1t_2{\displaystyle \frac{t_2}{2}}{\displaystyle \frac{1}{2}}\right)+d()(t_1+t_2)+b()t_2\chi (),`$ (5.18)
where $`d()`$ and $`b()`$ are the two integers defined by the decomposition of $`c_1()`$:
$$c_1()=b()[\text{F}]+d()([\text{B}]+[\text{F}])+\underset{i=1}{\overset{8}{}}\lambda _i()𝒘^i.$$
(5.19)
That is, they are obtained via $`b()=[\text{B}]c_1()_{B_9}`$, $`d()=[\text{F}]c_1()_{B_9}`$. In terms of these variables, the intersection form $`\chi _9:=\chi _{B_9}`$ can be written as
$$\chi _9(_1,_2)=r_1r_2+r_1k_2+r_2k_1d_1d_2b_1d_2b_2d_1+\lambda _1\lambda _2+\frac{1}{2}(r_1d_2r_2d_1).$$
(5.20)
A period $`\varpi `$ obeys the Picard–Fuchs differential equations: $`^{(1)}\varpi =0,^{(2)}\varpi =0,`$ with
$`^{(1)}`$ $`=\left(\vartheta _1z_1\left(\vartheta _1\vartheta _2\right)\right)\vartheta _1,`$ (5.21)
$`^{(2)}`$ $`=\vartheta _2\left(\vartheta _2\vartheta _1\right)z_2(\vartheta _2+\frac{1}{6})(\vartheta _2+\frac{5}{6}),`$ (5.22)
which can be obtained from the standard procedure of the local mirror principle using a realization of $`B_9`$ as a hypersurface in a toric threefold. It is easy to see that we can take the two periods of the $`\widehat{E}_8`$ torus:
$$\varpi _0(z_2)={}_{2}{}^{}F_{1}^{}(\frac{1}{6},\frac{5}{6};1;z_2),\varpi _1^{(2)}(z_2)=\frac{1}{2\pi i}\left(\varpi _0(z_2)\mathrm{log}(\frac{z_2}{432})+{}_{2}{}^{}F_{1}^{}{}_{}{}^{}(\frac{1}{6},\frac{5}{6};1;z_2)\right),$$
as two of the four periods of the Picard–Fuchs system (5.21), (5.22) with $`t_i=\varpi _1^{(i)}/\varpi _0`$ $`(i=1,2)`$; the mirror map. As for the D4-brane period $`\varpi _2`$, we can take it to have the following form at the large radius limit:
$$t_d:=\frac{\varpi _2}{\varpi _0}=\left(\frac{t_2^2}{2}+t_1t_2\frac{t_2}{2}\frac{1}{2}\right)+O(e^{2\pi it_1},e^{2\pi it_2}),$$
(5.23)
because this is the classical part of the only four-cycle period, modulo addition of periods of lower dimensional cycles, that remains finite under the limit of infinite elliptic fiber of a Calabi–Yau threefold which contains $`B_9`$ as a section of elliptic fibration . Moreover, rewritten in the new variables $`𝒰=t_2`$, $`\stackrel{~}{\varphi }=t_1+t_2`$, (5.23) coincides with the period $`\stackrel{~}{\varphi }_D`$ of the Phase I in . Thus in terms of the basis of the solutions of the Picard–Fuchs equations: $`\{\varpi _0,\varpi _1^{(1)},\varpi _1^{(2)},\varpi _2\}`$, the quantum central charge of the coherent sheaf $``$ on $`B_9`$ measured by $`J`$ (5.17) can be expressed up to normalization factor as
$`Z()`$ $`={\displaystyle \frac{1}{\varpi _0}}\left(r()\varpi _2+d()(\varpi _1^{(1)}+\varpi _1^{(2)})+b()\varpi _1^{(2)}\chi ()\varpi _0\right),`$
$`=r()t_d+d()(t_1+t_2)+b()t_2\chi ().`$ (5.24)
We leave further detailed investigation of this Picard–Fuchs system to a future work.
### 5.3 Duality maps
For the $`\widehat{𝑬}_N`$ 7-brane configuration in the type IIB side, where we restrict ourselves to the cases $`3N8`$ for simplicity, let us recapitulate a string junction $`𝑱`$ as given in (2.62):
$$𝑱=\underset{i=1}{\overset{N}{}}\lambda _i𝝎^i+q𝝎^q+p𝝎^p+n𝜹^{(1,0)},$$
(5.25)
where $`(\lambda _i)`$ is the $`E_N`$ weight vector, $`(p,q)`$ the asymptotic charge, and $`n`$ the grade of the junction. The following intersection form $`\mathrm{\Phi }_N`$ on the junction lattice is adopted in :
$$\mathrm{\Phi }_N(𝑱_1,𝑱_2)=\lambda _1\lambda _2+n_1q_2+n_2q_1+\frac{p_1p_2}{9N}+q_1q_2+p_1q_2.$$
(5.26)
As discussed in sections 2 and 3, we have two realizations of BPS states in the $`E_N`$ theories on $`𝑹^4\times S^1`$: either by a coherent $`𝒪_{B_N}`$-module $``$ in the type IIA side or by a string junction $`𝑱`$ in the type IIB side, which raises a natural question: what is the correspondence between coherent sheaves on a del Pezzo surface $`B_N`$ and string junctions in the $`\widehat{𝑬}_N`$ 7-brane background? An answer has been given by Hauer and Iqbal , who found that the following map<sup>2</sup><sup>2</sup>2Strictly speaking, $`𝑱()`$ as well as $`Z()`$ depend only on $`\text{ch}()`$, but the notation adopted here will cause no confusions.
$$\text{ch}()𝑱()=\underset{i=1}{\overset{N}{}}\lambda _i()𝝎^i+r()𝝎^q+d()𝝎^p\chi ()𝜹^{(1,0)}$$
(5.27)
induces an isomorphism of the junction lattice and the homology lattice $`H_2(B_N)`$, which we identify with the RR charge lattice of D-branes on $`B_N`$, that is,
$$\mathrm{\Phi }_N(𝑱(_1),𝑱(_2))=\chi _N(_1,_2).$$
(5.28)
The inequality constraint (4.26) imposed on the self-intersection of a torsion-free semi-stable sheaf $``$ corresponding to a BPS brane is then converted into $`\mathrm{\Phi }_N(𝑱,𝑱)q^2`$, which is surely a necessary condition for a junction $`𝑱()`$ to represent a BPS state. This may serve as a physical consistency check of the map proposed above (5.27).
The map (5.27) has been obtained by the inspection of the two intersection forms: $`\mathrm{\Phi }_N`$ on the $`\widehat{E}_N`$ junction lattice and $`\chi _N`$ on the homology lattice $`H_2(B_N)`$. On the other hand, we have identified the $`\widehat{E}_N`$ string junction charges with the central charges of D-branes on $`B_N`$ in (5.13), which leads us to define a natural map $`\rho _N`$ from the string junctions (5.27) to the D-brane central charges on the del Pezzo surface (5.12) measured by the Kähler class $`tc_1(B_N)`$ by
$$\rho _N(𝝎^q,𝝎^p,𝝎^i,𝜹^{(1,0)})=((9N)t_d,t,0,1),$$
(5.29)
so that we have the correspondence:
$$\rho _N(𝑱())=Z_t(),$$
(5.30)
which we propose as another evidence for the map (5.27). Note that the junctions carrying $`E_N`$ weights $`(𝝎^i)`$ are in the kernel of $`\rho _N`$, only because our Kähler class $`tc_1(B_N)`$ cannot see $`E_N`$ quantum numbers; it should not be so difficult to incorporate $`E_N`$ quantum numbers in the central charge $`Z()`$, see .
Now that the correspondence between coherent $`𝒪_{B_N}`$-modules and $`\widehat{E}_N`$ string junctions has been found, our next task is to explicitly construct the string junctions for various coherent $`𝒪_{B_N}`$-modules. We give here the images of the map (5.27) of a few basic coherent sheaves on $`B_N`$:
$`𝑱(𝒪_{B_N})`$ $`=𝝎^q𝜹^{(1,0)},`$ (5.31)
$`𝑱(\iota _!L_C)`$ $`={\displaystyle \underset{i=1}{\overset{N}{}}}\lambda _i([C])𝝎^i+d([C])𝝎^p\chi (L_C)𝜹^{(1,0)},`$ (5.32)
$`𝑱(𝒪_p)`$ $`=𝜹^{(1,0)},`$ (5.33)
where (5.31) is a D4-brane wrapped on $`B_N`$; $`\iota :CB_N`$ in (5.32) is an embedding of curve, and $`L_C`$ an line bundle on $`C`$, which represents a D2-brane wrapped on $`C`$ bounded with several D0-branes; finally $`𝒪_p`$ in (5.33) is called the skyscraper of length one with support at a point $`pB_N`$, which clearly corresponds to a D0-brane at $`p`$.
To be more explicit in D2-branes (5.32), we want to describe two typical examples here: For the first example, let $`C`$ be an exceptional curve, that is, a rational curve with self-intersection $`1`$, so that $`d([C])=1`$ by the classical adjunction formula (4.15). The totality of the exceptional curves spans the fundamental Weyl orbit of $`E_N`$: $`(\mathrm{𝟑},\mathrm{𝟐})`$, $`\mathrm{𝟏𝟎}`$, $`\mathrm{𝟏𝟔}`$, $`\mathrm{𝟐𝟕}`$, $`\mathrm{𝟓𝟔}`$, $`\mathrm{𝟐𝟒𝟎}`$, for $`3N8`$. Next we take $`𝒪_C`$ as a line bundle on it; then the corresponding $`\widehat{E}_N`$ string junction becomes
$$𝑱(\iota _!𝒪_C)=\underset{i=1}{\overset{N}{}}\overline{\lambda }_i𝝎^i+𝝎^p𝜹^{(1,0)},\overline{\lambda }\overline{\lambda }=\frac{10N}{9N}.$$
(5.34)
For the second example, we take an elliptic curve $`E`$ with $`[E]=c_1(B_N)H_2(B_N)`$, which is known as an anti-canonical divisor in $`B_N`$; we consider also a degree zero line bundle $`L_E`$ on it, which is parametrized by the Jacobian of $`E`$: $`\text{Jac}(E)E`$; it is easy to see that the corresponding string junction becomes
$$𝑱(\iota _!L_E)=(9N)𝝎^p.$$
(5.35)
It is possible to extend the above considerations to the $`\widehat{E}_9`$ theory. The string junction and intersection form in this case are given by
$`𝑱`$ $`={\displaystyle \underset{i=1}{\overset{8}{}}}\lambda _i𝝎^i+p𝝎^p+q𝝎^q+n𝜹^{(1,0)}+m𝜹^{(0,1)},`$ (5.36)
$`\mathrm{\Phi }_9(𝑱_1,𝑱_2)`$ $`=\lambda _1\lambda _2+p_1p_2+p_1q_2+q_1q_2+m_1p_2+p_1m_2+q_1n_2+n_1q_2.`$ (5.37)
Hauer and Iqbal have shown that the following map from the homology lattice $`H_2(B_9)`$ in the type IIA side to the $`\widehat{E}_9`$ string junctions in the type IIB side:
$$\text{ch}()𝑱()=\underset{i=1}{\overset{8}{}}\lambda _i()𝝎^i+d()𝝎^p+r()𝝎^q+b()𝜹^{(0,1)}\chi ()𝜹^{(1,0)},$$
(5.38)
again defines the isomorphism of the lattices:
$$\mathrm{\Phi }_9(𝑱(_1),𝑱(_2))=\chi _9(_1,_2).$$
(5.39)
According to our point of view, on the other hand, comparison of (5.38) with (5.24) leads to define a natural map $`\rho _9`$ from the $`\widehat{E}_9`$ string junctions to the $`B_9`$ central charges by
$$\rho _9(𝝎^q,𝝎^p,𝝎^i,𝜹^{(0,1)},𝜹^{(1,0)})=(t_d,t_1+t_2,0,t_2,1),$$
(5.40)
which again induces the correspondence
$$\rho _9(𝑱())=Z_{t_1,t_2}().$$
(5.41)
To sum up, what we have done can be succinctly shown in the commutative diagram:
where the horizontal arrow is the isomorphisms (5.27), (5.38) proposed in , while the remaining two are ours; the left being the central charge formulae (5.12), (5.24), and the right the correspondences (5.30), (5.41).
## 6 Conclusions
In this article we have first derived the central charge formula for the D3-brane probe theory in the $`\widehat{𝑬}_{N=8,7,6}`$ 7-brane backgrounds. Employing local mirror symmetry we then translate the central charge for the IIB string junctions into that for IIA D-branes on del Pezzo surfaces $`B_{N=8,7,6}`$. To make this precise we have compared the quantum central charge (modulo instanton corrections) with the classical central charge verified by the geometric analysis of D-brane configurations. As a result, we have shown that the duality maps (5.27), (5.38) between the homology lattice of the del Pezzo surface $`H_2(B_N)`$ and the $`\widehat{E}_N`$ string junction lattice, originally found in based on the isomorphism of the lattices $`\mathrm{\Phi }_N`$ (5.28), (5.39), can be naturally recovered from the correspondence of the string junctions and the del Pezzo central charges (5.30), (5.41) for $`N=9,8,7,6`$ and the $`E_N`$ singlet part. The cases for $`N5`$ are also worth being examined in detail, and will be analyzed in our subsequent paper . Since the $`E_{N4}`$ theories on $`𝑹^4\times S^1`$ reduce to 4D asymptotically free theories in the $`R0`$ limit it will be interesting to compare with a recent work in which $`SU(2)`$ gauge theory with fundamental matters on $`𝑹^4\times S^1`$ is investigated by compactifying the type II theory on the local $`𝑭_2`$.
For further directions in future study, let us note that the duality maps (5.27), (5.38) concern only with the RR charges on the del Pezzo side or with the junction charges on the 7-brane side, while an actual coherent $`𝒪_{B_N}`$-module $``$ or an $`\widehat{E}_N`$ string junction have moduli parameters in general; for example, on the del Pezzo side, the structure sheaf $`𝒪_{B_N}`$ is rigid, while a torsion sheaf $`\iota _!L_E`$ with support on an anti-canonical divisor $`E`$ clearly has moduli parameters. Therefore it would be quite interesting to establish the duality map between coherent $`𝒪_{B_N}`$-modules and $`\widehat{E}_N`$ string junctions including their moduli parameters.
Another important issue is to analyze the stability of D-branes on del Pezzo surfaces. The stability of BPS-branes is the subject of current interest . Under the map (5.25) it will be possible to study the stability of certain D-brane configuration on $`B_N`$ in terms of the corresponding junction configuration in the $`\widehat{𝑬}_N`$ 7-brane background. Then the most interesting task is to determine curves of marginal stability (CMS) of BPS states and follow their decay processes. CMS of BPS junctions may be worked out numerically. Our preliminary computation indicates that there appear infinitely many CMS on the $`u`$-plane and their patterns look quite different from the ones observed for ordinary 4D $`𝒩=2`$ $`SU(2)`$ gauge theories. We hope to report the results elsewhere in the near future.
Acknowledgements
S.K.Y. would like to thank Y. Yamada for interesting discussions. The research of Y.O. is supported by JSPS Research Fellowship for Young Scientists. The research of K.M. and S.K.Y. was supported in part by Grant-in-Aid for Scientific Research on Priority Area 707 “Supersymmetry and Unified Theory of Elementary Particles”, Japan Ministry of Education, Science and Culture.
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# WITTEN DEFORMED EXTERIOR DERIVATIVE AND BESSEL FUNCTIONS
## 1 Introduction
Early studies \[e.g.,, \] proposed a unifying scheme for special functions showing that some of these functions may originate from the same structure .For Bessel functions of concern here , generating functions of integer orders are representation ”states” of derivative and integral operators of arbitrary orders.More precisely we have the ”inner” structure
$`_{|m|}`$ $`=`$ $`{\displaystyle \frac{}{zz}}.{\displaystyle \frac{}{zz}}\mathrm{}\mathrm{}\mathrm{}.{\displaystyle \frac{}{zz}}.`$
$`_{|m|}`$ $`=`$ $`{\displaystyle }zdz.{\displaystyle }zdz\mathrm{}\mathrm{}\mathrm{}\mathrm{}..{\displaystyle }zdz.`$
$`_m\mathrm{\Phi }(z,t)`$ $`=`$ $`(t)^m\mathrm{\Phi }(z,t),mϵZ`$ (1)
$`\mathrm{\Phi }(z,t)`$ $`=`$ $`{\displaystyle \underset{n=\mathrm{}}{\overset{n=\mathrm{}}{}}}\varphi _n(z)t^n`$
where we extend the index m to negative values by introducing the symbol $`𝑑z`$ to denote a “truncated ” primitive i.e. in defining the integral we omit the constant of integration $`\frac{df}{dz}𝑑z=f`$ and where $`\varphi _n(z)`$ stands for the reduced Bessel function $`\varphi _n(z)=\frac{J_n(z)}{z^n}`$ of integer order. For the polynomials such as Hermite and Laguerre for instance ,the generating functions only involve the realization of the set N of positive integers,with slight modifications of the derivative operators to account for the conventions used in defining these polynomials.It is important to note that although this common structure only set up the z-dependance of the generating functions ,it is the “ dynamical ” part of the scheme so to say.The t dependence is simply set by imposing some given desired properties.For Bessel functions for example we require a “symmetry” between positive and negative integer indices that is $`J_n=(1)^nJ_n`$ ,while for the polynomials it is the natural property of orthonormality that is invoked.
In a recent paper we intuitively applied a mechanism to generate real numbers out of integers in order to unify ( reduced ) Bessel functions and showed by direct analytic computation that indeed , Bessel functions fit into the scheme and therefore integer orders are mapped to real orders ( $`\lambda `$ is real ) through the formula
$$\frac{J_{n+\lambda }(z)}{z^{n+\lambda }}=exp\left[\lambda \underset{mϵZ/(0)}{}\frac{_m}{m}\right]\frac{J_n(z)}{z^n}$$
(2)
Let us summarize the mechanism to see how it works to convert an integer into a real.Suppose we are given an abstract state $`nnϵZ`$ and a set of raising $`\mathrm{\Pi }_m,m>0`$ and lowering $`m<0`$ operators .Then it is easy to show , given that data , that the state $`n+\lambda >`$ is related to the state $`n>`$ through the following formula
$$n+\lambda >=exp\left[\lambda \underset{mϵZ/(0)}{}\frac{(1)^m\mathrm{\Pi }_m}{m}\right]n>$$
(3)
Fourier transforming the $`n>`$ state as
$$n>=_\pi ^\pi \frac{d\theta }{2\pi }e^{in\theta }\theta >$$
(4)
where the $`\mathrm{\Pi }_m`$ operators act on the $`\theta >`$ state by simple multiplication by the factor $`e^{im\theta }`$ we get
$`n+\lambda >`$ $`=`$ $`{\displaystyle _\pi ^\pi }exp\left[\lambda {\displaystyle \underset{mϵZ/(0)}{}}{\displaystyle \frac{(1)^m\mathrm{\Pi }_m}{m}}\right]e^{in\theta }\theta >{\displaystyle \frac{d\theta }{2\pi }}`$ (5)
$`=`$ $`{\displaystyle _\pi ^\pi }e^{i(n+\lambda )\theta }\theta >{\displaystyle \frac{d\theta }{2\pi }}`$
Compare states in 4 to states in 5. In deriving the last line use have been made of the known formula
$$\underset{m=1}{\overset{\mathrm{}}{}}(1)^m\frac{sinm\theta }{m}=\frac{\theta }{2},\pi <\theta <\pi $$
(6)
Formula 2 is a new formula ( to be added to the huge literature on Bessel functions ) which is shown to apply to Neumann and Hankel functions as well .It is to be noted that although formula 2 is shown to be true ,we didn’t know why the above mechanism should apply to Bessel functions.Our guess of the above relation was based on the following correspondence
$`\varphi _n(z)`$ $``$ $`n>`$
$`\varphi _{n+\lambda \text{ }}(z)`$ $``$ $`n+\lambda >`$
A hint to this correspondence came from the fact that we indeed have a set of raising and lowering operators $`\mathrm{\Pi }_m=(1)^m_mmϵN`$
$$(1)^m\frac{d^m}{(zdz)^m}\varphi _n(z)=\varphi _{n+m}(z)mϵN,nϵZ$$
(7)
and for negative $`m^{}_{}{}^{}s`$ we just have to replace derivative operators by integral operators defined in 1.That this correspondence works is quite intriguing and therefore further investigations of it are needed .In this paper we answer the point. In section 2 we will study the deformed exterior derivative on the punctured plane and will realize that the converting( or deforming) mechanism is closely tied to the deformation of the exterior derivative .In section 3 we will demonstrate directly that real order reduced Bessel functions are the deformed versions of integer order reduced Bessel functions in the same fashion as $`d_\lambda `$ is the deformed version of the exterior derivative $`d`$ .As a consequence formula 2 will show up more elegantly.
## 2 Deformed exterior derivative on $`R^2/(0)`$
Let $`M`$ be a Riemannian manifold of dimension n .let $`V_p`$ $`p=0,1\mathrm{}.n`$ be the space of $`p`$-forms .Let $`d`$ and $`d^\text{ }`$ be the usual exterior derivative which define the De Rham cohomology of $`M`$ and its adjoint.Let $`V`$ be a smooth function $`V:MR`$ $`orC`$ ( called prepotential in the language of topological quantum field theories ) and $`\lambda `$ a real number .Define
$$d_\lambda [V]=e^{\lambda V}de^{\lambda V}$$
(8)
Evidently we have $`d_\lambda ^2=d_\lambda ^2=0`$. E.Witten had shown that $`V`$ plays the role of a Morse function and his consideration of the system
$$H_\lambda =d_\lambda d_\lambda ^{}+d_\lambda ^{}d_\lambda $$
(9)
had led to a new proof of Morse inequalities.Let us note at this point that there exists another version of $`d`$-deformation which is related to the fixed point theorems for Killing vector fields
$$d_s=d+siK$$
where s is an arbitrary number and where $`K`$ is a killing vector field-the infinitesimal generator of an isometry of $`M`$.In this context $`K`$ is regarded as an operator $`iK`$ on differential forms acting by interior multiplication and hence maps a $`p`$-form into a $`(p1)`$-form .Since we are interested in functions ($`0`$-form) such a deformation is not relevant as $`d_s`$ coincide with $`d`$ on the space of functions.In this section and the subsequent section we investigate the simpler system
$$H=d_\lambda $$
and show that it gives informations on the index structure of Bessel functions.Let us note at once that the above system is topological in the sense that $`d_\lambda `$ ,like $`d`$ or $`d_s`$, can be defined purely in terms of differential topology without choosing a metric in $`M`$ .Now to proceed we need to know the appropriate form of $`V.`$The system in 9 has also been investigated , in another context , by Baulieu et all to get informations on topological invariants.Their analysis of topological quantum mechanics on the punctured plane $`R^2/(0)`$ had selected the prepotential $`V=k\theta `$ which we later generalized .We have shown that the most general prepotential compatible with the topology of the punctured plane ( first homotopy group $``$ $`Z`$ ) has necessary the form.
$$V(\theta )=k\theta +\varphi (\theta )$$
(10)
where $`\theta `$ is the polar angle on the plane , k a constant and $`\varphi (\theta )`$ any function but periodic , (recall that the polar angle $`\theta `$ is not a periodic function ).It is that form 10 that we plug into $`d_{\lambda \text{ }}`$.On the restricted space of functions which depend only on the angle ,the exterior derivative simplifies to $`d=d\theta _\theta `$( there is no r dependance on which d acts ) .Inserting the specific form of the prepotential $`V`$ into 8 and rewriting the twisted operator as $`d_\lambda =d\theta _\theta ^\lambda `$ we find
$$_\theta ^\lambda =e^{\lambda \varphi }_\theta e^{\lambda \varphi }+\lambda k=_\theta +\lambda k+\lambda _\theta \varphi $$
Fourier transforming the periodic function $`\varphi (\theta )`$
$$_\theta \varphi (\theta )=i\underset{mϵZ/0}{}\rho _me^{im\theta }$$
we get
$$_\theta ^\lambda =_\theta +i\lambda \underset{mϵZ}{}\rho _me^{im\theta }$$
(11)
where we inserted the constant $`k=i\rho _0`$ into the sum .In the punctured plane the operator $`_\theta `$ and $`_\theta ^\lambda `$ have the natural interpretation respectively of the winding number operator and the effective or perturbed winding number operator .We thus write them as $`W=i_\theta `$ and $`W_\lambda =i_\theta ^\lambda `$. We also introduce the operator $`\mathrm{\Pi }_m=e^{im\theta }`$ with evident action on the basis $`n>`$ defined in 4 . For the operator $`W`$ and $`\mathrm{\Pi }_m`$ we have $`Wn>=nn>`$ and $`\mathrm{\Pi }_mn>=n+m>`$.In this new basis 11 takes the form
$$W_\lambda =W+\lambda \underset{mϵZ}{}\rho _m\mathrm{\Pi }_m$$
This is an example of a very simple topological quantum mechanical system where $`W_\lambda `$ is the perturbed hamiltonian ,$`\mathrm{\Pi }_m`$ a set of operators responsible for the interactions, $`\lambda \rho _m`$ a set of coupling constants and $`W`$ is the unperturbed hamiltonian .The eigenstates of the effective winding number which we denote $`n,\lambda `$ $`,\rho >`$ are shown to be related to the unperturbed one ,through the formula
$$n,\lambda \text{ },\rho >=exp\left[\lambda \underset{mϵZ/(0)}{}\frac{\rho _m\mathrm{\Pi }_m}{m}\right]n>$$
(12)
Hermiticity of $`W_\lambda `$ restricts the real “ spectral ” function $`\rho `$ to be symmetric $`\rho _m=\rho _m`$ .Comparing with the previous result 3 we see that our application of the formalism to Bessel functions requires the simple choice of the function $`\rho _m=(1)^m`$ .
## 3 Relation of $`\varphi _{n+\lambda }`$to $`d_\lambda `$
To show the relation , first write the generating functions of integer and of real orders
$`\mathrm{\Phi }(z,t)`$ $`=`$ $`{\displaystyle \underset{n=\mathrm{}}{\overset{n=\mathrm{}}{}}}\varphi _n(z)t^n`$
$`t^\lambda \mathrm{\Phi }(z,t)`$ $`=`$ $`{\displaystyle \underset{n=\mathrm{}}{\overset{n=\mathrm{}}{}}}\varphi _{n+\lambda }(z)t^n`$ (13)
We have learned in the particular case of section 2 that eigenstates of $`d_\lambda `$ ( $`W_\lambda `$) are deformed versions of the eigenstates of $`d`$ ( $`W`$ )and that the deformation consists in converting the index $`n`$ into $`n+\lambda `$ .We therefore have to look for eigenstates of $`d`$ and of $`d_\lambda `$ .The generating function $`\mathrm{\Phi }(z,t)`$ is an eigenstate of the exterior derivative by use of the recursion formula ( fixed t )
$$d\mathrm{\Phi }(z,t)=(\frac{t}{z})^1\mathrm{\Phi }(z,t)dz$$
For the eigenstate of $`d_\lambda `$ , the function of interest to look at is $`e^{\lambda V}\mathrm{\Phi }(z,t)`$ .We will show that ,with an appropriate $`V`$ ,it is indeed an eigenstate of the deformed exterior derivative and in the same time generating function of real order Bessel functions.The judicious choice of the operator $`V`$ so as to identify this function with the generating function for real orders (remember that the same crucial point of which $`V`$ to choose arose in the last section) can be shown to be <sup>1</sup><sup>1</sup>1In choosing a ( differential ) operator for $`V`$ instead of a simple function as in 8 we have implicitly generalized the deformed exterior derivative on flat space.This is enough for our purpose .We do not however, know , the expression of the generalized $`d_\lambda `$ in the case of a general manifod $`M`$.Such expression should be defined so as to be covariant and not to depend on the metric on $`M`$ like $`d=dzD_z=dz_z.`$
$$V=\underset{mϵZ/0}{}\frac{_m}{m}$$
In fact we have
$`e^{\lambda V}\mathrm{\Phi }`$ $`=`$ $`\mathrm{exp}\left[\lambda {\displaystyle \underset{mϵZ/0}{}}{\displaystyle \frac{_m}{m}}\right]\mathrm{\Phi }`$
$`=`$ $`\mathrm{exp}\left[\lambda {\displaystyle \underset{mϵZ/0}{}}{\displaystyle \frac{(1)^m(t)^m}{m}}\right]\mathrm{\Phi }`$
$`=`$ $`t^\lambda \mathrm{\Phi }(z,t)`$
This is the generating function of real order Bessel functions.To come to the second line we applied the recursion formula 1 and from the second line to the last we put $`t=e^{i\theta \text{ }}`$and made use of 6.Then acting by $`d_\lambda `$ on $`e^{\lambda V}\mathrm{\Phi }(z,t)`$ we find
$$d_\lambda (e^{\lambda V}\mathrm{\Phi })=e^{\lambda V}d\mathrm{\Phi }=(t)^1dze^{\lambda V}(z\mathrm{\Phi }(z,t))=ϝ(z,t,\lambda )dze^{\lambda V}\mathrm{\Phi }$$
(15)
The function $`ϝ(z,t,\lambda )`$ is a more involved expression which we will not work out as this is not necessary .Inspection of the third term in 15 shows that $`e^{\lambda V}(z\mathrm{\Phi })\mathrm{\Phi }e^{\lambda V}\mathrm{\Phi }`$where the last proportionality comes from the result in 3 .
Thus using the fact that the generating function of integer orders $`\mathrm{\Phi }(z,t)`$ is an eigenstate of $`d`$ ,we show that the function $`e^{\lambda V}\mathrm{\Phi }`$ is indeed an eigenstate of $`d_\lambda `$ and generating function of real orders .
Expanding both sides of 3 as in 13 the above relationship extends to Bessel functions themselves as the operator $`V`$ acts only on the $`z`$ variable.Hence we recover the unifying formula
$$\varphi _{n+\lambda }(z)=\mathrm{exp}\left[\lambda \underset{mϵZ/0}{}\frac{_m}{m}\right]\varphi _n(z)$$
The method outlined in this section has the advantage of giving a new check to the unifying formula ,in addition it shades light on the inner structure of Bessel functions showing that the modern concept of deformed or twisted exterior derivative ,first introduced by Witten ( which has been at the origin of the launching of topological field theories ) is already encoded in the index structure of Bessel functions .
> Acknowledgment
I would like to personally thank the head of the high energy section at the Abdus Salam international centre for theoretical physics ICTP Dr S.Randjbar -Daemi for inviting me to the centre scientific activities at various occasions.
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# A Variant Approach to the Overlap Action
## I Introduction
There has been a great deal of interest recently in lattice fermion actions which implement an exact chiral symmetry without doubling. The one explicit realization of such an action is the overlap action of Neuberger. It obeys the simplest version of the Ginsparg-Wilson (G-W) relation. All published studies of the overlap action in four dimensions use the Wilson fermion action as their starting point. These actions are apparently very expensive to simulate. However, it is a common expectation that a better starting action would simplify the evaluation of an overlap action, if many fewer steps of iteration would compensate for the increased cost per step of the more complicated action. That expectation is realized by an action (or family of actions) I describe here. With it, some physics calculations using an overlap action in quenched approximations can be done without recourse to supercomputers. Ingredients in my scheme include a fat link to suppress coupling to dislocations and a free field action with a spectrum which resembles an overlap; much of the gain comes from the use of eigenmodes of the approximate action to begin the calculation of overlap eigenmodes.
These techniques are obviously inspired by the fixed point action program for constructing classically perfect fermion actions. The best action I have found is, however, not, as far as I know, a fixed point action of any renormalization group transformation.
Versions of these ideas have been presented many times in two dimensional models.
After setting some conventions in Sec. II, I will describe the candidate actions and their tests in Sec. III. I will then present in Sec. IV a calculation of the quark condensate in finite volume with one of the new actions, along the lines of the calculation of Ref. .
## II Notation and Conventions
I will call a generic lattice Dirac operator $`d`$ and the overlap Dirac operator $`D`$. The eigenvalues of the simplest implementation of a G-W action lie on a circle of radius $`r_0`$ and the (massless) overlap Dirac operator is
$$D(0)=r_0(1+\frac{z}{\sqrt{z^{}z}})$$
(1)
where $`z=d(r_0)/r_0=(dr_0)/r_0`$ and $`d(m)=d+m`$ is the massive Dirac operator for mass $`m`$ (i.e. $`r_0`$ is equivalent to a negative mass term.) The overall multiplicative factor of $`r_0`$ is a useful convention; when the Dirac operator $`d`$ is thought of as “small” and Eqn. 1 is expanded for small $`d`$, $`Dd`$. Apart from this overall factor of $`r_0`$, my conventions are those of Ref. .
The Hermitian Dirac operator for mass $`m`$ is defined as $`h(m)=\gamma _5d(m)`$ and the overlap Hermitian Dirac operator is denoted as $`H(m)`$. Specifically,
$$H(0)=r_0(\gamma _5+ϵ(h(r_0)))$$
(2)
where $`ϵ(x)`$ is the step function, $`ϵ=1`$ if $`x<0`$, $`ϵ=1`$ if $`x>0`$. I will refer to the argument of the step function as the “kernel” of the overlap.
It is convenient to define the massive overlap Dirac operator as
$$D(m)=(1\frac{m}{2r_0})D(0)+m$$
(3)
and then the squared massive Hermitian Dirac operator is
$$H(m)^2=m^2+(1\frac{m^2}{4r_0^2})H(0)^2.$$
(4)
The zero eigenvalue eigenmodes of $`H(0)`$ are chiral, with $`\varphi |\gamma _5|\varphi =\pm 1`$, and the nonzero eigenvalue eigenmodes of $`H(0)`$ come in pairs of equal and opposite eigenvalues. For these modes, $`\varphi _\lambda |\gamma _5|\varphi _\lambda =\varphi _\lambda |H(0)|\varphi _\lambda /(2r_0)=\lambda /(2r_0)`$.
In a background gauge field carrying a topological charge $`Q`$, $`D(0)`$ (and $`H(0)`$) will have $`Q`$ pairs of real eigenmodes with eigenvalues 0 and $`2r_0`$. In computing propagators (for example for $`\overline{\psi }\psi `$), it is convenient to clip out the eigenmode with real eigenvalue $`2r_0`$, and to define the subtracted propagator as
$$\stackrel{~}{D}(m)^1=\frac{1}{1\frac{m}{2r_0}}[D(m)^1\frac{1}{2r_0}].$$
(5)
## III The overlap with a better kernel
### A Ingredients
A good $`d`$ or $`h`$ for use as a kernel in Eq. 1 or 2 should already “look like” $`D`$ or $`H`$. This means that the eigenvalues of its eigenmodes should lie (approximately) on a circle and its low lying eigenmodes should be approximately chiral–in a spectroscopy calculation, the additive renormalization of the bare quark mass, as measured, for example, through the variation of the pion mass with bare quark mass, should be small. Most implementations of the overlap action in the literature use the Wilson action, which does not satisfy either of these criteria: the eigenvalues of the free Wilson action sit like beads on a string around a set of four circular arcs (with real parts of their eigenvalues ranging from 0 to 8), and the additive mass renormalization of the interacting theory is $`\mathrm{\Delta }am_q1`$ for simulations at lattice spacings near 0.15 fm. Use of the clover action instead of the Wilson action improves the chiral properties but does nothing for its (free field) eigenmode spectrum.
Improving the action involves two ingredients.
First, to improve the chiral properties of the action, thin links are replaced by fat links. Actions with fat links are already quite chiral as shown by their small mass renormalization and $`Z_A1`$ . In this work I have studied fat link actions with two blockings: The first uses APE-blocking:
$`V_\mu ^{(n)}(x)=`$ $`(1\alpha )V_\mu ^{(n1)}(x)`$ (6)
$`+\alpha /6{\displaystyle \underset{\nu \mu }{}}(V_\nu ^{(n1)}(x)V_\mu ^{(n1)}(x+\widehat{\nu })V_\nu ^{(n1)}(x+\widehat{\mu })^{}`$ (7)
$`+V_\nu ^{(n1)}(x\widehat{\nu })^{}V_\mu ^{(n1)}(x\widehat{\nu })V_\nu ^{(n1)}(x\widehat{\nu }+\widehat{\mu })),`$ (8)
with $`V_\mu ^{(n)}(x)`$ projected back onto $`SU(3)`$ after each step, and $`V_\mu ^{(0)}(n)=U_\mu (n)`$ the original link variable. I have mostly studied a large amount of fattening, $`\alpha =0.45`$, $`N=10`$, but have looked at the smaller value $`\alpha =0.25`$, $`N=7`$. The second blocking is a “hypercubic” blocking devised by A. Hasenfratz. The fat links of this action are confined to a hypercube, so this action is more local than the APE-blocked actions. The mean plaquette at $`\beta =5.9`$ for (0.45,10) APE blocking is about 2.98, for (0.25,7) APE blocking it is 2.88, and for the hypercubic blocking, 2.84.
I believe that what is important here is that the fat gauge links decouple the fermions from short distance structure in the gauge field, not that the links are fattened in a particular way.
The eigenvalues of a GW action lie on a circle. I determine the best $`d`$ by taking a free field test action and varying its parameterization to optimize its eigenvalue spectrum (in the least-squares sense) for circularity, for some $`r_0`$. It was convenient to let the parameter $`r_0`$ also be a free parameter. The simplest extension of a nearest neighbor action which can have its eigenvalues lying approximately on a circle is ”planar:” it has scalar and vector couplings $`S=_{x,r}\overline{\psi }(x)(\lambda (r)+i\gamma _\mu \rho _\mu (r))\psi (x+r)`$ for $`r`$ connecting nearest neighbors ($`\stackrel{}{r}=\pm \widehat{\mu }`$; $`\lambda =\lambda _1=0.170`$, $`\rho _\mu ^{(1)}=0.177`$) and diagonal neighbors ($`\stackrel{}{r}=\pm \widehat{\mu }\pm \widehat{\nu }`$, $`\nu \mu `$; $`\lambda =\lambda _2=0.061`$, $`\rho _\mu ^{(2)}=\rho _\nu ^{(2)}=0.0538`$. The constraint $`\lambda (r=0)=8\lambda _124\lambda _2`$ enforces masslessness on the spectrum, and $`1=2\rho _\mu ^{(1)}+12\rho _\mu ^{(2)}`$ normalizes the action to $`\overline{\psi }i\gamma _\mu _\mu \psi `$ in the naive continuum limit. The corresponding value of $`r_0`$ is 1.6. A plot of the free field eigenvalues on a finite lattice is shown in Fig. 1. Presumably the places where its eigenvalues are purely real (but not zero) correspond in the continuum limit to theories with various numbers of massless free fermions, exactly as for the Wilson action.
The action also includes a clover term, with its coefficient set to the tree-level value appropriate to this action of $`C_{SW}=1.029941`$. Thus $`d`$ has no $`O(a)`$ artifacts. All the links, including the ones in the clover term, are replaced by fat links. The gauge connections to the diagonal neighbors are the average of the two shortest path connections. The cost of the action is about 6.5 times that of the usual clover action.
While it is probably not germane to the present discussion, the dispersion relation of this action is improved compared to that of clover or Wilson actions.
No claims are made for uniqueness–or even optimality. The optimizations do not have to be done with enormous precision since the overlap formula itself deforms almost any kernel action into a chiral action. The amount of fattening is also a free parameter, and there will almost certainly be tradeoffs between ease of implementing the overlap and the desire that the action show good scaling behavior. At large levels of fattening, there appears to be little renormalization of the parameters of a free action, and so little tuning is required to produce actions which behave well in simulations. Thus, the overlap parameter $`r_0`$ will be kept at its free field value of 1.6.
### B Testing actions
In what follows I will focus on comparisons of the overlap action with the Wilson action as a kernel, the “Wilson overlap” and with the fat link planar action as a kernel, or “planar overlap.” I have looked at six actions in all. Most tests involve the Wilson overlap or the (0.45,10) APE-blocked planar overlap. I have briefly investigated an overlap action with the (0.45,10) APE-blocked fat link clover action as a kernel. I have also investigated planar overlaps with (0.45,10) APE-blocked links and with hypercube-blocked links. Finally, I tested the “Gaussian” action, which has been proposed as a candidate approximate FP action.
The step function is approximated by either the polar formula introduced by Neuberger
$$ϵ(z)ϵ_N(z)=z\frac{1}{N}\underset{j=1}{\overset{N}{}}\frac{1}{c_jz^2+s_j}$$
(9)
($`c_j=\mathrm{cos}^2(\pi (j+1/2)/(2N))`$, $`s_j=\mathrm{sin}^2(\pi (j+1/2)/(2N))`$), or by a fourteenth order Remes polynomial, following the work of Edwards, Heller and Narayanan. In practice, the polar formula works very robustly for the planar overlap, when I rescale $`h(r_0)`$ to $`h(r_0)/r_0`$ in Eq. 9 and take $`N=6`$ to 10. The Wilson overlap is much more delicate and I have used the Remes approximation exclusively for it. I rescaled the operator $`h(r_0)`$ by a factor of 1/2.5 to map its eigenvalues into a range where the Remes algorithm has small errors. The inverses of the terms in the sum of Eq. 9 are found using a multi-mass conjugate gradient routine.
Let’s begin by looking at the actions. The range of the action is computed using a variation of the calculation of Hérnandez, Jansen, and Lüscher: Shown in Figs. 2 and 3 is a comparison of Wilson and planar overlap $`\sqrt{|D(0)\chi |^2}`$ for a delta-function source $`\chi `$ at the origin, as a function of distance $`r=\sqrt{x_\mu ^2}`$, for the free field action as well as on a set of $`\beta =5.9`$ $`8^4`$ configurations. The two actions show similar exponential falloff with distance. Of course, these pictures do not show anything about the locality of the gauge connections in the actions, just the fermions. In perturbation theory the fat link action has a cloud of glue of size $`\sqrt{\alpha N/3}`$ convoluted over every fermion offset.
It might be relevant that the plaquette for the fat link action has a value $`\text{Tr}(1U_p)/3)12.97/3=0.01`$, which should be sufficient for the argument of Hérnandez, Jansen, and Lüscher to insure the locality of the overlap action.
All my tests of the overlap begin with finding the eigenmodes of $`H`$ with the smallest eigenmodes (actually the eigenmodes of $`H`$ diagonalized in a basis which is composed of the smallest eigenmodes of $`H^2`$). Low lying eigenmodes of $`h(m)`$ and $`H(0)`$ are found using an adaption of a conjugate gradient algorithm of Bunk et. al. and Kalkreuter and Simma; I modified a code originally written for staggered fermions by K. Orginos.
I have implemented most of the standard tricks for efficient evaluation of eigenmodes of $`H`$. I compute a set of the $`N_0`$ low lying eigenmodes of $`h(r_0)`$ (typically $`N_0=1020`$), and project them out during the calculation of the operator $`ϵ(h(r_0)))\chi `$.
A first sign that the fat link will make the overlap better behaved comes from looking at the eigenmodes of the kernel function $`h(r_0)`$. Fig. 4 shows a histogram of the ten smallest eigenmodes of $`h(r_0)`$ reconstructed from the ten smallest eigenmodes of $`h(r_0)^2`$ on a set of ten $`8^4`$ $`\beta =5.9`$ configurations. The fat link actions have many fewer small eigenmodes. Eigenmodes near zero are associated with dislocations, gauge configurations on which topological objects are about to disappear. It is difficult for approximations to the step function to process these modes. The absence of these modes in the spectrum of $`h(r_0)`$ is why the polar form of $`ϵ_N`$ can perform well, even for small $`N`$. There is a lot of discussion in the literature (see for example Ref. ) about improving overlap or domain wall fermions by using gauge actions which have fewer dislocations than the Wilson gauge action, but a fermion action which could not see dislocation would work just as well. A fat link action does just that. While the fat link clover action is as much as an improvement as the planar action, it has a much higher conditioning number than the planar action. This is reflected in the amount of work needed to find the eigenvalues (and will also affect the number of inner conjugate gradient steps needed to evaluate the step function): the average number of Rayleigh iterations needed to collect the ten lowest eigenmodes is 3600 for the Wilson action, 5100 for the fat link clover action, and 1900 for the (fat link) planar action.
Because $`[H(0)^2,\gamma _5]=0`$, eigenmodes of $`H(0)^2`$ can be simultaneously eigenmodes of $`\gamma _5`$. All eigenmodes are constructed in a basis of chiral eigenfunctions. In that basis, eigenmodes of $`H(0)^2`$ are expected to be doubly degenerate. This degeneracy can be used to test algorithms. A low order polar approximation to the step function does a poor job of producing degenerate pairs of eigenmodes of $`H(0)^2`$ for the Wilson overlap and at this level of test it must be discarded.
The calculation of the condensate in Sec. IV requires the quark propagator $`D(m)^1`$, with the contribution from chiral modes removed. I do this by taking the set of eigenmodes of $`h`$ and finding the two lowest eigenmodes $`H(0)^2`$ in a basis containing one state of each chirality. The chirality of the highest state identifies the chiral sector which has no zero modes. I then perform the inversion of $`H^2=D^{}D`$ in this sector. This calculation is also accelerated by finding and projecting out some large number ($`O(10)`$) of eigenmodes of $`H(0)^2`$ during this inversion.
A trick which I have not noticed to have been emphasized in the literature is actually the source of much of the gain in efficiency of the planar overlap. This is to begin the calculation of eigenmodes of $`H`$ by first finding a set of low lying eigenmodes of $`h(m)`$ for some potentially useful $`m`$, and beginning the calculation for $`H`$ using the eigenmodes of $`h`$ (rather than, say, beginning with a set of random vectors). The combination of the small number of small eigenmodes of $`h(r_0)`$ for the fat link action and the use of good trial functions makes the planar overlap quite efficient compared to the Wilson overlap.
Most of my tests are on a set of $`8^4`$ $`\beta =5.9`$ lattices. Let us consider a set of examples which illustrate the differences. In the planar overlap I take the $`N=10`$ polar approximation to the step function, scale the argument by $`r_0=1.6`$, and project out $`N_0=10`$ eigenmodes of $`h(r_0)`$ in the step function. The Wilson overlap uses the fourteenth-order Remes approximation, scales its argument by 1/2.5, and projects 20 eigenvalues of $`h(r_0)`$. In both cases I wish to find the lowest two eigenmodes of $`H(0)^2`$, regardless of chirality, and I begin by finding and utilizing the four smallest eigenmodes of $`h(0)`$. I investigate a configuration which happens to have topological charge $`Q=1`$. The planar overlap calculation needs 109 Rayleigh iterations and 2529 inner conjugate gradient steps to find the two lowest modes (at a cost of 6.5 times the equivalent Wilson action step). The Wilson overlap needs 1220 Rayleigh iterations, and 86,723 inner conjugate gradient steps, to reach the same level of accuracy (about $`10^5`$ for the eigenvalues). Had we wanted the ten smallest eigenmodes of $`H(0)^2`$, using 20 trial eigenmodes of $`h(0)^2`$, the cost would be 306 Rayleigh iterations and 8189 inner conjugate gradient steps for the planar overlap, or about 26 inner steps per Rayleigh iteration, while the Wilson overlap uses 5820 Rayleigh iterations and a million inner conjugate gradient steps (about 170 inner steps per iteration). This is a savings in computer time of about a factor of eighteen for the planar overlap compared to the Wilson overlap. On other configurations, the number of inner CG’s per Rayleigh step for the planar overlap ranges between 20 and 30, and the number of Rayleigh iterations ranges from 300 to 500.
The difference in the number of inner conjugate gradient steps is due to the fact that the Wilson $`h(r_0)`$ has many more small eigenvalues than the fat link planar action. The difference in the number of Rayleigh iterations arises because the initial trial vectors, eigenmodes of $`h(0)`$, are close to being eigenmodes of $`H(0)`$. I illustrate this with a scatter plot of the change in an eigenmode, plotted as a function of the value of the overlap eigenmode, in Fig. 5. A graph of $`\gamma _5`$ vs. eigenmode for this action, at various stages of the overlap calculation, is shown in Fig. 6. It is also useful to look at the scatter of chirality vs. eigenvalue for several trial actions which might be candidates for an overlap kernel. This is shown in Fig. 7. I show $`\gamma _5`$ vs. eigenmode for the Wilson action, non-perturbative thin link clover action (with $`C_{SW}=1.85`$), fat link planar action, and fat link Gaussian action. From this picture, one would suspect that the thin link actions would not provide good trial wave functions for the overlap.
I have tested another fat link action, a planar action with 7 APE-blocking steps and $`\alpha =0.25`$. I kept $`r_0=1.6`$ since tests of the fat link clover action with this level of fattening at $`\beta =5.9`$ also showed little additive mass renormalization. The number of small eigenvalues of $`h(r_0)`$ increases slightly compared to the (0.45,10) case, and the number of inner CG steps grows, correspondingly, to about about 30. But the eigenmodes of $`h(0)`$ do not seem to be as good a trial basis as they were for the fatter link action, and the number of Rayleigh iterations grows from about 10,000 to 12,000-15,000 for the same calculation as was done above. A set of pictures showing relevant information for this action is shown in Fig. 8. This action does not see some of the instantons that the (0.45,10) fattening saw. That is reflected in the scatter of the eigenvalues.
One might hope that a better trial function would help. A simple way of changing the trial wave function is to vary the mass $`m`$ in $`h(m)`$. Trials at $`m=0.05`$ and $`0.1`$ did not produce any dramatic changes. This is an obvious place for future work.
The planar action with hypercubic fat links shows nearly identical behavior. Both these actions need about fifty per cent more inner CG steps than the (0.45,10) planar action but are still a large improvement over the Wilson overlap–and one could still run them on work stations. They might be useful in contexts where a very fat link might have undesirable properties.
Finally, the Gaussian action with (0.45,10) APE blocking requires about half as many inner CG calls as the (0.45,10) planar action, typically 3000-5000. However, this action requires about 2.5 times as much storage and 2.5 times as much CPU time per inner CG step as the planar overlap action, since its couplings span a hypercube (80 neighbors rather than 24). Thus it does not seem to be competitive with the planar action without clever coding .
The (0.45,10) fat clover action, with $`C_{SW}=1`$, was also briefly investigated as a kernel for the overlap. It seemed to need 100,000-300,000 inner CG steps in the fiducial calculation, in 2500 to 8000 Rayleigh iterations, beginning from $`m=0`$ trial wave functions. (There is an extreme variation in these numbers from lattice to lattice tested.) Eigenmodes of the fat link clover action are apparently quite different from those of the corresponding overlap action. The number of inner CG steps reflects the greater conditioning number of $`h(r_0)`$.
These tests are certainly incomplete. It may be that I have simply written a very inefficient Wilson overlap, though most of the code for the two methods is common. My Wilson overlap is so expensive that it cannot be tuned, whereas it is easy to test variations of the planar overlap. The “design philosophy” of beginning with an action which is close to an overlap action certainly produces actions which are inexpensive enough to run on small computers (work stations in my case).
## IV The quark condensate in finite volume
In order to do a little physics in what is otherwise a paper about technique, I present a calculation of the chiral condensate in the quenched approximation with the (0.45,10) fattened planar overlap action. The method is (almost) exactly that of the pioneering calculation of Hernandez and Jansen and Lellouch: one computes the condensate in background gauge field configurations of fixed topology labeled by winding number $`\nu `$. The condensate is
$$\mathrm{\Sigma }_\nu =m\mathrm{\Sigma }^2V(I_\nu (m\mathrm{\Sigma }V)K_\nu (m\mathrm{\Sigma }V)+I_{\nu +1}(m\mathrm{\Sigma }V)K_{\nu 1}(m\mathrm{\Sigma }V))+\frac{\nu }{mV}$$
(10)
where $`m`$ is the quark mass, $`V`$ is the volume, and $`\mathrm{\Sigma }`$ is the infinite volume condensate. $`I_\nu (z)`$ and $`K_\nu (z)`$ are modified Bessel functions. In practice, one computes $`\mathrm{\Sigma }`$ “without topology,” by working in the chiral sector which has no zero eigenmodes (and doubling the result). Then one measures $`\mathrm{\Sigma }_\nu \nu /(mV)`$. In practice, the lattice number needs an additive renormalization: in zero mass, this is removed by the replacement of $`D^1`$ by $`\stackrel{~}{D}^1`$ as in Eq. 5. There can also be contribution which vary with the quark mass, so the lattice number will need to be fit to
$$\overline{\psi }\psi _{sub}=\frac{1}{V}\underset{x}{}\stackrel{~}{D}_{x,x}=(\mathrm{\Sigma }_\nu \nu /(mV))+Cm+\mathrm{}$$
(11)
One small difference between my calculation and that of Ref. is in the normalization of the propagator: they do not have the prefactor $`1/(1\frac{m}{2r_0})`$ of Eq. 5. It is hard to believe that the condensate should not be an odd function of the symmetry breaking term (the quark mass), and the data does not show anything but a simple linear dependence on the mass.
Finally, $`\mathrm{\Sigma }`$ is scheme-dependent, and an overall lattice-to-continuum regulator needs to be computed.
Warned by the results of Ref. , I restricted my calculation to the $`\nu =\pm 1`$ sector. I generated a set of 40 lattices at each of three volumes $`8^4`$, $`10^4`$, and $`12^4`$ at $`\beta =5.9`$. I filtered them to find candidate $`\nu =\pm 1`$ configurations using a pure gauge measurement of topological charge, which I had previously calibrated against a set of 10 $`8^4`$ lattices. I checked that these lattices in fact had $`\nu =\pm 1`$ during the calculation of the condensate; only one lattice failed this test. This obviously leaves out configurations which the fermion observable identifies as carrying topological charge, but the gauge observable does not, but the alternative is to process every configuration through the overlap program. I was left with a set of 10 $`8^4`$, 13 $`10^4`$, and 9 $`12^4`$ lattices on which I computed $`\overline{\psi }\psi `$. I computed propagators at five quark masses, $`am=0.001`$, 0.002, 0.005, 0.01, and 0.02 (again using a multi-mass Conjugate Gradient solver). I used twelve random sources per lattice (though I blocked data together before averaging). This needed about 40, 70, and 125 steps at the three volumes, at about 20-30 inner CG’s per outer step, to reach a fractional squared residue of $`10^{12}`$. As a check, I also calculated $`\mathrm{\Sigma }`$ from the GMOR relation $`_x\pi (x)\pi (0)=\mathrm{\Sigma }/m`$; as expected in the overlap it reproduced the direct calculation of $`\mathrm{\Sigma }/m`$ to within $`10^7`$ (see Ref. ).
My lattice results are shown in Fig. 9.
The data at each volume are quite correlated. A single-elimination jackknife fit to
$$\overline{\psi }\psi _{sub}=m\mathrm{\Sigma }^2V(I_1(m\mathrm{\Sigma }V)K_1(m\mathrm{\Sigma }V)+I_2(m\mathrm{\Sigma }V)K_0(m\mathrm{\Sigma }V))+Cm$$
(12)
gives (re-inserting the lattice spacing) $`\mathrm{\Sigma }a^3=0.00394(16)`$, $`C=0.304(8)`$. The lattice number is about two standard deviations higher than the number reported by Ref. at $`\beta =5.85`$ of $`\mathrm{\Sigma }a^3=.0032(4)`$. We can sharpen this disagreement by trading the lattice spacing $`a`$ for the Sommer radius $`r_0`$, using the interpolating formula of Ref. : $`\mathrm{\Sigma }r_0^3=0.215(27)`$ for Ref. , $`\mathrm{\Sigma }r_0^3=0.354(14)`$ here.
A real comparison requires computing the $`Z`$ factor. I have not done that yet, but I can make a heuristic attempt at the calculation, by exploiting the fact that perturbation theory for fat link actions becomes simple in the limit of large fattening. This argument is implicit in the discussion of fat link perturbation theory in Ref. . Since it falls outside the main thrust of the paper, I relegate the discussion to an Appendix.
Taking $`r_0`$ to be 0.5 fm, $`\mathrm{\Sigma }_{\overline{MS}}(\mu =2`$ GeV) = 0.0244(10)(37) GeV<sup>3</sup>. The two errors are from the lattice fit and from an assumed lattice spacing uncertainty of five per cent.
This calculation produces a number which disagrees badly with a calculation of the condensate using clover fermions and the GMOR relation, $`\mathrm{\Sigma }_{\overline{MS}}(\mu =2`$ GeV) = 0.0147(8)(16)(12) GeV<sup>3</sup> . It is done at smaller lattice spacing and has completely different systematic uncertainties. Oddly enough, however, my result agrees well with an earlier analysis of data by Gupta and Bhattacharya: $`\mathrm{\Sigma }_{\overline{MS}}(\mu =2`$ GeV) = 0.024(2)(2) GeV<sup>3</sup> (statistical and lattice spacing uncertainties).
Of course, all the potential weaknesses of this calculation, reported by Ref. apply here, too: the lattices are small, and the lattice spacing is large. And it is a quenched calculation.
All of the computations of $`\mathrm{\Sigma }`$ were done over about a month of running on a set of four or five 450 Mhz Pentium II work stations.
## V Conclusions
None of the results shown here are particularly surprising, but that does not mean that they might be totally devoid of interest. A good approximate overlap action is easier to convert into an exact overlap, than any random action like the Wilson action. Crucial ingredients seem to be some kind of gauge connection which suppresses dislocations, a free action which “resembles” a free field overlap action, and the use of as much information from the trial $`h(m)`$ as possible to begin the calculation of the overlap action. One should be able to do better.
## Acknowledgements
I am deeply indebted to Urs Heller for many conversations about the overlap action, to Robert Edwards for a table of Remes coefficients, to Kostas Orginos for a copy of his eigenvalue code, to Anna Hasenfratz for discussion about blocking, and to Archie Paulson for help coding. This work was supported by the U. S. Department of Energy.
## A Z-factors for very fat actions
The evaluation of the Z-factor is straightforward. I will use many of the techniques of Ref. , but collect them here for completeness. We expect $`\mathrm{\Sigma }(\mu )_{\overline{MS}}=Z(a\mu )\mathrm{\Sigma }(a)_{latt}`$ with $`Z_S=1/Z_m`$ and $`Z_m`$ is the quark mass renormalization factor. In perturbation theory, $`Z_m=1+bz`$, and
$$z=\mathrm{\Delta }_{latt}\mathrm{\Delta }_{cont}=6\mathrm{ln}(1/(\mu a))+x,$$
(A1)
where $`b=\alpha _s(q)/(3\pi )`$. I choose to use $`\alpha _s=\alpha _{\overline{MS}}(q^{})`$. I will use the plaquette to define a coupling $`\alpha _V(3.41/a)`$ (=0.16054 at $`\beta =5.9`$), which will then be converted using two-loop perturbation theory to $`\alpha _{\overline{MS}}(q^{})`$, to give $`\mathrm{\Sigma }(\mu =1/a)_{\overline{MS}}`$. I will then run this result to $`\mu =2`$ GeV, using the (inverse of the) two-loop running formula for the $`\overline{MS}`$ quark mass.
$`z`$ is the difference between a lattice integral
$$\mathrm{\Delta }_{latt}=_kI_{latt}F_N;$$
(A2)
and a continuum integral
$$\mathrm{\Delta }_{cont}=_{k,c}I_{cont}$$
(A3)
I have factored the lattice integrand into a thin link piece $`I_{latt}`$ and the form factor of the fat quark-gluon vertices $`F_N`$; for APE-smeared links $`F_N=(1\alpha \widehat{q}^2/6)^{2N}`$ where $`\widehat{q}^2=_\mu (4/a^2)\mathrm{sin}^2(k_\mu a/2)`$. I am using the notation $`_k=d^4k/(2\pi )^4`$ over the hypercube, $`\pi /a<k_\mu <\pi /a`$; $`_{k,c}=d^4k/(2\pi )^4`$. Introducing a gluon mass $`\lambda `$ to regularize its IR divergence, $`\mathrm{\Delta }_{cont}=5/23\mathrm{ln}(\lambda ^2\mu ^2)`$ (in $`\overline{MS}`$). Adding and subtracting continuum-like terms to isolate the divergence in the lattice integral, we write
$`\mathrm{\Delta }_{latt}=`$ $`{\displaystyle _k}(I_{latt}3{\displaystyle \frac{\theta (\pi ^2k^2)}{k^4}})F_N3{\displaystyle _k}{\displaystyle \frac{\theta (\pi ^2k^2)}{k^4}}(1F_N)+3{\displaystyle _k}{\displaystyle \frac{\theta (\pi ^2k^2)}{k^4}}`$ (A4)
$`=`$ $`J3D+3\mathrm{ln}{\displaystyle \frac{\pi ^2}{a^2\lambda ^2}}.`$ (A5)
The first two terms ($`J`$ and $`D`$) are IR and UV finite. Now the point is that $`D`$ is a continuum-like integral, and all the dependence of $`\mathrm{\Delta }_{latt}`$ on the lattice action is in the first term. But if the fattening is large, the form factor $`F_N`$, is nonzero only at tiny $`k`$, where it is unity. At that point, for all practical purposes, $`I_{latt}3/k^4`$, and so we expect that $`J`$ will be small. All the dependence of $`Z_m`$ on the lattice action, other than its form factor, vanishes when we set $`J=0`$ in Eq. A5. $`q^{}`$ (defined according to the prescription of Lepage and Mackenzie) is calculated by a similar procedure.
While I can’t check this approximation for the overlap action, at (0.45,10) fattening, a complete calculation for clover fermions gives $`x=5.58`$ and $`q^{}=1.07`$, while the $`J=0`$ result is $`x=5.44`$ and $`q^{}=1.10`$. A similar argument would predict that finite renormalization factors, such as the vector and axial current renormalizations, are unity, which is a good approximation to what is seen at large fattening. Of course, for small fattening, this approximation fails badly and a better (non-perturbative?) calculation is necessary. (For example, the (0.25,7) APE-blocked clover action has $`x=0.975`$, $`q^{}=2.88`$; the approximate result is $`x=2.49`$, $`q^{}=1.59`$.)
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# Measurement of the 21Ne Zeeman frequency shift due to Rb-21Ne collisions
## I Introduction
Polarization of noble gases by spin exchange with optically pumped alkali metal vapors is an important technique for a variety of scientific and technological applications. Polarized noble gases been used for: precise measurements of the neutron spin structure function ; neutron polarizers and analyzers ; magnetic resonance imaging ; studies of gas diffusion in porous media ; gain media for Zeeman masers ; and tests of fundamental symmetries .
New tests of fundamental symmetries using polarized noble gases have recently been proposed. Kostelecky and Lane recently reported that a <sup>21</sup>Ne/<sup>3</sup>He differential magnetometer could be used to test local Lorentz invariance (LLI) . We have proposed to develop a <sup>21</sup>Ne/<sup>3</sup>He dual noble gas maser to carry out such an LLI test .
The spin exchange interaction between alkali and noble gas atoms has been extensively described in the literature . During spin exchange collisions, there is significant overlap of the alkali atom’s valence electron wavefunction with the noble gas nucleus which comprises a Fermi contact interaction of the form
$$H_{int}=\alpha \stackrel{}{K}\stackrel{}{S}$$
(1)
where $`\stackrel{}{S}`$ and $`\stackrel{}{K}`$ are the electron and nuclear spin angular momentum operators, respectively. This interaction leads to exchange of angular momentum between the alkali valence electron and the noble gas nucleus. The valence electron and nucleus experience during the collision a strong magnetic field from the close proximity of their spin magnetic moments. In spite of the small fraction of the time spent in collision, the average magnetic fields experienced by both the alkali valence electron and the noble gas nucleus can be affected if either the alkali or noble gas species is polarized. The resultant change in the Larmor precession frequency for either species, referred to as the contact shift, is expressed in terms of an enhancement factor $`\kappa `$. $`\kappa `$ is defined in terms of the *total* shift in the Larmor precession frequencies in co-located spherical distributions of magnetizations, induced both by the contact interaction and by the classical magnetic fields generated by the magnetizations. $`\kappa `$ is the ratio of that total frequency shift (contact plus classical magnetization) to the shift that would be induced by the classical magnetization alone. $`\kappa `$ can be much greater than unity, increasing with increasing $`Z`$ of the noble gas nucleus . In the limit where the formation of alkali-noble gas van der Waals molecules is a negligible part of the spin exchange interaction (true for the experiment reported in this paper), a single enhancement factor describes both the contact shift of the noble gas Zeeman splitting by the alkali atoms, and vice versa. . The frequency shift experienced by noble gas atoms in the presence of polarized alkali atoms, including both the contact shift and the shift due to classical magnetization fields, is
$$\delta \omega _{ng}=\left(\kappa 1+f_{DDF}\right)\gamma _{ng}\frac{8\pi }{3}g_S\mu _B\left[a\right]S_z$$
(2)
Recall the elementary result that $`\frac{8\pi }{3}M`$ is the magnetic field strength inside a magnetized sphere with magnetic dipole moment density $`M`$; the field is uniform and parallel to the magnetization. $`f_{DDF}`$ is a dimensionless factor that accounts for the effect of classical magnetization fields (”distant dipole fields” ), and is unity for a spherical distribution of magnetization. $`\left[a\right]`$ is the number per unit volume of alkali atoms, and $`\gamma _{ng}`$ is the gyromagnetic ratio of the noble gas species. For convenience, we define the component of alkali atom magnetization along the magnetic field (z) axis, $`M_a`$, as
$$M_a=\mathrm{}\gamma _a\left[a\right]S_z$$
(3)
where $`\gamma _a=g_S\mu _B\mathrm{/}\mathrm{}=`$ $`2\pi `$1.40 MHz/gauss is the electron gyromagnetic ratio. The alkali nuclear magnetic moments do not contribute to the contact shift, but do affect the distant dipole fields produced by the alkali magnetization distribution. We ignored these effects on the distant dipole fields, since they were too small to bear on the results reported in this work.
In precision measurements, it is important to be able to monitor the absolute polarization of the noble gases. Romalis and Cates demonstrated accurate polarimetry of <sup>3</sup>He using the Rb Zeeman frequency shift due to Rb-<sup>3</sup>He spin exchange collisions, and made a very precise measurement of the Rb-<sup>3</sup>He enhancement factor. Future use of <sup>21</sup>Ne in precision measurement applications would benefit from the use of contact shift polarimetry, but there is at present no published measurement of the Rb-<sup>21</sup>Ne enhancement factor. In this paper we present the first such measurement. We measured the net Larmor frequency shifts, simultaneously impressed by the presence of polarized Rb, on ensembles of polarized <sup>3</sup>He and <sup>21</sup>Ne precessing in the same \[nearly\] spherical glass cell. The measurement presented in this paper is not the first in which a ratio of contact shift enhancement factors was determined. Baranga *et. al.* measured the shifts of the Larmor precession frequencies of Rb and K vapors residing in the same cell, induced simultaneously by the presence of polarized <sup>3</sup>He .
The experiment is described in Section II. Determining the ratio of the Rb-induced shifts permitted us to divide out the dependence on Rb magnetization, and express the <sup>21</sup>Ne enhancement factor $`\kappa _{21}`$ in terms of the <sup>3</sup>He enhancement factor $`\kappa _3`$. The very accurate measurement of $`\kappa _3`$ by Romalis and Cates determined the absolute contact shift enhancement factor $`\kappa _{21}`$. Analysis of the data is presented in Section III. Suggestions for improved measurements are given in Section IV.
## II Apparatus, procedure, system parameters
We employed a <sup>21</sup>Ne/<sup>3</sup>He dual noble gas Zeeman maser to establish simultaneous precession of the <sup>21</sup>Ne and the <sup>3</sup>He in the same sealed glass cell. The <sup>21</sup>Ne/<sup>3</sup>He maser is a new device that will be described elsewhere; that description is summarized here. Figure 1 is a schematic representation of the maser. Circularly polarized light resonant with the Rb D1 transition (795 nm) is generated by a 30 W laser diode array (LDA); the light is directed onto a 1.4 cm dia \[near-\] spherical sealed glass cell containing 228 torr of <sup>21</sup>Ne, 266 torr of <sup>3</sup>He, 31 torr of N<sub>2</sub>, and Rb metal of natural isotopic composition . The LDA light propagates along the axis of a uniform magnetic field of strength 3 Gauss. The field is generated by a precision solenoid/$`\mu `$-metal-shield system capable of producing a 3 G magnetic field with $`35`$ $`\mu `$G/cm magnetic field gradients, with no gradient shimming. The main solenoid was driven using a current source capable of $`3`$ ppm stability over timescales of order 1000 sec. The glass cell was heated both by the LDA light and blown hot air to a temperature of ($`124.8\pm 1.9`$)<sup>o</sup>C (see Sec. III). The temperature of the air surrounding the cell was controlled to about 43 mK RMS. The total output power of the LDA was actively controlled to about a part in $`10^4`$.
The Rb was polarized by optical pumping; the resulting Rb electron polarization was very close to unity. The output power of the LDA was about 30W, so that even though it had a very broad ($`1.5`$ nm) linewidth, the on-resonant optical power was roughly 0.5 W. Thus, the optical pumping rate was $`>10^5`$ photon absorptions per sec per Rb atom. The rate of Rb polarization destruction events per Rb atom in the cell was $`<10^3`$ sec<sup>-1</sup> so that the Rb electron polarization was greater than 0.99. Pickup coils wound in series were placed in close proximity to the cell. The coils were excited by the magnetic flux associated with the precession of the polarized noble gases. The coils were part of a tuned resonant circuit with resonances at both noble gas precession frequencies. The large current flow induced in the pickup coils by resonant excitation resulted in magnetic fields oscillating at the noble gas nuclear Larmor frequencies being impressed back on the precessing atoms. This feedback resulted in steady state maser operation.
<sup>21</sup>Ne has nuclear spin 3/2, and thus a quadrupole moment. The quadrupole moments of the <sup>21</sup>Ne ensemble interact coherently with electric field gradients at the glass cell walls, shifting the four Zeeman energy levels so as to split the three \[otherwise degenerate\] dipole transitions . The splitting is proportional to $`\delta \omega _Q3\mathrm{cos}^2\theta 1`$, where $`\theta `$ is the angle of the cell’s symmetry axis with the magnetic field. We oriented the cell as close as possible to the ”magic angle” $`\mathrm{cos}\theta =\sqrt{1\mathrm{/}3}`$ to minimize the quadrupole splitting. The quadrupole wall shifts in the cell were not measured; however, computer simulations indicate that without proper cell orientation, realistic-sized splittings could disrupt <sup>21</sup>Ne maser oscillation. In any case, <sup>21</sup>Ne maser oscillation with steady-state amplitude was attained. We observed a small ($`10`$ $`\mu `$Hz) sinusoidal modulation of the <sup>21</sup>Ne maser frequency. The period of the modulation was about 16 hr. We speculate that the modulation was a consequence of remnant quadrupole splitting resulting from cell orientation error. In any case, the effect was small, with a period much longer than the duration of our contact shift measurements, and thus was not a source of significant systematic error.
The <sup>3</sup>He and <sup>21</sup>Ne maser signals were presented to a low-noise preamplier. The preamp output was then analyzed using lockin amplifiers, which were referenced to a set of synchronized signal generators tuned near the noble gas maser frequencies and locked to an H-maser-derived 5 MHz signal. The phases of both noble gas masers’ precession were recorded digitally at a sample rate of 1 Hz. The A/D sampling trigger was also derived from the master 5 MHz reference.
The measurements consisted of observing the maser phases during a series of reversals of the direction of the Rb polarization. The Rb polarization was reversed by rotation of the quarter-wave plates (Fig. 1) and was typically achieved in $``$ 20 sec. The total duration of a set of reversals was about 1000 sec., limited by magnetic field stability as well as the extent to which the noble gas masers’ steady states were disrupted by the alternating Rb polarization. Two sets of reversal data were acquired. Clear changes in the slopes of the phase curves were precisely correlated with reversal of the Rb polarization direction; these frequency changes were the result of the contact shift combined with the \[classical\] shift due to the Rb magnetization distribution. The change of noble gas maser oscillation frequency induced by Rb polarization reversal was about 44 mHz for <sup>21</sup>Ne maser and 71 mHz for the <sup>3</sup>He maser.
## III Data and analysis
Phase data for the second of the two measurement scans are shown in Figs. 2 and 3. We model the <sup>21</sup>Ne phase $`\phi _{21}`$ (see Fig. 2) and the <sup>3</sup>He phase $`\phi _3`$ (see Fig. 3) acquired by the lockin amplifiers as evolving according to the following equations:
$$\phi _{21}=\phi _{21,o}+\gamma _{21}B_ot2\pi \upsilon _{ref,21}+\gamma _{21}_0^t𝑑t^{}\delta B\left(t^{}\right)+\left(\kappa _{21}1+f_{DDF}\right)\gamma _{21}\frac{8\pi }{3}_0^t𝑑t^{}M_{Rb}\left(t^{}\right)$$
$$\phi _3=\phi _{3,o}+\gamma _3B_ot2\pi \upsilon _{ref,3}+\gamma _3_0^t𝑑t^{}\delta B\left(t^{}\right)+\left(\kappa _31+f_{DDF}\right)\gamma _3\frac{8\pi }{3}_0^t𝑑t^{}M_{Rb}\left(t^{}\right)$$
(4)
where: $`\gamma _{21}`$ ($`\gamma _3`$) is the gyromagnetic ratio of Ne (He); $`\kappa _{21}`$ ($`\kappa _3`$) is the enhancement factor for Ne (He); $`f_{DDF}`$ is the dimensionless factor to account for the classical magnetic field due to the Rb magnetization; $`\nu _{ref,21}`$ and $`\nu _{ref,3}`$ are the reference frequencies presented to the lockin amplifiers; $`B_o`$ is the time-average magnetic field strength over the time interval; $`\delta B`$ is the time-dependent deviation of the magnetic field from its mean value; and $`M_{Rb}`$ is the Rb magnetization. Note that independent maser frequency measurements determined the ratio of the two gyromagnetic ratios, $`\gamma _3\mathrm{/}\gamma _{21}`$, to an accuracy much better than 1 ppm. Also, the effects of magnetic fields due to noble gas magnetization fields on the phases have been neglected in comparison to the contact shift effects; this is reasonable for the noble gas magnetizations of this experiment (see below).
Magnetic field drift occurred over the course of the measurement scan. However, the <sup>21</sup>Ne and <sup>3</sup>He masers are co-located magnetometers, so that we can construct a data set from which the effects of magnetic field drift and Rb distant dipole fields have been removed. We define
$$\mathrm{\Delta }\phi \phi _3\frac{\gamma _3}{\gamma _{21}}\phi _{21}$$
(5)
It is easily shown that
$$\mathrm{\Delta }\phi =\phi _{3,o}\phi _{21,o}+2\pi \left(\frac{\gamma _3}{\gamma _{21}}\nu _{ref,21}\nu _{ref,3}\right)t+\left(\kappa _3\kappa _{21}\right)\gamma _3\frac{8\pi }{3}_0^t𝑑t^{}M_{Rb}\left(t^{}\right)$$
(6)
Aside from a trivial phase offset and known linear phase evolution, $`\mathrm{\Delta }\phi `$ reflects the action of the Rb contact shifts only: phase evolution due to magnetic fields is subtracted off. This combined phase profile is plotted in Fig. 4. As seen in Fig. 4, the phase evolution due to the contact shift was a piecewise-linear ”sawtooth” profile. Since our measurement is derived from computing the differences in slopes between adjacent linear regions of the curve, removal of constant and linear dependence from the phase data had no effect on the measurement. Thus, our plots of the phase data shown in Figs. 2, 3, and 4 are displayed with mean value and mean slope of zero.
If the Rb magnetization were known, one could immediately deduce $`\kappa _{21}`$ using eqn (6). Even though the Rb electron spin polarization was very close to unity, the density was not known *a priori*. It was therefore necessary to consider the ratio of contact shifts impressed on the $`\phi _{21}`$ and the $`\mathrm{\Delta }\phi `$ phase curves, and then determine the temperature self-consistently.
The data were analyzed in two ways which yielded essentially identical results. First, we computed the average changes in slope of the phase curves $`\phi _{21}`$ and $`\mathrm{\Delta }\phi `$, caused by the reversals of the Rb magnetization. That is, a least-squares linear fitting routine was used to extract the slope of each piecewise linear region of the phase curves. The differences in slope of adjacent regions was then recorded for each of the phase curves. The mean absolute values of the slope differences were then computed for each data set. We define $`M_{21}`$ as the mean of the absolute value of the slope differences for $`\phi _{21}`$, and $`\mathrm{\Delta }M`$ as that for the $`\mathrm{\Delta }\phi `$ profile. The uncertainties were estimated as the standard deviation of the mean in the slope differences’ absolute values.
We define the ratio of mean slope changes as $`M_{21}\mathrm{/}\mathrm{\Delta }M`$. The weighted average value for $``$ from the two data sets was $`=0.1262\pm 0.0025`$ where a one-sigma uncertainty is reported. The uncertainty was almost entirely due to that in $`M_{21}`$. The fractional uncertainty in $`\mathrm{\Delta }M`$ for both of the data sets was $`<0.1\%`$ This indicates that Rb polarization reversal was achieved to high precision and that the Rb magnetization was stable to better than 0.1$`\%`$ over the course of the measurements. The fractional uncertainty in $`M_{21}`$ was much larger, $`2\%`$, caused by the effects of magnetic field drift. In principle, we could have analyzed the $`\phi _3`$ profile rather than $`\phi _{21}`$, but magnetic field drift effects were $`\gamma _3\mathrm{/}\gamma _{21}=9.65`$ times larger on $`\phi _3`$ than $`\phi _{21}`$, making it impractical to extract slope data from $`\phi _3`$ directly (see Fig. 3).
In our analysis we did not account for possible variations in frequency due to changes in the noble gas longitudinal polarization that occurred in the course of a measurement. The magnetization of the noble gases created magnetic fields which varied in time during the course of the each of the two measurement scans, as evidenced by \[small\] maser amplitude changes observed during the scans. A given maser’s magnetization field could not affect its own frequency, since magnetization fields were parallel to the magnetization in the near-spherical cell. However, each maser’s frequency was shifted by the other’s magnetization, and such shifts did not subtract out of $`\mathrm{\Delta }\phi `$ as did magnetic field drift. The noble gas polarizations underwent both slow drift as well as variations correlated with the Rb polarization reversals. The small uncertainty in the $`\mathrm{\Delta }\phi `$ slope differences proves that frequency variations due to slow drift in the noble gas magnetizations had negligible effect. We also determined that the variations of the noble gas polarization caused by the periodic Rb polarization reversal had negligible effect on the contact shift measurement. Periodic reversal of the Rb polarization (i.e., ”square wave” modulation of the Rb polarization) induced a ”sawtooth wave” modulation (i.e., integrated square wave) on the noble gas polarizations of peak-to-peak amplitude approximately $`\gamma _{SE}\tau `$, where $`\gamma _{SE}`$ is the Rb-noble gas spin exchange rate per noble gas atom, and $`\tau `$ = time between Rb polarization reversals. Thus, the varying noble gas polarizations induced large changes in the maser frequencies over a measurement period: we estimate that a 2.8 mHz shift was induced on the <sup>3</sup>He maser by the <sup>21</sup>Ne polarization variation and a 0.7 mHz shift was induced on the <sup>21</sup>Ne maser by the <sup>3</sup>He polarization variation. However, given a constant time interval $`\tau `$ between Rb polarization reversals, the *average* frequency shifts due to the Rb-driven noble gas polarization variations did not change from one period to the next, thus they subtracted out in the calculation of the period-to-period frequency changes $`M_{21}`$ and $`\mathrm{\Delta }M`$. We conclude that noble gas polarization-induced magnetic fields had no systematic effect on the measurement of $``$.
It is easy to show from eqns (4) and (6) that the ratio $``$ can be interpreted as
$$=\left|\frac{\gamma _{21}}{\gamma _3}\frac{\kappa _{21}+\left(f_{DDF}1\right)}{\kappa _3\kappa _{21}}\right|$$
(7)
where it is seen that dependence on the Rb magnetization divided out. The value and estimated uncertainty for $``$ can then be used to compute the value and uncertainty of the desired contact shift enhancement factor for <sup>21</sup>Ne, $`\kappa _{21}`$, though the system temperature must be known because <sup>3</sup>He contact shift $`\kappa _3`$ is temperature-dependent. $`\kappa _{21}`$ and the system temperature must be determined self-consistently, and the effect of classical magnetization fields must be accounted for: these questions are addressed below.
We now consider the second method of analyzing the data to obtain a value for $``$, which yielded essentially the same results as the above approach. While this second method is slightly more complex, it is useful because it provides a framework for estimating the size of magnetic field drift in the present measurement. This second method of extracting a value of $``$ from the data considers the phase data sets to be $`N`$-component vectors (denoted by a tilde), where $`N`$ number of data in the set ($`10^3`$ for each of the two scans). We use the $`\mathrm{\Delta }\phi `$ profile as a model of the shape of the Rb magnetization-induced phase (see eqn 6). We then seek to describe the $`\stackrel{~}{\phi }_{21}`$ profile vector as a linear combination of elements of a partial basis of orthonormal vectors. The partial basis included an element corresponding to the \[normalized\] $`\mathrm{\Delta }\phi `$ profile ($`\stackrel{~}{V}_o`$). Additional basis vectors orthogonal to it, $`\left\{\stackrel{~}{V}_i\right\}_{i=1,N_b1}`$ were derived from $`\mathrm{\Delta }\phi `$ and Legendre polynomials $`\left\{P_i\right\}_{i=0,N_b1}`$ using Gram-Schmidt orthogonalization . The number of basis vectors $`N_b`$ was of order $`10^1`$, whereas the dimension of the vector space $`N`$ was of order $`10^3`$, so $`N_b<<N`$, and we find:
$$\stackrel{~}{\phi }_{21}=\underset{i=0}{\overset{N_b1}{}}a_i\stackrel{~}{V}_i$$
(8)
The coefficient of $`\stackrel{~}{V}_0`$ in this linear combination, $`a_0`$, was proportional to the value for $``$ obtained from that data set. The coefficients of the linear combination were determined using linear least squares fitting and by simple projection of $`\stackrel{~}{\phi }_{21}`$ onto the vectors of the partial basis, with identical results. The number of elements in the basis set was varied, but the coefficients obtained from the least squares fit remained the same regardless of the size of the basis set.
Using this second data analysis method, we thus determined two values of $``$, one from each data set. It was not possible to estimate uncertainties for the individual values of $``$ from the vector analysis method. Uncertainty estimates from the least squares analyses were unrealistically small because the vector analysis could not account for the effect of a non-zero projection of the magnetic field drifts along the Rb magnetization-induced phase profile vectors. We take the value of $``$ to be the mean of the results for the two data sets, and we coarsely estimate the uncertainty of $``$ as the deviation of the two values, to obtain $`=0.1287\pm 0.0036`$. This result agrees well with that of the first method of data analysis. For computing $`\kappa _{21}`$ we will use the result from the first method of analysis since the error estimate is probably more reliable.
We can compute the magnetic field drift profile by subtracting from $`\stackrel{~}{\phi }_{21}`$ its projection on the Rb magnetization basis profile $`\stackrel{~}{V}_0`$:
$$\stackrel{~}{\delta B}=\frac{1}{\gamma _{21}}\frac{d}{dt}\left(\stackrel{~}{\phi }_{21}\left(\stackrel{~}{\phi }_{21}\stackrel{~}{V}_o\right)\stackrel{~}{V}_o\right)$$
(9)
Note that the profile $`\stackrel{~}{\delta B}`$ is the magnetic field drift *orthogonal* to $`\stackrel{~}{V}_0`$. The drift profile from the second data set is plotted in Fig. 5; the time derivative of the \[discretely sampled\] profile was estimated by finite difference. Note that the field was stable to about $`\pm 3`$ ppm RMS.
Having obtained a value for $``$, we can then solve eqn (7) for $`\kappa _{21}`$ in terms of measured parameters, since $`\kappa _3`$ is known . However, $`\kappa _{21}`$ is expected to be temperature-dependent and so the temperature at which the experiment was carried out must be reported along with the value of $`\kappa _{21}`$. The Rb and noble gas temperature indicated by the hot air temperature control system was in error, because: the absolute calibration was poorly known; and the cell was heated by the LDA to a temperature higher than that indicated by the control system. Thus, it was necessary to deduce the noble gas/Rb vapor temperature from the contact shift data, and the known <sup>3</sup>He enhancement factor (with its known temperature dependence), in a self-consistent way. This was achieved by first noting that the mean slope change $`\mathrm{\Delta }M`$ is related to the enhancement factors and the Rb magnetization via
$$\mathrm{\Delta }M=2\gamma _3\frac{8\pi }{3}\left|\left(\kappa _3\kappa _{21}\right)M_{Rb}\left(T\right)\right|$$
(10)
where $`\left|M_{Rb}\left(T\right)\right|`$ (see eqn (3) above) is a known function of temperature. The Rb number density $`\left[Rb\right]`$ is related to temperature by
$$\left[Rb\right]=\frac{10^{9.318\frac{4040}{T}}}{\left(1.3810^{17}\right)T}cm^3\left(Tin{}_{}{}^{o}K\right)$$
(11)
The measured temperature dependence of the <sup>3</sup>He enhancement factor is
$$\kappa _3=4.52+0.00934\left[T\left({}_{}{}^{o}C\right)\right]$$
(12)
Substituting (11) and (3) into (10) yields an equation relating $`\kappa _{21}`$ and temperature (and known/measured parameters). Substituting (12) into (7) yields another equation relating $`\kappa _{21}`$ and temperature. This system of two \[nonlinear\] equations and two unknowns can be solved to yield the temperature and $`\kappa _{21}`$. The temperature determination depends strongly on the enhancement factor, while the enhancement factor determination depends only weakly on the temperature (via the known temperature dependence of $`\kappa _3`$). Thus, an iterative method to compute the temperature rapidly converges: an initial guess for the temperature enables calculation of $`\kappa _{21}`$, which in turn permits calculation of the Rb magnetization and thus the system temperature. This corrected temperature is used to re-compute a better estimate for $`\kappa _{21}`$, which is used to re-compute a better estimate of the Rb magnetization, etc.
Using the result $`=0.1262\pm 0.0025`$ from the slope difference analysis, we determine the contact shift enhancement factor for Rb and <sup>21</sup>Ne to be $`\kappa _{21}=31.8\pm 2.8`$ at a temperature of $`(124.8\pm 1.9)^o`$C. Note that the noble gas/Rb vapor temperature is much higher than the value of $`112^o`$C indicated by the hot air temperature control system, presumably the result of heating by the laser diode array. We are not aware of a previous measurement of the Rb-<sup>21</sup>Ne contact shift. Walker estimated the Rb-<sup>21</sup>Ne enhancement factor to be 38, and the ratio of the Rb-<sup>21</sup>Ne and Rb-<sup>3</sup>He enhancement factors to be 4.3 . We report this ratio to be $`5.6\pm 0.5`$.
The $`3\%`$ uncertainty in $``$ dominates other sources of error, inducing a $`9\%`$ uncertainty in $`\kappa _{21}`$. The contribution to the uncertainty in $`\kappa _{21}`$ resulting from the $`1.5\%`$ uncertainty in the He enhancement factor $`\kappa _3`$ is only $`1.5\%`$, which is negligible. The uncertainty in temperature induced a $`0.3\%`$ uncertainty in $`\kappa _{21}`$. To the extent that the cell could be described as an ellipsoid of revolution, its deviation from sphericality induced an uncertainty of $`0.15\%`$ in $`\kappa _{21}`$ with zero systematic correction.
The claim of zero systematic correction of the classical magnetization field shift for a ellipsoidal cell is justified by noting that: i) the classical magnetization field inside an ellipsoid of uniform magnetization is uniform, with its magnitude depending on the ellipsoid’s shape and the relative orientation of the ellipsoid major axis and the magnetization vector; and ii) when the ellipsoid major axis is oriented at the magic angle $`\theta _m=\mathrm{cos}^1\left(1\mathrm{/}\sqrt{3}\right)`$ with respect to the magnetization direction, the component of the classical magnetization field parallel to the magnetization is exactly that of a spherical distribution with the same magnetization, regardless of the ellipsoid’s shape. We oriented the cell at $`\theta _m`$ in order to suppress quadrupole wall shifts, with an uncertainty of about $`10^o`$. The $`10^o`$ uncertainty in orientation leads to the estimated uncertainty of $`0.15\%`$ in $`\kappa _{21}`$ due to measured asphericality of the cell.
## IV Suggestions for improved measurements; conclusions
There are straightforward improvements that could be implemented to reduce the uncertainty in the present measurement. The experiment was conducted in the course of developing a <sup>21</sup>Ne/<sup>3</sup>He dual Zeeman maser, so that opportunity for taking contact shift data was very limited. A factor of three improvement in the result could easily be achieved just by taking ten times more data. Active field stabilization using an independent magnetometer would essentially eliminate field drift, reducing statistical uncertainty in the measurement to that due to phase noise, a factor of $`>0.03/0.001=30`$ reduction. Active field stabilization would also permit a reduction in the phase noise since the signal acquisition bandwidth could be reduced by at least a factor of 3, so that the statistical uncertainty could be reduced by about a factor of 300 compared to the present experiment. This uncertainty would already be less than the uncertainty contributed by cell asphericality ($`0.15\%`$). However, the absolute uncertainty in the Rb-<sup>21</sup>Ne contact shift enhancement factor $`\kappa _{21}`$ would then be limited by the uncertainty in that for He, $`\kappa _3`$, $`1.5\%`$, whereas the ratio $`\kappa _{21}/\kappa _3`$ would be known to an uncertainty of only a few tenths of one percent.
It should be possible to measure ratios of alkali-noble gas enhancement factors to a level well below a part in $`10^3`$ if one could completely eliminate effects of the classical magnetization field shift. We have already shown that one can use differential magnetometry to obtain data from which classical magnetization field shifts are eliminated to high precision. It can be shown that using a three-species noble gas maser, one could measure the contact shift enhancement factor of a third species if the other species’ enhancement factors are known, *independent of the classical magnetization field shift*. Even though it has not yet been done, we believe it is a straightforward matter to implement a three-species maser in a single-bulb glass cell.
A comprehensive set of contact shift data would serve as a probe of noble gas-alkali interaction potentials . A program using multi-species noble gas masers to measure contact shift ratios could be implemented using <sup>129</sup>Xe, <sup>83</sup>Kr, <sup>21</sup>Ne, and <sup>3</sup>He along with K, Cs, and Rb in turn. Na would be a poor candidate for use in a noble gas maser because of its relatively low vapor pressure. The remaining stable noble gas isotope, <sup>131</sup>Xe, cannot be used with the other noble gases in a multi-species maser because its nuclear dipole moment has a sign opposite to that of the others (making it impossible to achieve a population inversion simultaneously with the other species). However, it may be possible to measure the <sup>131</sup>Xe precession frequency in free induction decay in the presence of two other noble gas masers. There are no other stable noble gas isotopes having non-zero nuclear spin.
Contact shift ratios would be converted to absolute values for enhancement factors by using the measured values for <sup>3</sup>He/Rb and <sup>3</sup>He/K . The enhancement factor for <sup>3</sup>He/Cs has not yet been measured.
Measuring contact shift ratios to a few parts in $`10^4`$ could proceed as follows: one would first measure the ratio of the enhancement factors of <sup>83</sup>Kr and <sup>129</sup>Xe to each other. The contact shift enhancement factors for <sup>83</sup>Kr and <sup>129</sup>Xe are estimated by Walker to be about 230 and 730, respectively . Using an independent magnetometer with the colocated Xe-Kr masers, the ratio of the \[large\] enhancement factors of Xe and Kr could be measured to a few parts in $`10^4`$ even if the classical magnetization field were accounted for only to the $`10\%`$ level (which can be done trivially- a $`1\%`$ accounting is claimed in this work). Triple masers using Kr and Xe, along with <sup>21</sup>Ne and <sup>3</sup>He in turn would yield high-precision measurements of the ratios of all the enhancement factors; the known enhancement factor for He would then determine the absolute values of the other enhancement factors. Temperature dependences could of course also be measured. In particular, enhancement factor ratios as a function of density (normalized to some experimental reference density) could be determined to very high precision. Improved cell thermometry would eliminate the need to depend on a temperature-density calibration to determine the cell temperature. Also, for very high precision measurements, control of the noble gas magnetization fields would have to be maintained more carefully than in the present work; in particular, the time between polarization reversals could be decreased in order to increase the relative size of the contact shifts in comparison to noble gas polarization induced frequency shifts. It might also be important to account carefully for the contribution of van der Waals molecule formation to the contact frequency shifts.
Finally, we remark that a two species maser using <sup>3</sup>He and <sup>21</sup>Ne has been proposed for a high precision test of local Lorentz invariance (LLI) , in which Rb-noble gas contact shifts would be an important systematic effect. Adding a co-located <sup>129</sup>Xe maser would provide a high precision *in situ* measurement of the Rb magnetization. With part-in-$`10^4`$ knowledge of the enhancement factors, the systematic effects of contact shifts on a search for LLI violation could be rendered negligible with the use of the <sup>129</sup>Xe maser ”co-polarimeter”.
We thank T. Chupp for the use of the <sup>21</sup>Ne/<sup>3</sup>He cell, and gratefully acknowledge D. Phillips and D. Bear for useful conversations.
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# HIP-2000-32/TH R-parity violation in (𝑡+𝑡̄)𝑔̃ production at LHC and Tevatron 11footnote 1The authors thank the Academy of Finland (project number 163394) for financial support.
ABSTRACT
We study the production of $`(t+\overline{t})\stackrel{~}{g}`$ at the hadron colliders in an R-parity ($`R_p`$) violating supersymmetric model. This process provides us with information not only about $`R_p`$ violation, but may also help us in detecting the supersymmetry itself. It is possible to detect an $`R_p`$ violating signal (with single gluino production) at the future hadron colliders, such as Fermilab Tevatron Run II or CERN Large Hadron Collider (LHC), if the parameters in the supersymmetric $`\text{/}R_p`$ interactions are not too small, e.g. for $`m_{\stackrel{~}{g}}=1`$ TeV, $`\lambda ^{^{\prime \prime }}=0.1`$, still hundreds of events are produced at LHC with luminosity $`30fb^1`$. Even if we could not detect a signal of $`\text{/}R_p`$ in the experiment, we get stringent constraints on the heavy flavour $`\text{/}R_p`$ couplings. In addition to the minimal supersymmetric standard model we have also considered some models with a heavy gluino as the lightest supersymmetric particle.
PACS number(s): 13.65.+i, 13.88.+e, 14.65.-q, 14.80.Dq, 14.80.Gt
I. Introduction
It is well known that the minimal supersymmetric standard model (MSSM) in its most general form contains lepton number ($`L`$) and baryon number ($`B`$) violating couplings. The resulting catastrophic proton decay can be avoided by imposing R-parity symmetry $`R_p`$,
$$R_p=(1)^{3B+L+2S},$$
(1)
where $`S`$ is the spin of the particle. In the models with R-parity conservation, superparticles can only be pair produced and the lightest supersymmetric particle (LSP) will be stable. However, $`R_p`$ conservation is not necessary for forbidding proton decays, instead of that, we just need either B-conservation or L-conservation . Models with $`R_p`$ violation ($`\text{/}R_p`$) can provide many interesting phenomena, such as neutrino masses and mixing. Partly because of that, $`\text{/}R_p`$ has attracted much attention . Many constraints from low-energy phenomenology are collected in Ref. .
The supersymmetry (SUSY) must be broken, since it has not been observed so far. Two kinds of breaking mechanisms of supersymmetry have been extensively studied phenomenologically, namely the minimal supergravity (mSUGRA) and minimal gauge-mediated SUSY breaking (GMSB) . These predict a different pattern of masses especially for the partners of gauge bosons, gauginos. Thus, finding the signals of gauginos is an important way to probe SUSY. In mSUGRA and GMSB, it is assumed that masses of gauginos will be unified at Grand Unified Theory (GUT) scale. From the evolution of parameters, gluinos should be the heaviest gauginos at low scale, since the ratios of gaugino masses to coupling constants do not change with scale in one-loop approximation . Thus gluinos can decay to other gauginos with jets in the $`R_p`$-conserving model .
However, heavy gluino as LSP (or NLSP) may still exist in some GMSB models, as S. Raby suggested in . He introduced the Higgs-Messenger mixing in GMSB model. This would lead triplet-doublet messengers split from triplet-doublet Higgs splitting. Since gluino obtains mass from SUSY-breaking induced by triplet messengers, heavier triplet messengers could suppress the mass of gluino so that gluino could be the LSP (or NLSP). The experimental limits on the masses of gluinos have been discussed in $`R_p`$-conserving model . In $`R_p`$-violating model, gluinos can decay through $`R_p`$-violating channels, which obviously changes both the detection strategy and present mass limits .
In the high energy hadron colliders especially squarks and gluinos will be produced plentifully, if they are not too heavy. It is hoped for that information on SUSY is found in Fermilab Tevatron Run II, or in the future CERN Large Hadron Collider, where squarks and gluinos with masses below 1.5 TeV should be detected . In mSUGRA and GMSB models, masses of squarks and gluinos are usually of the same order. However, in some special mechanism, such as O-II model , it has been suggested that squarks can be much heavier than gauginos. In this model the gluinos will be produced at much lower energies than squarks, so that it can be the first detected SUSY particle at hadron colliders. Since gluinos would be almost degenerate with the lightest neutralino and chargino, $`R_p`$ violating decay of gluinos will be significant .
The single production of gluinos, neutralinos and charginos has already been considered in the general case , and the single production of squarks, which is also significant to detect SUSY and $`R_p`$-violation, has been considered recently . In this work we will consider $`\text{/}R_p`$ at hadron colliders in the process
$`PP(P\overline{P})t+\stackrel{~}{g}(\overline{t}+\stackrel{~}{g}).`$ (2)
This process occurs via B-violating terms in the $`\text{/}R_p`$ model. In terms inducing heavy flavours, the $`\text{/}R_p`$ couplings can be very large from the present upper limits . For example, $`\lambda _{2ij}^{^{\prime \prime }}`$ and $`\lambda _{3ij}^{^{\prime \prime }}`$, getting their strongest constraints from the ratio of widths of $`Z`$ to leptons and hadrons, can be of order one ($`𝒪`$(1)) for the sfermion mass $`𝒪`$(100 GeV).
The pair production and decay of gluinos at hadron colliders have already been researched in Ref. . It has been shown that detecting gluinos is very difficult. It may become easier with an accompanying top quark in the process which we consider.
On the other hand, the single top quark production is an interesting topic itself . The single top quark production in the $`\text{/}R_p`$ model has been considered in Ref. . There, the possible cross section will depend on the $`\text{/}R_p`$ parameters as $`|\lambda ^{^{\prime \prime }}|^4`$, while the process we will consider depends on $`|\lambda ^{^{\prime \prime }}|^2`$. Although the mass of the gluino is not known presently, it is still possible to get stronger constraints on $`R_p`$-violating parameters from the $`(t+\overline{t})\stackrel{~}{g}`$ production.
In the following, we will give the analytical calculations of $`PP(P\overline{P})t\stackrel{~}{g}(\overline{t}\stackrel{~}{g})`$ in section II. In section III gluino and top decays are considered and in section IV the numerical results are presented. The conclusions are given in section V and some details of the expressions are listed in the appendix.
II. Production of gluinos in $`PPt\stackrel{~}{g}`$
The superpotential for $`\text{/}R_p`$ violating, but gauge and supersymmetry preserving interactions is written as
$$\begin{array}{ccc}W_{\text{/}R_p}\hfill & =\lambda _{[ij]k}L_i.L_j\overline{E}_k+\lambda _{ijk}^{^{}}L_i.Q_j\overline{D}_k+\lambda _{i[jk]}^{^{\prime \prime }}\overline{U}_i\overline{D}_j\overline{D}_k+ϵ_iL_iH_u\hfill & \end{array}$$
$`(2.1)`$
where $`L_i`$, $`Q_i`$ and $`H_u`$ are SU(2) doublets containing lepton, quark and Higgs superfields respectively, $`\overline{E}_j`$ ($`\overline{D}_j`$, $`\overline{U}_j`$) are the singlet lepton superfields (down-quark and up-quark). The square brackets around the generation indices $`i,j`$ denote antisymmetry of the bracketted indices.
We ignored the last term in Eq. (2.1) because its effects are assumed small in our process . Thus, we have 9 $`\lambda `$-type, 27 $`\lambda ^{^{}}`$-type and 9 $`\lambda ^{^{\prime \prime }}`$-type independent parameters left. The constraints on the couplings ,
$$|(\lambda \mathrm{or}\lambda ^{^{}})\lambda ^{^{\prime \prime }}|<10^{10}\left(\frac{\stackrel{~}{m}}{100\mathrm{GeV}}\right)^2.$$
$`(2.2)`$
is usually taken to indicate that only L- or B-number violating couplings exist.
In our work we will only consider the baryon number violating couplings, i.e., the third term in Eq. (2.1). In the following calculations we assume the parameters $`\lambda ^{^{\prime \prime }}`$ to be real.
We define the Mandelstam variables as usual
$$s=(p_1+p_2)^2=(p_3+p_4)^2,$$
$`(2.3.a)`$
$$t=(p_1p_3)^2=(p_4p_2)^2,$$
$`(2.3.b)`$
$$u=(p_1p_4)^2=(p_3p_2)^2.$$
$`(2.3.c)`$
The amplitude (Feynmann diagrams in Fig.1) of $`q_jq_k^{}\overline{t}\stackrel{~}{g}`$ is given by:
$$\begin{array}{ccc}M\hfill & =\hfill & M_s+M_t+M_u,\hfill \end{array}$$
$`(2.4)`$
with
$$\begin{array}{ccc}M_s=\hfill & & \overline{u^a}(p_3)(i\sqrt{2}g_sT_{\alpha \alpha ^{}}^a)(\mathrm{cos}\theta P_L+\mathrm{sin}\theta P_R)v_\alpha ^{}(p_4)\frac{i}{(sm_{\stackrel{~}{t}_1}^2+im_{\stackrel{~}{t}_1}\mathrm{\Gamma }_{\stackrel{~}{t}_1})}\hfill \\ & \times \hfill & \overline{v_\gamma ^c}(p_2)(2iϵ^{\alpha \gamma \beta }\lambda _{3kj}^{^{\prime \prime }})\mathrm{sin}\theta P_Ru_\beta (p_1)\hfill \\ & +\hfill & \overline{u^a}(p_3)(i\sqrt{2}g_sT_{\alpha \alpha ^{}}^a)(\mathrm{sin}\theta P_L\mathrm{cos}\theta P_R)v_\alpha ^{}(p_4)\frac{i}{(sm_{\stackrel{~}{t}_2}^2+im_{\stackrel{~}{t}_2}\mathrm{\Gamma }_{\stackrel{~}{t}_2})}\hfill \\ & \times \hfill & \overline{v_\gamma ^c}(p_2)(2iϵ^{\alpha \gamma \beta }\lambda _{3kj}^{^{\prime \prime }})\mathrm{cos}\theta P_Ru_\beta (p_1),\hfill \end{array}$$
$`(2.5.a)`$
$$\begin{array}{ccc}M_t\hfill & =\hfill & \overline{u^a}(p_3)(i\sqrt{2}g_sT_{\beta \rho }^aP_R)u_\rho (p_1)\frac{i}{(tm_{\stackrel{~}{d}}^2)}\overline{u_\alpha ^c}(p_4)(2iϵ^{\alpha \beta \gamma }\lambda _{3jk}^{^{\prime \prime }}P_R)u_\gamma (p_2),\hfill \end{array}$$
$`(2.5.b)`$
$$\begin{array}{ccc}M_u\hfill & =\hfill & \overline{u^a}(p_3)(i\sqrt{2}g_sT_{\beta \gamma }^aP_R)u_\gamma (p_2)\frac{i}{(um_{\stackrel{~}{d}}^2)}\overline{u_\alpha ^c}(p_4)(2iϵ^{\alpha \beta \rho }\lambda _{3kj}^{^{\prime \prime }}P_R)u_\rho (p_1),\hfill \end{array}$$
$`(2.5.c)`$
where $`P_{L,R}`$ are left- and right-helicity projections respectively, $`\theta `$ is the mixing angle of stop quarks (see Appendix for details) and $`\mathrm{\Gamma }_{\stackrel{~}{t}_{1,2}}`$ are decay widths of stop quarks $`\stackrel{~}{t}_{1,2}`$. The families of down-type quarks are marked by $`j,k=1,2,3`$ and the upper index $`c`$ means charge conjugate. In the calculations, we have neglected all mixing angles of scalar quarks except stop quarks. The amplitude depends on the $`R_p`$-violating parameters $`\lambda _{3jk}^{^{\prime \prime }}`$ ($`j,k=1,2,3`$) in the process.
The total cross section for the process $`q_jq_k^{^{}}\overline{t}\stackrel{~}{g}`$ is:
$$\begin{array}{c}\widehat{\sigma }(\widehat{s})=\frac{1}{16N_c^2\pi \widehat{s}^2}_{\widehat{t}^{}}^{\widehat{t}^+}𝑑\widehat{t}\overline{}_{spins}|M|^2,\hfill \end{array}$$
$`(2.6)`$
where $`\widehat{t}^\pm =\frac{1}{2}\left[(m_t^2+m_{\stackrel{~}{g}}^2\widehat{s})\pm \sqrt{\widehat{s}^2+m_t^4+m_{\stackrel{~}{g}}^42\widehat{s}m_t^22\widehat{s}m_{\stackrel{~}{g}}^22m_t^2m_{\stackrel{~}{g}}^2}\right]`$ and $`M`$ is the amplitude. Here we have neglected the masses of incoming quarks. $`N_c=3`$ is the color factor and the bar over summation means averaging over the initial spins.
In a similar way, the cross section for $`\overline{q}_j\overline{q}_k^{^{}}t\stackrel{~}{g}`$ can be calculated. The possible effects of $`q_jq_k^{^{}}\overline{t}\stackrel{~}{g}`$ and $`\overline{q}_j\overline{q}_k^{^{}}t\stackrel{~}{g}`$ should be observed in $`P\overline{P}`$ or $`PP`$ colliders. The cross section for the process $`P(P_1)P(P_2)(t+\overline{t})\stackrel{~}{g}X`$ can be obtained by convoluting the subprocess with quark distribution functions ,
$$\begin{array}{ccc}\sigma (s)\hfill & =𝑑x_1𝑑x_2f_i(x_1,Q)f_j(x_2,Q)\widehat{\sigma }(\widehat{s},\alpha _s(\mu ))\hfill & \end{array}$$
$`(2.7)`$
with $`p_1=x_1P_1`$, $`\tau =x_1x_2=\widehat{s}/s`$. $`f_{i,j}(x_n,Q)(n=1,2)`$ are the corresponding quark distribution functions of protons. We take $`Q=\mu =300`$ GeV. Similarly we can find numerically the cross section of $`P(P_1)\overline{P}(P_2)(t+\overline{t})\stackrel{~}{g}X^{}`$.
III. Top and gluino decays
In the process of Eq. (2), we have in the final state two heavy particles which will possibly decay inside the detector. Here we will shortly review the relevant decay modes of the top quark and gluino .
The top decays in the MSSM have been considered by several authors . The main decay mode is $`tbW`$, but $`tbH^+`$ can compete with it if mass of the charged Higgs is lighter than $`m_tm_b`$. Top quark decay to R-odd particles will also be important if those superparticles are light enough . However, in our case, with heavy squarks, the decays to real superparticles are impossible, except for the light LSP gluino with light squark . We plot the ratio $`\mathrm{\Gamma }_{\stackrel{~}{t}_1\stackrel{~}{g}}:\mathrm{\Gamma }_{bW}`$ as a function of the stop quark $`\stackrel{~}{t}_1`$ mass, $`m_{\stackrel{~}{t}_1}`$, in Fig.2 for the gluino mass $`m_{\stackrel{~}{g}}=30`$ GeV. In order to guarantee the purported standard top quark events at the Tevatron, $`BR(tbW)`$ should be larger than $`4050\%`$ as lower bound . So with the assumption of light LSP gluino (about $`2535`$ GeV), lower limit on the mass of stop quark can be obtained. For the top quark decay through $`R_p`$-violating interactions , the branching ratio of those decay modes will be very small compared with $`tbW`$ in our case. The only decay channel which we will consider in detecting top quark is the $`tbWbl\nu _l`$, where $`l=e,\mu `$. We will confine us to these decays (with $`BR(Wl\nu _l)22\%`$, its branching ratio in top quark decay should be at least $`8.811\%`$) since it is assumed that they have less background and are thus easier to detect than the hadronic decay modes.
We consider in this work a heavy gluino, which may be the LSP. The decay modes of a heavy gluino have been looked at e.g. in in $`R_p`$ violating case. The decay channel $`\stackrel{~}{g}q\stackrel{~}{q}`$ will dominate if kinematically allowed. A lighter than squarks gluino will decay to
$`\stackrel{~}{g}q\overline{q}\stackrel{~}{\chi }_i^0,q\overline{q}^{}\stackrel{~}{\chi }_j^\pm ,g\stackrel{~}{\chi }_i^0,`$ (3)
where $`\stackrel{~}{\chi }_i^0`$ and $`\stackrel{~}{\chi }_i^\pm `$ are neutralinos and charginos, respectively. The R-parity breaking decay modes become important for large $`\lambda ^{^{\prime \prime }}`$:
$`\stackrel{~}{g}q_iq_jq_k.`$ (4)
In Fig.3, we draw the $`\stackrel{~}{g}`$ decay with parameters in mSUGRA model. We consider the decay of gluinos in Fig.3 (a) and (b) with $`R_p`$ conservation and violation, respectively. In our calculations, only the two lightest neutralinos and the lightest chargino are considered, since the other neutralinos and the heavier chargino are too heavy for gluino to decay into. It is shown that the gluino branching ratio to two heavy quarks (top quark or bottom quark) and neutralinos or charginos (or jets from $`R_p`$-violating interactions) can be very large and increases with mass of gluino. This is reasonable because stop quarks and sbottom quarks can be much lighter than the other scalar quarks due to the large Yukawa couplings. Especially, in $`R_p`$-violating terms, the process through virtual stop quark $`\stackrel{~}{t}_1`$ will dominate the decay width because in the mSUGRA model it is the lightest squark and its mixing angle is near $`\pi /2`$. Thus, there are two $`b`$ (or $`t`$) quarks as a signal of gluino in our process (branching ratio can be read from Fig.3, as about $`6070\%`$). Combined with other top quark, the three heavy quarks, leads to the final state $`3b+n(l+\nu )`$ ($`n=1,2,3`$). This final state can be detected in the future CERN LHC and distinguished from background (with an assumed b-tagging efficiency $`ϵ50\%`$ in LHC ). For the much heavier gluinos which can decay directly to top quark and stop quark, decay width is shown in Fig.3 (c).
When gluino is the LSP, the only available decay modes are the R-parity violating ones. If the R-parity violating couplings are not exceedingly small, the gluino will decay in the detector through these channels. However, if the couplings are very small, the gluinos will form so-called R-hadrons before decaying (see e.g. ).
IV. Numerical results
In the mSUGRA model we are interested in the region of the parameter space where the gluino is lighter than the squarks, with the possible exception of the lighter stop quark. This choice allows gluino to be produced via the third generation $`R_p`$-breaking coupling with relatively large cross section and on the other hand it is not complicated by gluino decay to other squarks than possibly the lightest stop.
As a representative example of this part of the parameter space, we take $`m_0=1000\mathrm{GeV},A_0=1000\mathrm{GeV},\mathrm{tan}\beta =10`$ and $`sign(\mu )=+1`$. Varying $`m_{1/2}`$ suitably gives us the relevant gluino masses. The resulting masses for supersymmetric particles with the varied $`m_{1/2}`$ are listed in Table 1<sup>2</sup><sup>2</sup>2We thank A. Wodecki for providing us his program, which calculates the sparticle masses in MSSM. The program checks also some phenomenological constraints..
Table.1 MSSM parameters with $`m_0=1`$ TeV, $`A_0=1`$ TeV, $`\mathrm{tan}\beta =10`$ and $`sign(\mu )=+`$, units of values are GeV in the table.
| $`m_{\frac{1}{2}}`$ | $`m_{\stackrel{~}{g}}`$ | $`m_{\stackrel{~}{t}_1}`$ | $`m_{\stackrel{~}{t}_2}`$ | $`m_{\stackrel{~}{b}_1}`$ | $`m_{\stackrel{~}{q}}`$ | $`m_{\stackrel{~}{\chi }_1^0}`$ | $`m_{\stackrel{~}{\chi }_1^\pm }m_{\stackrel{~}{\chi }_2^0}`$ |
| --- | --- | --- | --- | --- | --- | --- | --- |
| 120 | 308 | 457 | 813 | 790 | 1020-1040 | 49 | 95 |
| 140 | 360 | 468 | 828 | 803 | 1028-1055 | 58 | 111 |
| 180 | 463 | 499 | 862 | 836 | 1060-1090 | 75 | 144 |
| 200 | 515 | 517 | 882 | 855 | 1075-1110 | 83 | 161 |
| 220 | 566 | 537 | 904 | 876 | 1100-1130 | 92 | 177 |
| 250 | 643 | 571 | 939 | 910 | 1120-1160 | 104 | 202 |
| 280 | 721 | 608 | 977 | 947 | 1160-1200 | 117 | 227 |
| 300 | 773 | 634 | 1004 | 974 | 1190-1230 | 125 | 243 |
| 330 | 850 | 675 | 1047 | 1016 | 1220-1270 | 138 | 268 |
| 340 | 876 | 689 | 1061 | 1031 | 1240-1280 | 142 | 277 |
| 350 | 902 | 703 | 1076 | 1045 | 1250-1300 | 146 | 285 |
| 370 | 954 | 733 | 1106 | 1076 | 1280-1330 | 155 | 302 |
| 390 | 1006 | 763 | 1138 | 1107 | 1310-1360 | 163 | 318 |
In Fig.4 (a), we show the cross section of $`P\overline{P}(t+\overline{t})\stackrel{~}{g}X`$ as a function of mass of gluino ($`m_{\stackrel{~}{g}}`$) at Tevatron Run II energy, i.e. with center-of-mass energy of the collision $`\sqrt{s}=2`$ TeV. There, the solid line corresponds to $`\lambda _{312}^{^{\prime \prime }}=\lambda _{313}^{^{\prime \prime }}=\lambda _{323}^{^{\prime \prime }}=1`$ and dashed line to $`\lambda _{312}^{^{\prime \prime }}=\lambda _{313}^{^{\prime \prime }}=\lambda _{323}^{^{\prime \prime }}=0.1`$. At LHC with $`\sqrt{s}=14`$ TeV, the cross section for $`PP(t+\overline{t})\stackrel{~}{g}X`$ as a function of $`m_{\stackrel{~}{g}}`$ is shown in Fig.4 (b) with $`\lambda _{312}^{^{\prime \prime }}=\lambda _{313}^{^{\prime \prime }}=\lambda _{323}^{^{\prime \prime }}=0.1`$ for solid line and $`\lambda _{312}^{^{\prime \prime }}=\lambda _{313}^{^{\prime \prime }}=\lambda _{323}^{^{\prime \prime }}=0.01`$ for dashed line. In the calculations, we have neglected the decay width of $`\stackrel{~}{t}_1`$, since mass of $`\stackrel{~}{t}_1`$ is far from the center-of-mass energy.
$`\mathrm{\Gamma }_{\stackrel{~}{t}_2}`$ including $`R_p`$-violating contribution is shown in Fig.5, where solid line corresponding to $`\lambda _{312}^{^{\prime \prime }}=\lambda _{313}^{^{\prime \prime }}=\lambda _{323}^{^{\prime \prime }}=1`$ and dashed line to $`\lambda _{312}^{^{\prime \prime }}=\lambda _{313}^{^{\prime \prime }}=\lambda _{323}^{^{\prime \prime }}=0.1`$. The results show that the cross sections can be very large if $`m_{\stackrel{~}{g}}`$ is smaller than $`400`$ GeV and $`\lambda ^{^{\prime \prime }}`$ are close to the present limits ($`\lambda ^{^{\prime \prime }}1`$) in the Tevatron Run II (with luminosity about $`2\mathrm{fb}^1`$, it corresponds to hundreds of events). At LHC with luminosity $`30\mathrm{fb}^1`$, the process can potentially be seen with mass of gluino less than $`1`$ TeV even if $`\lambda ^{^{\prime \prime }}0.01`$.
We have also considered production of the LSP heavy gluino. It has been shown in the model of that masses of gluinos below $`115`$ GeV are excluded except a narrow window between $`m_{\stackrel{~}{g}}=2535`$ GeV. We found that the cross section of the process in Tevatron Run I can be very large with masses of gluinos staying in that narrow window (about $`10`$ pb with $`m_{\stackrel{~}{g}}=30\mathrm{GeV},\lambda _{3ij}^{^{\prime \prime }}=1`$ and all $`m_{\stackrel{~}{f}}=150`$ GeV). Unlike in the discussion of , gluinos can decay through $`R_p`$-violating interactions (e.g. $`\stackrel{~}{g}bcs`$) if $`\lambda ^{^{\prime \prime }}`$ is nonzero. The gluino decay added by single production of top-quark, should have been visible already in Tevatron Run I using above results. Therefore, the narrow LSP gluino window can give much stronger constraints on $`R_p`$-violating parameters $`\lambda ^{^{\prime \prime }}`$ (with $`\lambda ^{^{\prime \prime }}=0.1`$, still several events are produced with luminosity $`19pb^1`$ in Fermilab), otherwise we can close this narrow LSP heavy gluino window.
In the mSUGRA model, the gaugino soft-SUSY-breaking masses are assumed to be equal at GUT scale and will lead to $`m_{\stackrel{~}{g}}:m_{\stackrel{~}{\chi }_1^\pm }:m_{\stackrel{~}{\chi }_1^0}7:2:1`$ at low energy. However, there are other mechanisms such as O-II model in which ratios of masses can be given as $`m_3:m_2:m_1(3+\delta _{GS}):(1\delta _{GS}):(33/5\delta _{GS})`$ at GUT scale, where $`\delta _{GS}`$ is the Green-Schwarz mixing term. Then at low energy, we can obtain light gluinos, which are almost degenerate with the lightest neutralino and chargino, with the heavy scalar quarks. As a representative of this model we take $`\delta _{GS}=4.1,m_0=1800\mathrm{GeV},m_{1/2}=10\mathrm{GeV},A_0=0,\mathrm{tan}\beta =3`$ and $`sign(\mu )=+1`$ (leading to $`m_{\stackrel{~}{g}}102`$ GeV and $`m_{\stackrel{~}{q}}1.5`$ TeV). Setting $`\lambda _{3ij}^{^{\prime \prime }}=1`$, the cross section of the process $`P\overline{P}(t+\overline{t})\stackrel{~}{g}X`$ can be about $`3.6`$ fb with $`\sqrt{s}=2`$ TeV and of the process $`PP(t+\overline{t})\stackrel{~}{g}X`$ about $`27`$ pb with $`\sqrt{s}=14`$ TeV. So in this model the single production of gluino can provide a very significant signal for detecting SUSY and $`R_p`$ violation. Detecting pair production of gluino in this model was already discussed in and it was shown that the $`R_p`$-violating decay of gluino will dominate if $`\lambda ^{^{\prime \prime }}`$ are close to the present upper limits. Thus, the production of heavy quarks from gluino decay may provide a good signal for detecting gluinos. Similarly, when we detect the production of $`(t+\overline{t})\stackrel{~}{g}`$, we have at least two heavy quarks for tagging.
IV. Conclusion
We have studied the processes $`PP(P\overline{P})(t+\overline{t})\stackrel{~}{g}X`$ in supersymmetric models with explicit $`R_p`$-violation. We have seen that it is possible to test the models at future Fermilab Tevatron Run-II and CERN LHC experiments, provided the $`\lambda ^{^{\prime \prime }}`$-type couplings are large enough. We suggest also to check the old data in Tevatron Run I and get stronger constraints on $`R_p`$ violation or exclude the narrow window of LSP heavy gluinos.
The process, with single production of top quark and gluino, can give signals both for SUSY and $`R_p`$ violation. Specifically, in a typical O-II model, the process will be an important one to check SUSY and $`R_p`$ violation. It is also shown that gluinos should be detected through their $`R_p`$-violating decay if $`R_p`$ violating parameters $`\lambda ^{^{\prime \prime }}`$ are close to the upper bounds.
Acknowledgement
Z.-H. Yu thanks the World Laboratory, Lausanne, for the scholarship.
Appendix
In this appendix we present the mass matrices which we need in the calculations.
A. The sfermion sector
In the mSUGRA model, the masses of squarks and sleptons follow from the GUT scale parameters $`m_0,m_{1/2},A_0,\mathrm{tan}\beta `$ and $`sign(\mu )`$ . The sfermion mass matrices are given by
$$\begin{array}{ccc}M_{\stackrel{~}{f}}^2=\left(\begin{array}{cc}m_{\stackrel{~}{f}_L}^2+m_f^2m_f(A_f\mu r_f)& \\ m_f(A_f\mu r_f)m_{\stackrel{~}{f}_R}^2+m_f^2& \end{array}\right),\hfill & & \end{array}$$
$`(A.1)`$
where $`m_{\stackrel{~}{f}_{L,R}}`$ are given in , $`m_f`$ are the masses of the partner fermions and $`r_{d,s,b}=r_{e,\mu ,\tau }=1/r_{u,c,t}=\mathrm{tan}\beta `$. From these matrices we can get the mixing angles $`\theta _f`$ and masses of the sfermions:
$$\begin{array}{ccc}\mathrm{sin}2\theta _f=\frac{2m_f(A_f\mu r_f)}{m_{\stackrel{~}{f}_1}^2m_{\stackrel{~}{f}_2}^2},\mathrm{cos}2\theta _f=\frac{m_{f_L}^2m_{f_R}^2}{m_{\stackrel{~}{f}_1}^2m_{\stackrel{~}{f}_2}^2},\hfill & & \end{array}$$
$`(A.2)`$
$$\begin{array}{ccc}m_{\stackrel{~}{f}_{1,2}}^2=m_f^2+\frac{1}{2}\left[m_{f_L}^2+m_{f_R}^2\sqrt{(m_{\stackrel{~}{f}_L}^2m_{\stackrel{~}{f}_R}^2)^2+4m_f^2(A_f\mu r_f)^2}\right].\hfill & & \end{array}$$
$`(A.3)`$
B. The chargino/neutralino sector
In order to calculate the decay of gluino, we need to consider the chargino and neutralino mass terms. The general chargino mass matrix is given as follows :
$$\begin{array}{ccc}M_C=\left(\begin{array}{cc}m_2\sqrt{2}m_w\mathrm{sin}\beta & \\ \sqrt{2}m_w\mathrm{cos}\beta \mu & \end{array}\right).\hfill & & \end{array}$$
$`(B.1)`$
It can be diagonalized by two real matrices U and V,
$$\begin{array}{ccc}U^{}M_CV^1,\hfill & & \end{array}$$
$`(B.2)`$
where $`U=𝒪_{}`$, $`V=𝒪_+`$ if $`detM_C0`$ and $`V=\sigma _3𝒪_+`$ if $`detM_C<0`$, with
$$\begin{array}{ccc}\sigma _3=\left(\begin{array}{cc}+\mathrm{1\; 0}& \\ 01& \end{array}\right),𝒪_\pm =\left(\begin{array}{cc}\mathrm{cos}\theta _\pm \mathrm{sin}\theta _\pm & \\ \mathrm{sin}\theta _\pm \mathrm{cos}\theta _\pm & \end{array}\right),\hfill & & \end{array}$$
$`(B.3)`$
and
$$\begin{array}{ccc}\mathrm{tan}2\theta _{}=\frac{2\sqrt{2}m_w(m_2\mathrm{cos}\beta +\mu \mathrm{sin}\beta )}{m_2^2\mu ^22m_w^2\mathrm{cos}\beta },\hfill & & \\ \mathrm{tan}2\theta _+=\frac{2\sqrt{2}m_w(m_2\mathrm{sin}\beta +\mu \mathrm{cos}\beta )}{m_2^2\mu ^2+2m_w^2\mathrm{cos}\beta }.\hfill & & \end{array}$$
$`(B.4)`$
Then the chargino masses are,
$$\begin{array}{ccc}m_{x_{1,2}^+}\hfill & =\hfill & \frac{1}{\sqrt{2}}[m_2^2+\mu ^2+2m_w^2\hfill \\ & & \{(m_2^2\mu ^2)^2+4m_w^4\mathrm{cos}^22\beta +4m_w^2(m_2^2+\mu ^2+2m_2\mu \mathrm{sin}2\beta )\}^{\frac{1}{2}}]^{\frac{1}{2}}\hfill \end{array}$$
$`(B.5)`$
In the case of the neutralinos, the mass matrix is given by :
$$\begin{array}{ccc}M_N=\left(\begin{array}{cc}m_10m_zs_w\mathrm{cos}\beta m_zs_w\mathrm{sin}\beta & \\ 0m_2m_zs_w\mathrm{cos}\beta m_zs_w\mathrm{sin}\beta & \\ m_zs_w\mathrm{sin}\beta m_zc_w\mathrm{sin}\beta 0\mu & \\ m_zs_w\mathrm{sin}\beta m_zc_w\mathrm{sin}\beta \mu 0& \end{array}\right).\hfill & & \end{array}$$
$`(B.6)`$
It can be diagonalized by a unitary matrix $`N`$ ($`4\times 4`$) as:
$$\begin{array}{ccc}N^TM_NN.\hfill & & \end{array}$$
$`(B.7)`$
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# 1 Introduction
## 1 Introduction
The problem of surface effects on conformation statistics of long flexible polymer chains has been widely studied both because of its merit as an interesting problem in statistical mechanics and because of its important role in many physical processes like colloidal stabilization, adhesion, or lubrication, etc . The statistical mechanics approach to this problem has successfully been applied, particularly in case of good solvent that contains only one linear polymer chain interacting with an impenetrable wall \[3-7\]. The essential physics is derived from a model of self-avoiding walk (SAW) on a semi-infinite lattice, with an energy contribution $`ϵ_a`$ for each step of the walk along the lattice boundary. This leads to an increased probability characterized by the Boltzmann factor $`\omega =\mathrm{exp}(ϵ_a/k_\beta T)`$ of making a step along the attractive wall, since $`ϵ_a<0`$, $`\omega >1`$ for any finite temperature $`T`$. At low temperatures due to the attraction between the polymer chain and the surface, the chain gets adsorbed on the surface while at high temperatures all polymer configurations have almost same weight and a non-adsorbed behaviour prevails. The transition between these two regions is marked by a critical adsorption temperature $`T_a`$, with a desorbed phase for $`T>T_a`$ and an adsorbed phase for $`T<T_a`$. The asymptotic behaviour of the average number $`M`$ of steps of the walk along the boundary can be summarized in the following way
$`M\{\begin{array}{cc}(T_aT)^{\frac{1}{\varphi }1}\hfill & T<T_a\hfill \\ N^\varphi \hfill & T=T_a\hfill \\ (TT_a)^1\hfill & T>T_a\hfill \end{array}`$ (1.4)
where $`N`$ denotes the average number of monomers and $`\varphi `$ is the crossover exponent.
When the polymer chain is in a poor solvent, it exhibits a phase diagram characterized by many different universality domains of critical behaviour. This is due to the competition between solvent - induced monomer-monomer attraction and the surface-monomer interaction. In this case the essential physics is derived using a model of self-attracting self-avoiding walk (SASAW) on a semi-infinite lattice .
Theoretical methods which have been used to study polymer adsorption include renormalization group , transfer matrix , Monte Carlo , exact enumeration and series expansion techniques . In case of two dimensions many exact results have been found through conformal field theory \[17-20\] and using conformal invariance prediction in conjunction with the Bethe ansatz solution of associated lattice models . Many exact results have also been found for the case of fractal lattices using real space renormalization group (RSRG) \[8, 23-25\] . In this article, we consider the adsorption and collapsed transition of a long flexible polymer chain interacting with another long flexible polymer chain adsorbed on a surface. The monomers of the adsorbed chain act as pinning (or interaction) sites for the chain of our interest. The situation is shown in Figure 1. A polymer chain shown by zigzag line lies on the surface. Its monomers shown by black circles act as a pinning sites for the other polymer chain which is shown by wiggle line. The monomers of this floating chain shown by open circles get attracted to the surface because of the pinning sites. When a monomer of the floating chain gets adsorbed on the pinning sites (shown by shade circle) there is a gain of energy $`ϵ_a`$. The model which we describe below also has two other interactions; (i) attraction between nearest neighbours formed by monomers of floating chain which are separated by relatively large distance along the chain and (ii) the interaction between nearest neighbours formed by a pinning site and a monomer of the floating chain occupying the neighbouring sites. These interactions are shown by $`ϵ_u`$ and $`ϵ_t`$, respectively, in the Figure 1. This situation is similar to the problem of adsorption of a long flexible polymer chain onto a (cell) membrane along which protein stick out in a spatially uncorrelated manner.
The paper is organized as follows. In section 2, we describe a model of two interacting chains in which one is adsorbed on a surface of a fractal container. The motivation of choosing a fractal space is that the model can be solved exactly apart from having its own practical applications. The results found by solving the model exactly are given in section 3. The paper ends with conclusions listed in section 4.
## 2 Model and its solution
We consider a fractal space mapped by truncated 4-simplex lattice. The basic geometrical unit of construction of this lattice is a tetrahedron with 4-corner vertices and bonds between every pair of vertices. Each vertex connected through a direct bond is termed as nearest neighbour. The fractal and spectral dimensions of the 4-simplex lattice are 2 and 1.5474, respectively. The tetrahedron of first and $`(r+1)`$th order are shown in Figure 2. The shaded region represents the surface. The surface is a truncated 3-simplex lattice with fractal and spectral dimensions, 1.5849 and 1.3652 respectively. The bulk critical behaviour including $`\theta `$-point and the phase diagram of surface interacting polymer chain using SASAW model has been studied recently .
In the model proposed here we represent the pinning sites by the monomers of an adsorbed polymer chain on a 3-simplex lattice. The configuration of the pinning sites can, therefore, be found by the statistics of single polymer chain on the 3-simplex lattice. For convenience we represent this configuration (of adsorbed polymer chain) by $`P_2`$. The floating polymer chain in a fractal container whose adsorption we want to study is represented by $`P_1`$. Therefore the problem is projected on to a model of two interacting crossed walks \[26-28\] in which one walk is confined to the surface and act as pinning sites for the other chain.
We assign the weight $`x_1(x_2)`$ to each step of the walk in the bulk (on the surface) and the weight $`\sqrt{x_1x_2}\omega `$ to each step taken on the pinning sites on the surface. In other words, when a monomer of chain $`P_1`$ visits a site on the surface occupied by chain $`P_2`$ a weight $`\sqrt{x_1x_2}\omega `$ is assigned. If a monomer visits a site on the surface not occupied by the chain $`P_2`$ (i.e a pure site) the weight assigned to it is $`x_1`$. The monomers of chain $`P_1`$ may attract each other. We denote the Boltzmann factor associated with this interaction by $`u`$, where $`u=\mathrm{exp}(ϵ_u/k_\beta T)`$ ($`ϵ_u<0`$ being the attractive energy associated with a pair of near-neighbour bonds on the lattice). In order to promote competition between the adsorbed and desorbed phases of chain $`P_1`$ it is desirable to introduce a parameter $`t=\mathrm{exp}(ϵ_t/k_\beta T)`$ in such a way that $`\sqrt{x_1x_2}t`$ is the weight of those steps that are performed on the lattice points which are the nearest neighbour but in a adjacent sites to the pinning sites. Here $`ϵ_t`$ is the interaction energy between a pair formed by a monomer of chain $`P_1`$ in the adjacent site and a pinning site on the surface.
The global generating function of this model can be written in the form
$`\mathrm{\Omega }(x_1,x_2,\omega ,u,t)`$ $`=`$ $`{\displaystyle \underset{allwalks}{}}x_1^{N_1}x_2^{N_2}\omega ^{N_s}t^{N_c}u^{N_m}`$
$`=`$ $`{\displaystyle \underset{N_1N_2N_sN_cN_m}{}}C(N_1N_2N_sN_cN_m)x_1^{N_1}x_2^{N_2}\omega ^{N_s}t^{N_c}u^{N_m}`$
where $`C(N_1,N_2,N_s,N_c,N_m)`$ represents the total number of configurations of all walks. Here $`N_1`$ is the total number of monomers in chain $`P_1`$, $`N_2`$ represents the total number of monomers in adsorbed polymer chain ($`P_2`$), $`N_s`$ denotes the number of monomers of polymer chain $`P_1`$ adsorbed on the pinning sites. $`N_c`$ and $`N_m`$ are number of monomers lying adjacent to the surface and forming the nearest neighbour with the pinning sites, and number of nearest neighbours in chain $`P_1`$, respectively.
The generating function for the 4-simplex lattice can be expressed in terms of finite number of restricted partition functions . We show in Figure 3 all the possible restricted partition functions which appear in this problem. This can be seen from Figure 4 in which we draw one of the possible configurations of walks on the $`r`$-th order of the lattice. The surface adsorbed chain $`P_2`$ is shown by zigzag line while the chain $`P_1`$ by wiggle line. These partition functions are defined recursively as weighted sum over all possible configurations for a given stage of the iterative construction of the 4-simplex lattice. The variables in these equations are just the partial generating functions corresponding to different polymer configurations for a given size of the fractal. Linearizing the recursion equations near the fixed points, the one reached by the system depending on the initial conditions, we can find the eigenvalues of the transformation matrix which give the characteristic exponents of the system.
The recursion relations for the restricted partition functions can be written as (see Figure 5)
$`A_{r+1}`$ $`=`$ $`A^2+2A^3+2A^4+4A^3B+6A^2B^2`$ (2.6)
$`B_{r+1}`$ $`=`$ $`A^4+4A^3B+22B^4`$ (2.7)
$`H_{r+1}`$ $`=`$ $`H^2+H^3`$ (2.8)
$`C_{r+1}`$ $`=`$ $`C^2+C^3+AD^2(H+2C+2F+2G+2+2A)+`$
$`4A^2G(C+F)+2AG^2(A+C)+4ACFG+HD^2`$
$`D_{r+1}`$ $`=`$ $`ABD(4G+2F+2E)+A^2D(G+F+E+C+H)+`$
$`AD(C+H+C^2+H^2)+BDE(2F+4G)+`$
$`ACD(F+G+H)+AHDE+D^3(2B+A)`$
$`E_{r+1}`$ $`=`$ $`AD^2(4B+2G+2E+2H+2A)+BD^2(2F+4G)+`$
$`AH^2(A+E)+6B^2E^2+2BE^3+2A^2HE`$
$`F_{r+1}`$ $`=`$ $`AD^2(2B+E+2G)+B^2(8G^2+6F^2+`$
$`B(8G^3+2F^3)+8BFG(B+G)+ACF(2A+C)+`$
$`BED^2+4BGF^2+A^2C^2`$
$`G_{r+1}`$ $`=`$ $`AD^2(2B+A+G+F+C)+AG(2AC+C^2)+`$
$`12B^2FG+6BGF^2+10BFG^2+10B^2G^2+`$
$`6BG^3+BD^2E`$
A notational simplification in which the index $`r`$ is dropped from the right hand side of the recursion relations is adopted here. It may be emphasized here that the recursion relations written above are exact for the model defined above.
Eq.(2.5) which represents the recursion relation for the adsorbed chain $`P_2`$ is independent of configuration of chain $`P_1`$. The effect of $`P_2`$ on chain $`P_1`$ is taken through $`C`$, $`D`$, $`E`$, $`F`$ and $`G`$. Since all interactions involved in the problem are restricted to bonds within a first order unit of the fractal lattice, $`\omega `$, $`t`$, and $`u`$ do not appear explicitly in the recursion relations. They appear only in initial values given below.
$`A_1`$ $`=`$ $`x_1^2+2x_1^3u+2x_1^4u^3`$ (2.14)
$`B_1`$ $`=`$ $`x_1^4u^4`$ (2.15)
$`C_1`$ $`=`$ $`x_1^2x_2^3\omega ^2t^2+x_1^3x_2^3(\omega ^6u+t^2\omega ^2u)+2x_1^4x_2^3\omega ^4t^2u^3+x_1^2x_2^2\omega ^4+`$
$`2x_1^3x_2^2t^2\omega ^2u+2x_1^4x_2^2t^2\omega ^2u^3`$
$`D_1`$ $`=`$ $`x_1^2x_2^3t\omega +x_1^3x_2^3(t\omega u+t\omega ^3u)+x_1^4x_2^3(u^3\omega ^3t^3+u^3\omega ^5t)+`$
$`x_1^2x_2^2t\omega +x_1^3x_2^2(\omega ^3tu+\omega tu)+x_1^4x_2^2(\omega ^3tu^3+\omega tu^3)`$
$`E_1`$ $`=`$ $`x_1^2x_2^3+2x_1^3x_2^3\omega ^2t^2u+2x_1^4x_2^3t^2\omega ^2u^3+x_1^2x_2^2+2x_1^3x_2^2u+`$
$`2x_1^4x_2^2u^3t^2\omega ^2`$
$`F_1`$ $`=`$ $`x_1^4x_2^2\omega ^4u^4+x_1^4x_2^3\omega ^2t^2u^4`$ (2.19)
$`G_1`$ $`=`$ $`x_1^4x_2^2t^2\omega ^2u^4+x_1^4x_2^3\omega ^4t^2u^4`$ (2.20)
Here index 1 and 2 on the right hand side of Eq.(2.11)-(2.17) correspond to chain $`P_1`$ and pinning sites (chain $`P_2`$) respectively. The fixed points corresponding to different configurations of polymers in the asymptotic limit are found by solving Eq.(2.3) and (2.4) for polymer chain $`P_1`$ and Eq.(2.5) for chain $`P_2`$ respectively. A complete phase diagram obtained from Eqs.(2.3) - (2.4) and from Eq.(2.5) are given elsewhere .
The state of polymer chain $`P_1`$ depends on the quality of the solvent and on the temperature and can therefore be in any of three states; swollen, compact globule and at $`\theta `$-point described in the asymptotic limit by the fixed points $`(A^{},B^{})`$ = (0.4294, 0.0498), (0.0,$`22^{1/3})`$ and (1/3,1/3) respectively. The fixed point corresponding to the swollen state is reached for all values of $`u<u_\theta `$ at $`x_1>x_\theta `$. The end to end distance for a chain of $`N_1`$ monomers of $`P_1`$ in this state varies as $`N_1^{\nu _1}`$ with $`\nu _1`$ =0.7294 and connectivity constant $`\mu =1/x`$ is found to be 1.5474. The fixed point corresponding to the compact globule state is reached for all values of $`u>u_\theta `$ at $`x_1(u)<x_\theta `$. At $`u_\theta `$ = 3.31607.. and $`x_\theta `$ = 0.22913.. the system is found to be at its tricritical point or $`\theta `$-point. The fixed point $`H^{}`$ = 0.61803.. is found by solving Eq.(2.5). It corresponds to a pattern with fractal dimension equal to 1.266.
In a system of polymer chain interacting with a surface adsorbed chain, we have three different combinations of the individual state. Using the fixed points of $`(A^{},B^{})`$ corresponding to three different states of polymer chain $`P_1`$ and $`H^{}=0.61803`$ (for surface adsorbed chain $`P_2`$), we solve the coupled non-linear equations \[Eqs. (2.6) - (2.10)\].
## 3 Results
For all values of $`\omega <\omega _c(u,t)`$, the polymer chain $`P_1`$ lies in the bulk. The critical behaviour of polymer chain $`P_1`$ does not get affected by the presence of the surface or pinning sites. In the phase diagram shown in Figure 6, $`\omega `$ is plotted as a function of $`u`$ for three values of $`t=0`$, 0.5 and 1.0. The $`\theta `$-line which separates the bulk swollen and collapsed phases and terminates at the surface adsorption line $`\omega =\omega _c(u,t)`$ is found at $`u=u_c=3.316074`$. The $`\theta `$-line in Figure 6 is shown by dashed line. $`\omega =\omega _c(u,t)`$ separates the bulk from the adsorbed phase. Below this line the desorbed phase does not get affected by the presence of the surface attraction, except for the $`\theta `$-point turning into a $`\theta `$-line and the critical lines corresponding to swollen and collapsed phases into respective regions.
When $`\omega >\omega _c(u,t)`$ the polymer gets adsorbed. Depending on the value of $`t`$ the chain may lie on the pinning sites or avoids them. For $`t1`$ the polymer chain may get adsorbed on the pinning sites acquiring the configuration as that of the pinning sites.
When $`\omega =\omega _c(u,t)`$ the chain $`P_1`$ is on the adsorption special line. The different regions of this line characterize different multicritical behavior as a function of $`u`$ and $`t`$. Note that the parameter $`u`$ and $`t`$ measure, respectively, the strengths of nearest neighbour interaction and the repulsive strength of a monomer on the adjacent layer and forming a nearest neighbour with a pinning site on the surface.
(I) The fixed point $`A^{}`$, $`B^{}`$,$`C^{}`$, $`D^{}`$, $`E^{}`$, $`F^{}`$, $`G^{}`$ ) = (0.4294, 0.04998, 0.2654, 0.2654, 0.2654, 0.0391, 0.0391) is reached for all values of $`u<u_c`$ and $`\omega =\omega _c`$ at $`t>0`$. Linearization around this fixed point gives one eigenvalue (other than $`\lambda _b=2.7965`$) $`\lambda _\varphi =1.7914`$ greater than one. From the eigenvalues $`\lambda _b`$ and $`\lambda _\varphi `$ one gets the crossover exponent
$$\varphi =\frac{\mathrm{ln}\lambda _\varphi }{\mathrm{ln}\lambda _b}=0.56690.57$$
This fixed point corresponds to special line $`\omega =\omega _c(u,t)`$ for $`u<u_c`$ and $`t>0`$. The polymer chain at the adsorption line (tricritical line) fluctuates among configurations corresponding to $`A`$ (bulk), $`C`$, $`D`$, $`E`$ (surface) \[see Figure 3\].
For $`t=0`$ the fixed point ($`A^{}`$, $`B^{}`$, $`C^{}`$, $`D^{}`$, $`E^{}`$, $`F^{}`$, $`G^{}`$ ) =(0.4294, 0.04998, 0.0, 0.0, 0.1164, 0.0, 0.0) is reached for all values of $`u<u_c`$. Linearization about this fixed point does not yield any eigenvalue (except $`\lambda _b`$) greater than one. Therefore there is no crossover from bulk to surface.
(II) When $`u>u_\theta `$ (= 3.316074) and $`x<x_\theta `$ (= 0.229157…) the polymer chain is in the collapsed state in the bulk. At $`t=0`$, we choose $`u=3.3333..`$ and $`x=0.2282..`$ and solve Eqs.(2.3) - (2.9). This leads to a fixed point ($`A^{},B^{},C^{},D^{},E^{},F^{},G^{}`$) = (0.0, 0.3568, 0.0, 0.0, 0.7652, 0.0, 0.0) at $`\omega =\omega _c=1.5751\mathrm{}.`$. This corresponds to configuration in which compact globule formed in the bulk container gets stuck to the surface in such a way that monomers avoid touching the sites of the surface occupied by the chain $`P_2`$. Linearization around this fixed point gives in addition to $`\lambda _b`$ an eigenvalue greater than one i.e $`\lambda _\varphi =2.4233`$. With this eigenvalue the crossover exponent is found to be
$$\varphi =\frac{\mathrm{ln}\lambda _\varphi }{\mathrm{ln}\lambda _b}0.64.$$
(III) When $`t>0`$ and $`\omega =\omega _c`$ (for $`u>u_\theta `$ and $`x<x_\theta `$), we get a fixed point ($`A^{},B^{},C^{},D^{},E^{},F^{},G^{}`$) = (0.0, 0.3568, 0.0, 0.0, 0.0, 0.2206,0.2206). Linearization around this fixed point gives $`\lambda _\varphi =2.3073`$. The crossover exponent in this case is $`\varphi 0.60`$. In this case the surface configurations attained by the polymer chain are those which correspond to the configurations $`F`$ and $`G`$ shown in Figure 3. It means that the globule gets attached to the site occupied by the chain $`P_2`$ on the surface.
(IV) When $`u=u_\theta =3.316074\mathrm{}..`$ and $`x=x_\theta =0.229137\mathrm{}`$ the polymer chain (in bulk) is at $`\theta `$-point. For this we have three distinct fixed points depending on the values of $`t`$. We now discuss these three fixed points.
(A) For $`t=0`$ and $`\omega =\omega _c=1.5772`$ the fixed point achieved by the system is ($`A^{}`$, $`B^{}`$, $`C^{}`$, $`D^{}`$, $`E^{}`$, $`F^{}`$, $`G^{}`$) = (1/3, 1/3, 0.0, 0.0, 0.613, 0.8228.., 0.0). Linearized equations around this fixed point give three eigenvalues greater than one. These values are
$$\lambda _1=2.4511,$$
and
$$\lambda _{b1}=3.7037,\lambda _{b2}=2.2222$$
Note that the last two eigenvalues are the same as those found for the bulk $`\theta `$-point. The crossover exponent is
$$\varphi =\frac{\mathrm{ln}\lambda _1}{\mathrm{ln}\lambda _{b1}}0.65.$$
This is a tetracritical point. The adsorbed phase corresponding to configurations which forms a layer on the surface occupied by the pinning sites.
(B) For $`0<t<1`$ we have fixed point ($`A^{}`$, $`B^{}`$, $`C^{}`$, $`D^{}`$, $`E^{}`$, $`F^{}`$, $`G^{}`$) = (1/3, 1/3, 0.0510.., 0.0, 0.613, 0.2364.., 0.2364). This point has been found to have three eigenvalues greater than one. These values are
$$\lambda _1=2.3311,$$
and the other two are $`\lambda _{b1}`$ and $`\lambda _{b2}`$ given above. The crossover exponent, $`\varphi `$ found in this case is equal to 0.61. This tetracritical point differs from the previous one in the sense that the adsorbed phase has configuration which is combination of both $`F`$ and $`G`$ as shown in Figure 3 whereas at $`t=0`$ the adsorbed phase has the configuration corresponding to $`E`$.
(C) At $`t=1`$ and $`\omega =1`$ we find the symmetrical fixed point ($`A^{}`$, $`B^{}`$, $`C^{}`$, $`D^{}`$, $`E^{}`$, $`F^{}`$, $`G^{}`$) = (1/3, 1/3, 0.2061, 0.2061, 0.2061, 0.2061, 0.2061). This point has four eigenvalues greater than one. Apart from the two known eigenvalues ($`\lambda _{b1}`$, and $`\lambda _{b2}`$) the two additional eigenvalues are
$$\lambda _1=2.1269\mathrm{and}\lambda _2=1.1947$$
This is a pentacritical point. Note that for $`t>1`$, corresponds to $`\omega <1`$. Therefore the tetracritical line as a function of $`t`$ is symmetrical about the point $`t=1`$.
## 4 Conclusions
In this paper we studied the critical behaviour of polymer chain interacting with a surface adsorbed chain. It is shown that this model differs from the usual polymer adsorption and also from the problem of two interacting chains studied in past \[15,23-28\] . We find a very rich $`\omega u`$ phase diagram plotted in Figure 6. The adsorbed polymer chain representing the pinning sites always remain in swollen state with radius of gyration exponent equal to that of truncated 3-simplex lattice. We therefore have pinning sites forming pattern with fractal dimension 1.266. The bulk desorbed phase has two regions: the region of swollen state separated from the collapsed globule state (by a tricritical $`\theta `$-line). The $`\theta `$-line is at $`u=u_\theta `$ = 3.316074.. and runs parallel to the $`\omega `$ axis i.e. remains unaltered due to the surface interaction. The point where it meets the line $`\omega ^{}(t,u)`$ is a multicritical point. These multi-critical points are characterized by three different fixed points depending on the value of $`t`$. When the value of $`t=1`$ the value of $`\omega `$ is found to be 1. The $`\omega `$ line runs parallel to the $`u`$-axis and meets at the $`\theta `$-point. This point has four eigenvalues greater than one and corresponds to the pentacritical point. This is a point at which two tetracritical lines corresponding to $`0<t1`$ and $`t>1`$ meet. When $`t>1`$, $`\omega `$ has to be less than $`1`$.
When $`t=0`$, the value of surface interaction increases with $`u`$ and meet at $`\theta `$-line. The Figure 6 gives the impression of the existence of reentrant adsorbed phase as $`u`$ is increased. One should, however, remember that these figures are merely a projection on the $`\omega u`$ plane of three dimensional figures in which the third dimension is given by $`x`$.
When the value of $`t`$ lies in between 0 and 1 the slope increases, but we do not find any “frustrated phase” as observed in usual situation of surface adsorption. The behaviour of special adsorption line described above can be understood from contributions of different coexisting polymer configurations (see Figure 3) to the bulk and surface free energies. When both adsorbed phase (as by definition) and bulk phase are in swollen state, the adsorption line has same nature in $`\omega u`$ plane for all values of $`t`$, although the slope of line decreases as $`t`$ increases and at $`t=1`$ the slope becomes zero and line runs parallel to $`u`$ axis. This is due to the very fact that at $`t=1`$ and $`\omega =1`$ the surface behaves as a part of bulk and distribution of monomers are isotropic.
Acknowledgement
We thank the Department of Science and Technology of Govt. of India for financial assistance. We also thank D. Dhar for many helpful discussions.
FIGURE CAPTION
Figure 1 Schematic representation of all possible interactions appearing in Generating Function defined by Equation 2.2.
Figure 2 Graphical representation of a truncated 4-simplex lattice of first order and ($`r+1`$)th order. The shaded regions represent the surface.
Figure 3 Diagrammatic representation of all the restricted partition functions which appear in the generating function defined by Eq.(2.2). Polymer chain $`P_1`$ in the bulk is represented by wiggle line while adsorbed polymer chain $`P_2`$ is shown by zigzag line.
Figure 4 Diagrammatic representation of one of possible configurations of a polymer interacting with surface confined polymer chain on third order of $`4`$-simplex lattice. All the possible partition functions contributing to the generating function are shown.
Figure 5 Diagrams representing the recursion relation for the restricted partition functions given by Eqs. 2.3-2.9. Some of the possible configurations contributing to the recursion relations are shown.
Figure 6 Special adsorption lines are shown in $`\omega `$ -$`u`$ plane for $`t`$ = 0, 1 and 0.5.
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# TUM-HEP-379/00 CERN-TH/2000-190 Universal Unitarity Triangle and Physics Beyond the Standard Model
## 1 Introduction
One of the important goals of particle physics is the determination of the Cabibbo–Kobayashi–Maskawa (CKM) matrix . In addition to the leading tree level K and B decays, flavour changing neutral current processes generated at the one loop level in the Standard Model (SM) and sensitive to the top quark couplings $`V_{td(s)}`$ play a crucial role in this determination. This program is not only complicated by the presence of hadronic uncertainties but also by the possible existence of new physics that contributes to various quantities through diagrams involving new particles. These new contributions depend on unknown parameters, like the masses and couplings of new particles, that pollute the extraction of the CKM parameters.
We would like to point out that in a certain class of extensions of the SM it is possible to construct measurable quantities that depend on the CKM parameters but are not polluted by new physics contributions. This means that these quantities allow a direct determination of the “true” values of the CKM parameters which are common to the SM and this particular class of its extensions. Correspondingly there exists a universal unitarity triangle common to all these models. Interestingly the quantities required to construct the universal unitarity triangle are essentially free from hadronic uncertainties.
In order to explain our point we use the Wolfenstein parameterization of the CKM matrix and its generalization to include higher order terms in $`\lambda `$ .
Let us recall first that the four Wolfenstein parameters $`\lambda `$, $`A`$, $`\varrho `$ and $`\eta `$ can be determined in the standard manner as follows:
Step 1:
The parameters $`\lambda `$ and $`A`$ are determined from semileptonic K and B decays sensitive to the elements $`|V_{us}|`$ and $`|V_{cb}|`$ respectively:
$$\lambda =|V_{us}|=0.22,A=\frac{|V_{cb}|}{\lambda ^2}=0.826\pm 0.041.$$
(1)
As the decays in question are tree level decays with large branching ratios this determination is to an excellent approximation independent of any possible physics beyond the SM.
Step 2:
The parameters $`\varrho `$ and $`\eta `$ are determined by constructing with the help of various decays the unitarity triangle of fig. 1, where
$$\overline{\varrho }=\varrho (1\frac{\lambda ^2}{2}),\overline{\eta }=\eta (1\frac{\lambda ^2}{2})$$
(2)
describe the apex of this triangle. The lengths CB, CA and BA are equal respectively to
$$1,R_b\sqrt{\overline{\varrho }^2+\overline{\eta }^2}=(1\frac{\lambda ^2}{2})\frac{1}{\lambda }\left|\frac{V_{ub}}{V_{cb}}\right|,R_t\sqrt{(1\overline{\varrho })^2+\overline{\eta }^2}=\frac{1}{\lambda }\left|\frac{V_{td}}{V_{cb}}\right|.$$
(3)
The standard construction of this triangle involves the ratio $`|V_{ub}/V_{cb}|`$ extracted from inclusive and exclusive tree level B decays and flavour changing neutral current processes such as $`B_d^0\overline{B}_d^0`$ mixing (the mass difference $`(\mathrm{\Delta }M)_d`$) and indirect CP violation in $`K_L`$ decays (the parameter $`\epsilon `$), both sensitive to the CKM element $`V_{td}`$. There is also a constraint coming from the lower bound on the mass difference $`(\mathrm{\Delta }M)_s`$ describing $`B_s^0\overline{B}_s^0`$ mixing. In particular in the case of $`B_{d,s}^0\overline{B}_{d,s}^0`$ mixings the following formulae for $`(\mathrm{\Delta }M)_{d,s}`$ resulting from box diagrams are used:
$$(\mathrm{\Delta }M)_{d,s}=\frac{G_F^2}{6\pi ^2}\eta _Bm_{B_{d,s}}(\widehat{B}_{B_{d,s}}F_{B_{d,s}}^2)M_W^2F_{tt}|V_{t(d,s)}|^2$$
(4)
Here $`F_{tt}`$ is a function of $`m_\mathrm{t}`$ and $`M_\mathrm{W}`$ resulting from box diagrams with top quark exchanges, $`\widehat{B}_B`$ is a non-perturbative parameter, $`F_B`$ is the B meson decay constant and $`\eta _B`$ the short distance QCD factor common to $`(\mathrm{\Delta }M)_d`$ and $`(\mathrm{\Delta }M)_s`$.
Similarly, the experimental value for $`\epsilon `$ combined with the theoretical calculation of box diagrams describing $`K^0\overline{K}^0`$ mixing gives the constraint for $`(\overline{\varrho },\overline{\eta })`$ in the form of the following hyperbola :
$$\overline{\eta }\left[(1\overline{\varrho })A^2\eta _2F_{tt}+P_c(\epsilon )\right]A^2\widehat{B}_K=0.226.$$
(5)
Here $`\widehat{B}_K`$ is a non-perturbative parameter analogous to $`\widehat{B}_{B_{d,s}}`$, $`\eta _2`$ is a short distance QCD correction , $`F_{tt}`$ is the function present also in (4) and $`P_c(\epsilon )=0.31\pm 0.05`$ summarizes charm–charm and charm–top contributions.
Combining the two steps above one can determine the range of values of $`(\overline{\varrho },\overline{\eta })`$ consistent with all present data. Analyses of this type can be found in . In the future this procedure can be generalized to include CP asymmetries in B decays sensitive to the angles of the unitarity triangle and various branching ratios for K and B decays sensitive to the sides and the height of this triangle . If the SM is the correct theory all these measurements should result in a unique value of $`(\overline{\varrho },\overline{\eta })`$.
This procedure of testing the SM can be applied to its extensions as well. Step 1 remains unchanged as this determination, based on tree level decays, is insensitive to physics beyond the SM. On the other hand Step 2 can be affected by new physics due to:
* New contributions to box diagrams modifying the function $`F_{tt}`$ and to the analogous functions describing various penguin diagrams contributing to rare K and B decays. This introduces new parameters into the box function $`F_{tt}`$ and the penguin functions that in the SM depend only on $`m_\mathrm{t}`$ and $`M_\mathrm{W}`$.
* New contributions to box and penguin diagrams that are not proportional to the same combination of CKM matrix elements as the SM top contribution (for example, new contributions to $`P_c`$ in eq. (5) or new contributions to $`(\mathrm{\Delta }M)_{d,s}`$ proportional to $`|V_{cd(s)}^{}V_{cb}|^2`$).
* New complex phases beyond the one present in the CKM matrix.
* New local operators contributing to the relevant amplitudes beyond those present in the SM. This would introduce additional non-perturbative factors $`B_i`$ and new box and penguin functions.
It is evident from (4) and (5) that any modification of the function $`F_{tt}`$ will change the values of the extracted $`(\overline{\varrho },\overline{\eta })`$. A recent analysis of this type within the MSSM can be found in . Similar comments apply to the extraction of $`(\overline{\varrho },\overline{\eta })`$ from various branching ratios for rare K and B decays. Moreover if new phases are present in the extensions of the SM, CP violating asymmetries will generally measure different quantities than $`\alpha `$, $`\beta `$ and $`\gamma `$ in fig. 1. For instance the CP asymmetry in $`B\psi K_S`$ will no longer measure $`\beta `$ but $`\beta +\theta _{NP}`$ where $`\theta _{NP}`$ is a new phase. Strategies for dealing with such situations have been developed. See for instance and references therein.
The presence of new physics and of new phases will be signaled by inconsistencies in the $`(\overline{\varrho },\overline{\eta })`$ plane. In order to sort out which type of new physics is responsible for deviations from the SM expectations one has to study many loop induced decays and many CP asymmetries. Some ideas in this direction can be found in .
While in principle a global fit of all experimental data can be used to test the SM and its extensions it is desirable to develop strategies which allow to make these tests in a transparent manner.
Here we will concentrate on models like the SM, the Two Higgs Doublet Models (TDHM) I and II and the MSSM with minimal flavour violation, that do not have any new operators beyond those present in the SM and in which all flavour changing transitions are governed by the CKM matrix with no new phases beyond the CKM phase. Furthermore, in these models the only sizable new contributions are proportional to the same CKM parameters as the SM top contributions. That is, only the values of the functions describing top-mediated contributions to box and penguin diagrams are modified.
We would like to point out that the models in this class share a useful property. Namely, the CKM parameters in these models extracted from a particular set of data are independent of the contributing loop functions like $`F_{tt}`$, they are universal in this class of models. Correspondingly there exists a universal unitarity triangle. The determination of this universal unitarity triangle and of the corresponding CKM parameters has four virtues:
* The CKM matrix can be determined without the knowledge of new unknown parameters present in these particular extensions of the SM.
* Because the extracted CKM matrix is also valid in these models, the dependence of various quantities on the new parameters becomes more transparent. In short: the determination of the CKM matrix and of the new parameters can be separated from each other, as opposed to the present strategies discussed in step 2 above.
* The comparison of the predictions for a given observable in the SM and in this kind of extensions can then be done keeping the CKM parameters fixed.
* The extraction of the universal CKM parameters is essentially free from hadronic uncertainties.
In what follows we will list the set of quantities which allow a determination of the universal unitarity triangle. Subsequently we will indicate how the models in this class can be distinguished from each other and from more complicated models which bring in new complex phases and new operators.
## 2 Determination of $`R_t`$
In order to illustrate our point let us consider (4). Using this formula one finds
$$\frac{|V_{td}|}{|V_{ts}|}=\xi \sqrt{\frac{m_{B_s}}{m_{B_d}}}\sqrt{\frac{(\mathrm{\Delta }M)_d}{(\mathrm{\Delta }M)_s}}\kappa ,\xi =\frac{F_{B_s}\sqrt{\widehat{B}_{B_s}}}{F_{B_d}\sqrt{\widehat{B}_{B_d}}}.$$
(6)
This ratio depends only on measurable quantities $`(\mathrm{\Delta }M)_{d,s}`$, $`m_{B_{d,s}}`$ and the non-perturbative parameter $`\xi `$. Now to an excellent accuracy
$$|V_{td}|=|V_{cb}|\lambda R_t,|V_{ts}|=|V_{cb}|(1\frac{1}{2}\lambda ^2+\overline{\varrho }\lambda ^2)$$
(7)
with $`\overline{\varrho }`$ defined in (2). We note next that through the unitarity of the CKM matrix, the present experimental upper bound on $`(\mathrm{\Delta }M)_d/(\mathrm{\Delta }M)_s`$ and the value of $`|V_{ub}/V_{cb}|`$ one has in all these models $`0\overline{\varrho }\mathrm{\Gamma }<\mathrm{\hspace{0.33em}0.5}`$, where $`\xi =1.16\pm 0.07`$ has been used. Consequently $`|V_{ts}|`$ deviates from $`|V_{cb}|`$ by at most $`2.5\%`$. This means that to a very good accuracy $`R_t`$ is given by
$$R_t=\frac{\kappa }{\lambda }$$
(8)
independently of new parameters characteristic for a given model and of $`m_\mathrm{t}`$. If necessary the $`𝒪(\lambda ^2)`$ corrections in (7) can be incorporated in (8). This will be only required when the error on $`\xi `$ will be decreased below $`2\%`$, which is clearly a very difficult task.
While the ratio $`(\mathrm{\Delta }M)_d/(\mathrm{\Delta }M)_s`$ will be the first one to serve our purposes, there are at least two other quantities which allow a clean measurement of $`R_t`$ within the class of extensions of the SM considered. These are the ratios
$$\frac{Br(BX_d\nu \overline{\nu })}{Br(BX_s\nu \overline{\nu })}=\left|\frac{V_{td}}{V_{ts}}\right|^2$$
(9)
$$\frac{Br(B_d\mu ^+\mu ^{})}{Br(B_s\mu ^+\mu ^{})}=\frac{\tau _{B_d}}{\tau _{B_s}}\frac{m_{B_d}}{m_{B_s}}\frac{F_{B_d}^2}{F_{B_s}^2}\left|\frac{V_{td}}{V_{ts}}\right|^2$$
(10)
which similarly to $`(\mathrm{\Delta }M)_d/(\mathrm{\Delta }M)_s`$ measure
$$\left|\frac{V_{td}}{V_{ts}}\right|^2=\lambda ^2\frac{(1\overline{\varrho })^2+\overline{\eta }^2}{1+\lambda ^2(2\overline{\varrho }1)}\lambda ^2R_t^2.$$
(11)
Out of these three ratios the cleanest is (9), which is essentially free of hadronic uncertainties . Next comes (10), involving $`SU(3)`$ breaking effects in the ratio of $`B`$ meson decay constants. Finally, $`SU(3)`$ breaking in the ratio of bag parameters $`\widehat{B}_{B_d}/\widehat{B}_{B_s}`$ enters in addition in (6). These $`SU(3)`$ breaking effects should eventually be calculable with reasonable precision from lattice QCD.
It should be remarked that the branching ratio for the rare decay $`K^+\pi ^+\nu \overline{\nu }`$ is known to provide a clean measurement of $`V_{td}`$ and consequently of $`R_t`$ . However, this branching ratio alone cannot serve our purposes because it is sensitive to new physics contributions.
## 3 Determination of $`\beta `$ and $`\gamma `$
In order to complete the determination of $`\overline{\varrho }`$ and $`\overline{\eta }`$ in the universal unitarity triangle one can use $`\mathrm{sin}2\beta `$ extracted either from the CP asymmetry in $`B_d\psi K_S`$ or from $`K\pi \nu \overline{\nu }`$ decays . In the first case one has to measure the time dependent asymmetry
$$a_{CP}(t,\psi K_S)=\mathrm{sin}(2\beta )\mathrm{sin}((\mathrm{\Delta }M)_dt)$$
(12)
that allows a measurement of the angle $`\beta `$ without any hadronic uncertainties. In the second case the measurements of $`Br(K^+\pi ^+\nu \overline{\nu })`$ and $`Br(K_\mathrm{L}\pi ^0\nu \overline{\nu })`$ are required. Then :
$$\mathrm{sin}2\beta =\frac{2r_s}{1+r_s^2}$$
(13)
with
$$r_s(B_1,B_2)=\sqrt{\sigma }\frac{\sqrt{\sigma (B_1B_2)}P_c(\nu \overline{\nu })}{\sqrt{B_2}}.$$
(14)
Here $`\sigma =1/(1\lambda ^2/2)^2`$ and $`B_{1,2}`$ stand for the “reduced” branching ratios
$$B_1=\frac{Br(K^+\pi ^+\nu \overline{\nu })}{4.1110^{11}}B_2=\frac{Br(K_\mathrm{L}\pi ^0\nu \overline{\nu })}{1.8010^{10}}.$$
(15)
It should be stressed that $`\mathrm{sin}2\beta `$ determined in this manner depends only on two measurable branching ratios and on $`P_c(\nu \overline{\nu })=0.42\pm 0.06`$ which is completely calculable in perturbation theory . Moreover, hadronic uncertainties in these decays have been found to be negligibly small . As analyzed in , a measurement of both branching ratios within $`\pm 10\%`$ will allow the determination of $`\mathrm{sin}2\beta `$ within $`\pm 0.05`$.
Both extractions of $`\mathrm{sin}2\beta `$ are to an excellent accuracy independent of the new parameters characteristic for a given model. In particular $`P_c(\nu \overline{\nu })`$ being proportional to $`V_{cs}^{}V_{cd}`$ receives only negligible new contributions in the class of models considered .
Concerning the determination of the angle $`\gamma `$, the two theoretically cleanest methods are: i) the full time dependent analysis of $`B_sD_s^+K^{}`$ and $`\overline{B}_sD_s^{}K^+`$ and ii) the well known triangle construction due to Gronau and Wyler which uses six decay rates $`B^\pm D_{CP}^0K^\pm `$, $`B^+D^0K^+,\overline{D}^0K^+`$ and $`B^{}D^0K^{},\overline{D}^0K^{}`$. Variants of the latter method which could be more promising experimentally have been proposed in . Both methods involve only tree diagrams and are unaffected by new physics contributions in the class of models considered. It appears that these methods will give useful results at later stages of CP-B investigations. In particular the first method will be feasible only at LHC-B. Clearly any other method for the determination of $`\gamma `$ in which new physics of the type considered here can be eliminated could also be used. For a recent review of $`\gamma `$ determinations we refer to and references therein.
## 4 Determination of the Universal Unitarity Triangle
Once $`R_t`$ and $`\mathrm{sin}2\beta `$ have been determined as discussed above, $`\overline{\varrho }`$ and $`\overline{\eta }`$ can be found through
$$\overline{\eta }=a\frac{R_t}{\sqrt{2}}\sqrt{\mathrm{sin}2\beta r_b(\mathrm{sin}2\beta )},\overline{\varrho }=1\overline{\eta }r_b(\mathrm{sin}2\beta )$$
(16)
where
$$r_b(z)=(1+b\sqrt{1z^2})/z,a,b=\pm .$$
(17)
Thus for given values of $`(R_t,\mathrm{sin}2\beta )`$ there are four solutions for $`(\overline{\varrho },\overline{\eta })`$ corresponding to $`(a,b)=(+,+),(+,),(,+),(,)`$. As described in three of these solutions can be eliminated by using further information, for instance coming from $`|V_{ub}/V_{cb}|`$ and $`\epsilon `$, so that eventually the solution corresponding to $`(a,b)=(+,+)`$ is singled out
$$\overline{\eta }=\frac{R_t}{\sqrt{2}}\sqrt{\mathrm{sin}2\beta r_{}(\mathrm{sin}2\beta )},\overline{\varrho }=1\overline{\eta }r_+(\mathrm{sin}2\beta ).$$
(18)
We will illustrate this with an example below.
On the other hand $`\overline{\varrho }`$ and $`\overline{\eta }`$ following from $`R_t`$ and $`\gamma `$ are simply given by
$$\overline{\eta }=R_b\mathrm{sin}\gamma \overline{\varrho }=R_b\mathrm{cos}\gamma $$
(19)
with
$$R_b=\mathrm{cos}\gamma \pm \sqrt{R_t^2\mathrm{sin}^2\gamma }.$$
(20)
Comparing the resulting $`R_b`$ with the one extracted from $`|V_{ub}/V_{cb}|`$ (see (3)) one of the two solutions can be eliminated.
As an alternative to $`\mathrm{sin}2\beta `$ or $`\gamma `$ one could use the measurement of $`\sqrt{\overline{\varrho }^2+\overline{\eta }^2}`$ by means of $`|V_{ub}/V_{cb}|`$ but this strategy suffers from hadronic uncertainties in the extraction of $`|V_{ub}/V_{cb}|`$. Similarly using $`|V_{ub}/V_{cb}|`$ and $`\gamma `$ one can construct the the universal unitarity triangle by means of (19).
We observe that all these different methods determine the “true” values of $`\overline{\eta }`$ and $`\overline{\varrho }`$ independently of new physics contributions in the class of models considered. Since $`\lambda `$ and $`|V_{cb}|=A\lambda ^2`$ are determined from tree level K and B decays they are insensitive to new physics as well. Thus the full CKM matrix can be determined in this manner. The corresponding universal unitarity triangle common to all the models considered can be found directly from formulae like (16), (18) and (19).
As an example let us take $`(\mathrm{\Delta }M)_d=0.471/ps`$, $`(\mathrm{\Delta }M)_s=16.0/ps`$ and $`\xi =1.16`$. This gives $`R_t=0.92`$. Taking $`\mathrm{sin}2\beta =0.70`$ one finds then by means of (16) the four solutions for the universal unitarity triangle given in table 1. As from the data on $`|V_{ub}/V_{cb}|`$ we have $`R_b\mathrm{\Gamma }<\mathrm{\hspace{0.33em}0.5}`$ only the first solution is allowed.
Concentrating on the allowed solution, in table 2 we illustrate with a few examples the accuracy of the determination of the unitarity triangle. The first two rows give the assumed input parameters and their experimental errors. The remaining rows give the results for $`\overline{\eta }`$, $`\overline{\varrho }`$, $`\gamma `$ and $`R_b`$ where errors have been added in quadrature.
The accuracy in the scenario I should be achieved at B-factories, FNAL and HERA-B. Scenarios II and III correspond to B-physics at Fermilab during the Main Injector era, LHC-B and BTeV. It should be stressed that this high accuracy is achieved not only because of our assumptions about future experimental errors in the scenarios considered, but also because of the clean character of the quantities considered.
Having the allowed values of table 1 at hand one can calculate $`\epsilon `$, $`\epsilon ^{}/\epsilon `$, $`(\mathrm{\Delta }M)_d`$, $`(\mathrm{\Delta }M)_s`$ and branching ratios for rare decays. As these quantities depend on the parameters characteristic for a given model the results for the SM, the MSSM and other models of this class will generally differ from each other. Consequently by comparing these predictions with the data one will be able to find out which of these models is singled out by experiment. Equivalently, $`\epsilon `$, $`\epsilon ^{}/\epsilon `$, $`(\mathrm{\Delta }M)_d`$, $`(\mathrm{\Delta }M)_s`$ and branching ratios for rare decays allow to determine non-universal unitarity triangles that depend on the model considered. Only those unitarity triangles which are the same as the universal triangle survive the test.
It is of course possible that new physics is more complicated than discussed here and that new complex phases and new operators beyond those present in the SM have to be taken into account. These types of effects would be signaled by:
* Inconsistencies between different constructions of the universal triangle,
* Disagreements of the data with the $`(\mathrm{\Delta }M)_{d,s}`$ and the branching ratios for rare K and B decays predicted on the basis of the universal unitarity triangle for all models of the class considered here.
In our opinion the universal unitarity triangle provides a transparent strategy to distinguish between models belonging to the class considered in this paper and to search for physics beyond the SM. Its other virtues have been listed at the end of the Introduction. Presently we do not know this triangle as all the available measurements used for the construction of the unitarity triangle are sensitive to physics beyond the SM. It is exciting, however, that in the coming years this triangle will be known once $`(\mathrm{\Delta }M)_s`$ has been measured and $`\mathrm{sin}2\beta `$ extracted from the CP asymmetry in $`B_d^0\psi K_S`$. At later stages $`K\pi \nu \overline{\nu }`$, $`BX_{d,s}\nu \overline{\nu }`$, $`B_{d,s}\mu ^+\mu ^{}`$ and future determinations of $`\gamma `$ through CP asymmetries in B decays will also be very useful in this respect.
This work has been supported in part by the German Bundesministerium für Bildung and Forschung under the contract 05HT9WOA0.
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# Kondo effect in mesoscopic systems
## 1 Introduction
Since the discovery of the anomalous low temperature resistivity increase exhibited by some metallic samples these anomalies attracted considerable interest. The first theoretical work to explain them was due to Kondo, who demonstrated that the scattering rate of electrons in metals by magnetic impurities has an anomalous third order contribution, which increases logarithmically as the temperature is reduced and leads to the break-down of perturbation theory . Since then this phenomenon is known as the Kondo effect. Following Kondo’s original work a lot of theoretical effort has been devoted to understand this phenomenon in detail. Wilson’s numerical renormalization group to treat the strong coupling limit and Nozières’ Fermi liquid theory turned out to be the most important milestones in this development.
Recently, the number of papers related to the Kondo effect showed a significant increase with broader and broader applications of the model. Various dilute and dense $`U`$ and $`Ce`$ based metallic alloys have been suggested as Kondo systems with both magnetic and orbital features . In these systems at low temperature very strong correlations build up, hence they became known as strongly correlated systems. Other new developments were in the direction of the observation of Kondo effect in mesoscopic systems such as thin layers and point contacts, and also artificial mesoscopic atoms (quantum dots). In these latter nanofabricated devices the d-level of the magnetic impurity in the metal is mimicked by degenerate states of a quantum dot, which is coupled to metallic or semiconducting leads.
Nanotechnology itself is a very fast developing field, which opens up new perspectives and offers new possibilities to study magnetic impurities and strongly correlated systems. Of course, its extensive overview or a discussion of the physics of nanofabricated artificial atoms is out of the scope of our review. Here we only focus on the study of magnetic and dynamical impurities in mesoscopic systems.
The Kondo effect, in general, originates from the scattering of conduction electrons by a localized object (magnetic or substitutional impurity or some topological defect) with some internal degrees of freedom (e.g. spin, two close atomic positions, dislocation kink). The typical Hamiltonian of Kondo-like problems is
$`H`$ $`=`$ $`{\displaystyle \underset{k\mu }{}}ϵ_kc_{k\mu }^{}c_{k\mu }+{\displaystyle \underset{\alpha }{}}ϵ_\alpha b_\alpha ^{}b_\alpha `$ (1)
$`+`$ $`{\displaystyle \underset{k,k^{}}{}}{\displaystyle \underset{\mu \nu \alpha \beta }{}}V_{\mu \nu }^{\alpha \beta }c_{k\mu }^{}c_{k^{}\nu }b_\alpha ^{}b_\beta `$
where $`ϵ_k`$ is the electron kinetic energy with momentum $`k`$, $`c_{k\mu }^{}`$ creates an electron spherical wave with radial momentum $`k`$ and internal quantum numbers $`\mu `$ and $`b_\alpha ^{}`$ creates a heavy object with quantum number $`\alpha `$ ($`\alpha `$ being the spin, the position, or a crystal field label of the impurity). Note that the internal indices $`\mu `$, $`\nu `$ of the conduction electrons may also represent magnetic spin or orbital indices or a combination of them as well. $`V_{\mu \nu }^{\alpha \beta }`$ denotes the interaction potential and a band cutoff $`D`$ (usually of the order of Fermi energy) is applied for the conduction electrons.
The first corrections to the electron-impurity scattering matrix are given by the two time-ordered diagrams shown in Fig. 1. The direction of time corresponds to the direction of the lines on the heavy objects. Assuming an interaction independent of $`k,k^{}`$, the scattering amplitude for an incoming electron with energy $`\omega `$ is
$$V_{}^{(2)}{}_{\mu \nu }{}^{\alpha \beta }(\omega )=\underset{\rho \gamma }{}[V_{\mu \rho }^{\alpha \gamma }V_{\rho \nu }^{\gamma \beta }V_{\rho \mu }^{\alpha \gamma }V_{\nu \rho }^{\gamma \beta }]\mathrm{ln}(\frac{D}{\omega }),$$
(2)
where the quantum numbers of the internal lines are summed over and the negative sign arises from the fermion anticommutation relations (note the crossed lines in the second diagram). The logarithm above was first identified by Kondo. The divergence of this term as $`\omega 0`$ reflects the break down of perturbation theory.
Different impurity models can be classified by the value of the “commutator” of Eq. (2): (i) For the commutative model $`[V_{\mu \rho }^{\alpha \gamma }V_{\rho \nu }^{\gamma \beta }V_{\rho \mu }^{\alpha \gamma }V_{\nu \rho }^{\gamma \beta }]=0`$, and no Kondo logarithms appear. An example for this is the dissipative tunneling system where the interaction between the heavy object and conduction electrons is diagonal in the internal indices . (ii) For a non-commutative model $`[V_{\mu \rho }^{\alpha \gamma }V_{\rho \nu }^{\gamma \beta }V_{\rho \mu }^{\alpha \gamma }V_{\nu \rho }^{\gamma \beta }]0`$, and logarithmic terms appear in the scattering matrix. In the following we only concentrate to the second case.
In many cases, the general form of the interaction Hamiltonian can be simplified by introducing appropriate variables and truncating the Hilbert space to
$$H_{int}=\underset{\genfrac{}{}{0pt}{}{i=x,y,z}{\alpha ,\beta =\pm ,s=1\mathrm{}n}}{}V^ib_\alpha ^{}\sigma _{\alpha \beta }^ib_\beta c_{k\mu s}^{}\sigma _{\mu \nu }^ic_{k^{}\nu s},$$
where $`\sigma ^i`$’s are the Pauli operators, $`V^i`$’s are anisotropic couplings and the conduction electron may have an additional channel index $`s`$ which is conserved and does not occur in the coupling itself. As it has been pointed out by Nozières and Blandin its mere existence can change the nature of the low temperature behavior drastically. According to the possible different values of that channel indices we speak about n-channel Kondo problem ($`n=1,2,\mathrm{}`$).
In the original spin Kondo problem $`V^i=J/2`$ where $`J`$ is the exchange coupling between the localized spin described by the spin indices $`\alpha ,\beta `$ and the conduction electron spins labeled by $`\mu ,\nu `$. $`b_\alpha ^{}`$ creates the localized electron state with spin $`S`$, for simplicity $`S=1/2`$. In that model there is no additional channel, thus $`n=1`$. In this case the localized $`S=1/2`$ spin is screened by the “compensation cloud” of the electrons and finally a singlet is formed. The binding energy is proportional to the Kondo temperature $`T_K`$ which for the isotropical case can be written as
$$T_K=D(2\varrho _0J)^{1/2}\mathrm{exp}(1/2\varrho _0J),$$
(3)
with $`\varrho _0`$ the density of states of the conduction electrons at the Fermi level for one spin direction. At $`TT_K`$ thermally excited conduction electrons cannot break the singlet, which acts as a rigid potential scatterer and Nozières’ Fermi liquid theory holds .
The actual size of the compensation cloud is given by the Kondo coherence length
$$\xi _K\frac{v_F}{T_K},$$
(4)
with $`v_F`$ the Fermi velocity. For $`T_K1K`$ it can easily exceed $`1\mu m`$, which is in the range of the size of a mesoscopic sample or even can be much larger. This observation triggered the experimental study of the dilute Kondo alloys in thin films, wires and point contacts (see Sec. 2.1 and 2.2).
Contrary to the $`n=1`$ case, for the two channel Kondo model ($`n=2`$) the ground state is not a singlet because two electrons with different channel indices compete to screen the impurity spin. The complexity of the ground state is reflected in a residual entropy of $`S=k_B\mathrm{ln}\sqrt{2}`$, which shows that the impurity is only partially screened even at $`T=0`$ and that small energy excitations are not frozen out even there. As a result, the linear specific heat coefficient $`C(T)/T`$ and the impurity susceptibility were found to diverge logarithmically. A further surprising result arose from the conformal field theory approach: Affleck and Ludwig showed that (i) the impurity contribution to the resistivity shows a $`\sqrt{\mathrm{max}\{\omega ,T\}}`$ singularity (as opposed to the $`\omega ^2,T^2`$ Fermi liquid behavior) (ii) that the amplitude of scattering to a one-electron state vanishes at $`T=\omega =0`$, and an incoming electron “evaporates” to infinitely many electron-hole excitations once it hits the impurity (see Table 1).
It must be emphasized that for $`n=2`$ an arbitrarily small splitting of the impurity states (produced by magnetic field or a strain field for dynamical defects) provides a cutoff for the non-Fermi liquid behavior and ultimately leads to a crossover to a Fermi liquid state.
## 2 Spin Kondo effect in mesoscopic devices
### 2.1 Size dependence in films and wires
In the last decade, many experiments were performed on thin films and narrow wires of dilute magnetic alloys in search of the Kondo compensation cloud. In these experiments (see Fig. 2) no essential change in the Kondo temperature was observed, however, in most of them a suppression of the Kondo resistivity amplitude was observed for small sample sizes. Covering a thin layer of magnetic alloys by another pure metal layer, a partial recovery of the Kondo signal was found which was smaller for more disordered overlayers .
The first natural explanation concerning the compensation cloud was ruled out both theoretically and experimentally as the Kondo singlet is formed whenever the level spacing is small compared to the Kondo temperature: $`\delta ϵ<T_K`$. The effect of local density of states (LDOS) fluctuations close to the surface (discussed in the next subsection) is also probably relatively small for the investigated alloys .
Two theories have been developed, that seem to explain the two limiting cases in the experiments: The first, based on weak localization , may be valid in disordered samples , where the smallest system size is large compared to the elastic mean free path and the Kondo anomaly depends on the level of disorder. The other explanation, the theory of spin-orbit-induced surface anisotropy explains all the experiments performed in ballistic samples (i.e., when the size of the sample is in the ballistic region). This surface anisotropy is developed in samples with strong spin-orbit interaction on the non-magnetic host atoms . In this case electrons can mediate information about the geometry of the sample resulting in an anisotropy for the impurity spin nearby the surfaces, but only in those cases where the angular momenta of the localized orbital $`l0`$ (e.g. $`l=2`$). According to Ref. this anisotropy is non-oscillating in the leading order and inversely proportional to the distance $`d`$ measured from the surface. For flat surfaces it is described by the Hamiltonian
$$H_a=K_d(\mathrm{𝐧𝐒})^2$$
(5)
where $`𝐧`$ is the normal vector of the surface and $`𝐒`$ is the spin operator of the impurity. The anisotropy factor $`K_d`$ is positive and is in the range of $`\frac{0.01}{(\text{d}/\text{Å})}eV<K_d<\frac{1}{(\text{d}/\text{Å})}eV`$ . An elegant extension of these calculations to general geometries was performed by Fomin and coworkers who investigated the dependence of the anisotropy on the roughness of the surface as well.
The Kondo resistivity of thin films was calculated assuming that the two surfaces of the thin films contribute to the anisotropy in an additive way, e.g. $`K_{d,t}=\frac{\alpha }{d}+\frac{\alpha }{td}`$, which was justified later in Ref. . The Kondo temperature was found only slightly affected in a sample of finite size, in agreement with the experiments .
The Kondo signal becomes reduced because close to the surface the motion of the Kondo spins is hindered by the spin-anisotropy , but the size dependence in the Kondo resistivity amplitude $`B(t)`$ defined by $`\mathrm{\Delta }\rho _{\mathrm{Kondo}}=B(t)\mathrm{ln}T`$ is different for integer and half-integer spins. For integer spins (e.g. $`S=2`$ for $`Fe`$) an impurity close enough to the surface is frozen to the $`\mathrm{𝐧𝐒}=0`$ state and the Kondo effect is impossible. Therefore the amplitude $`B(t)`$ is reduced with respect to its bulk value and goes to zero as the film thickness is decreased (see Fig. 3) . For half integer spin (e.g. $`S=5/2`$ for $`Mn`$), on the other hand, the lowest energy state is the $`\mathrm{𝐧𝐒}=\pm 1/2`$ doublet, thus the impurity still produces Kondo resistivity. As a consequence, the size dependence is much weaker, and $`B(t)`$ remains finite even for $`t0`$ (see Fig. 3) . These results are in agreement with the experiments . For $`S=1/2`$ spin alloys (e.g. $`La_{1x}Ce_x`$ films) no size dependence is expected due to the surface anisotropy in agreement with the experiment of Ref. .
The proximity effects can be well explained by the surface anisotropy as the number of available spin-orbit scatterers is increased by the overlayer and the magnetic impurities are in further distances from the surface of the samples, but only if the overlayer is in the ballistic region as well. In a new experiment of Giordano different multilayers composed of $`Au`$ and $`Au(Fe)`$ films were examined (where the overlayer was positioned only on one side, or on both sides of the film), giving good agreement also quantitatively with the predictions of the theory of surface anisotropy.
There are experiments where quantities different from the Kondo resistivity were measured in order to test the theory of anisotropy as well. First Giordano measured the magnetoresistance of thin films and found also a size dependence as the splittings due to the magnetic field and the surface anisotropy compete. Magnetoresistance calculations gave excellent agreement with these measurements. Thermopower and impurity spin magnetization measurements on samples with reduced dimensions can also be explained by the theory of surface anisotropy.
### 2.2 Size dependence in point contacts
Parallel to the thin film experiments, a thorough study of the Kondo effect in ultra small $`CuMn`$ point contacts (PCs) has been carried out . Rather surprisingly, in this case not a suppression but an orders of magnitude increase of both the Kondo signal and the Kondo temperature has been reported.
As shown in Ref. , these anomalies can be well explained by the presence of LDOS fluctuations: For a small PC, even a weak channel quantization induces huge LDOS fluctuations which become larger and larger with decreasing contact sizes (see Fig. 4). As $`T_K`$ depends on the LDOS exponentially (see Eq. (3)), this may produce an extremely wide distribution of the Kondo temperatures for impurities in the contact region. The zero bias anomaly of the PC, however, turns out to be dominated by magnetic impurities with the largest $`T_K`$, since these are the ones that show a well-developed Kondo resonance. Indeed, in Ref. the effect of these fluctuations was taken into account through a modified renormalization procedure, and a perfect agreement was found between the calculated and experimentally determined anomalous amplitude of the Kondo signal (see Fig. 5).
It was also predicted by the theory that this effect should be much less pronounced for alloys with large $`T_K`$ as $`T_K`$ is less sensitive to the change in $`\varrho _0`$ in that case (see Eq. (3)), which has indeed been later confirmed by the experiments studying $`Cu(Fe)`$ alloys .
### 2.3 Kondo resonance in the density of states measured by STM
It has been known for a long time that the local electron density of states nearby a magnetic Kondo impurity has a specific structure due to the Kondo resonance. In the early experimental attempts a change in the electron density of states due to a layer of dilute magnetic alloys fabricated inside a metal has been measured by an oxide tunnel junction placed in a few atomic distances from that layer, and the Kondo structure was indeed observed.
Recently, several groups have demonstrated using scanning tunneling microscopy (STM) that a magnetic Kondo impurity adsorbed on the surface of a normal metal produces a narrow, resonance-like structure in the electronic surface density of states (DOS), whose asymmetric line shape resembles that of a Fano resonance . The experiments were performed with single $`Ce`$ atoms on $`Ag`$ as well as with single $`Co`$ atoms on $`Au`$ and $`Cu`$ surfaces by measuring the I-V characteristics of the tunneling current through the tip of a STM placed close to the surface and at a small distance $`R`$ from the magnetic atom (see Fig. 6 (a)).
A systematic study of the local electronic structure of individual transition-metal impurities on Au surfaces was performed by Jamneala et al. who showed that for elements near the end of the $`3d`$ row ($`Ti`$, $`Co`$, and $`Ni`$) the above mentioned narrow resonance structure appears, whereas for the elements around the center of the row ($`V`$, $`Cr`$, $`Mn`$, $`Fe`$) the electronic structure is found to be featureless. In these experiments the electron tunnels from the tip into the metal, travels to the impurity and, after scattering off it, goes back to the tip, resulting in an interference between the unperturbed and scattered electrons. There is also a possibility that the electrons tunnel from the tip directly to the magnetic impurity, more precisely into the d- or f-level of the atom. However, the tunneling rate for the latter process is probably very small, especially for f-levels, which are deeply inside the atom.
The first theory proposed in Ref. takes into account both processes. Recently, it has been shown in Ref , however, that the Fano resonance can develop even if one neglects the direct tunneling to the impurity. In this theory the physics is governed by the unperturbed one-electron Green’s function at the surface of the metal $`𝒢_{R,\sigma }^{(0)}(\omega i\delta )`$ and the scattering amplitude $`t_\sigma (\omega i\delta )`$ due to the impurity. The latter, given by $`\frac{\mathrm{\Delta }}{\pi \rho _0}G_{d,\sigma }(\omega i\delta )`$ in the Anderson model , can be approximated as
$`t_\sigma (\omega i\delta )={\displaystyle \frac{\mathrm{\Delta }}{\pi \rho _0}}({\displaystyle \frac{Z_d}{\omega \overline{\epsilon }_di\mathrm{\Delta }}}`$
$`+{\displaystyle \frac{Z_U}{\omega \overline{\epsilon }_d\overline{U}i\mathrm{\Delta }}}+{\displaystyle \frac{Z_K}{\omega \epsilon _KiT_K}}),`$ (6)
where $`G_{d,\sigma }(\omega i\delta )`$ is the d-electron Green’s function, $`Z_d`$, $`Z_U`$ and $`Z_K`$ are the appropriate strength of the poles, $`\overline{\epsilon }_d`$, $`\overline{U}`$, $`\mathrm{\Delta }`$, and $`\epsilon _K`$ are the energies of the singly and doubly occupied orbitals of the effective model , the broadening of the d-level, and the position of the Kondo resonance, respectively . The final expression for the tunneling density of states reads
$$\delta \rho _R(\omega )=\frac{[Im𝒢_R^{(0)}(\omega ^{})]^2}{\pi \rho _0}\left\{\frac{(q_R+\epsilon )^2}{\epsilon ^2+1}1+C_R\right\}$$
(7)
where the spin index $`\sigma `$ was dropped, $`\omega ^{}=\omega i\delta `$, $`\epsilon =(\omega \epsilon _K)/T_K`$ and $`q_R=Re𝒢_R^{(0)}(\omega ^{})/Im𝒢_R^{(0)}(\omega ^{})`$. $`C_R`$, which depends on $`Z_d`$, $`\overline{\epsilon }_d`$, $`\mathrm{\Delta }`$, and $`q_R`$ , arises from potential scattering on the d-level and corresponds to a weakly energy dependent Friedel oscillation. The first part of Eq. (7) coming from the scattering by the Kondo resonance gives a Fano line shape in the tunneling LDOS, controlled by the parameter $`q_R`$. The fit on the experimental data for a Co atom on a Au (111) surface gave excellent agreement with fitting parameters being consistent with the predictions of an NCA calculation combined with band structure results .
In the experiment of Jamneala et al. the Kondo resonances were not observed in case of $`V`$, $`Mn`$, $`Cr`$, and $`Fe`$ atoms. In case of $`Mn`$, the Kondo temperature is small for bulk samples and it is further reduced by the weaker exchange coupling at the surface, thus the resonance cannot be expected on meV scale. In case of $`Fe`$ and $`Cr`$ the surface anisotropy described by Eq. (5) may be also reduced, but even in that case, that may make impossible to see spin $`S=2`$. In case of $`Co`$ according to the electronic structure calculations , the spin on the surface is close to $`S=1/2`$, where the anisotropy does not play a role.
To calculate the distance dependence of $`q_R`$ and $`C_R`$, i.e., of the line shape, the tunneling of electrons from the tip (1) into the 3-dimensional $`Au`$ bulk states as well as (2) into the 2-dimensional $`Au(111)`$ surface band was considered. In both cases a free electron-like band structure was assumed . Whereas the periodic changes of the line shape between Fano and Lorentzian ones and the decrease in the overall amplitude with increasing distance were demonstrated (see Fig. 6 (b)), the precise dependence of the line shape on $`R`$ is not reproduced by our simplifying assumption of a free electron band structure . That will require taking into account the detailed band structure as well as the additional scattering phase shift induced by the charge of the Co ion and the charge distribution around the $`Co`$ ion.
Finally, it is worth to comment on the Fano parameter. In the original paper of Fano the parameter $`q`$ is defined in a way where it is proportional to $`q^2|(\mathrm{\Phi }|T|i)|^2/|(\mathrm{\Psi }_E|T|i)|^2`$. The discrete level and the continuum state of energy $`E`$ have the wave function $`\phi `$ and $`\mathrm{\Psi }_E`$, respectively, and there is hybridization between them with amplitude $`V_E`$. The initial electron state $`|i`$ has a transition described by operator $`T`$ to the exact states including the hybridization. The discrete state $`\mathrm{\Phi }`$ in definition of $`q^2`$ is, however, not the original unhybridized state $`\phi `$, but the state modified by the hybridization as $`\mathrm{\Phi }=\phi +𝒫𝑑E^{}V_E^{}\mathrm{\Psi }_E^{}(EE^{})^1`$ . Thus $`(\mathrm{\Phi }|T|i)`$ can be different from zero even if $`(\phi |T|i)=0`$, i.e., thus even without direct transition to the localized state $`\phi `$, a transition rate still exists into the state $`\mathrm{\Phi }`$. That clearly shows that Fano line shape can be obtained without direct transition to the localized d- or f-states in the present case.
## 3 Possibility of two-channel Kondo effect due to structural defects
It is well established by now that scattering on fast dynamical defects can produce Kondo-like anomalies. In the simplest model the defect atom tunnels between two positions and thus forms a two-level system (TLS). These two levels are typically split due to the spontaneous tunneling between the positions and the asymmetry of them, resulting in a typical splitting of $`\mathrm{\Delta }1100K`$ (see Fig. 7).
In the TLS Kondo model the coordinate of the dynamical impurity is coupled to the angular momentum of the conduction electrons through an effective exchange interaction, and the real spins of the conduction electrons act as silent channel indices. Consequently – in the absence of splitting – the physics of the TLS is described by the two-channel Kondo model predicting a NFL behavior below $`T_K`$. In this model the spin-flip scattering of the original Kondo model is replaced by electron assisted tunneling.
### 3.1 Point contacts
Several experiments have been reported where the observed low temperature anomalies were attributed to TLS Kondo defects (see Fig. 8).
In all these experiments the observed anomalies were unambiguously due to dynamical structural defects: they disappear under annealing and not or only slightly depended on magnetic field.
A logarithmic increase of the resistivity attributed to the presence of dislocations or substitutional tunneling impurities has been observed in various systems . However, the most spectacular experiments were carried out in $`Cu`$ and $`Ti`$ point contacts where a two-channel Kondo-like $`\sqrt{T}`$ and $`\sqrt{V}`$ non-Fermi liquid scaling behavior due to non-magnetic scatterers has been observed in the contact resistance . The widths of the zero bias anomalies were associated with the Kondo temperature, $`T_K5K`$. In another beautiful experiment a fluctuation of the zero bias anomaly between two curves due to some slow TLS’s has been observed in amorphous point contacts , which could be consistently explained assuming that a slow fluctuator influences the splitting of one or two fast Kondo two-level systems close to it . There is a further experiment where an alternating voltage was superimposed on a constant bias $`V_0`$, $`V(t)=V_0+V_1\mathrm{cos}(\omega t)`$. As far as the characteristic frequency (e.g. Kondo temperature) of the mechanism responsible for the zero bias anomaly is large compared to $`\mathrm{}\omega `$ the measured $`IV`$ characteristic is just the time average of the current: $`I(t)=I(V_0)+\frac{1}{4}\left(\frac{^2I}{V^2}\right)_{V=V_0}V_1^2`$. For frequencies higher than this scale an $`1/\omega `$ dependence is expected. No deviations have been observed experimentally even for $`\nu =60`$ GHz ($`2.4`$ K) implying that $`T_K>5\mathrm{K}`$. This lower bound is in agreement with the value of $`T_K`$ estimated from the width of the zero-bias anomaly.
There remain, however, a number of puzzles. In all these experiments the estimated Kondo temperature is in the range $`T_K10K`$. $`T_K`$ has been first estimated in Ref. assuming that TLS’s are formed by a heavy atom that tunnels within a distance of about $`0.4\AA `$, and was found to be in the range $`0.011K`$. It has been suggested that virtual hopping to the lowest excited states could increase $`T_K`$ substantially , however, in the above model this turned out to be wrong . It has been shown that, in reality, $`T_K`$ is reduced even further if one includes the effect of all the excited states. Thus, within the original simplistic TLS model, where the TLS is formed by some heavy atom tunneling between two close positions and electron-hole symmetry is assumed, it seems to be impossible to have $`T_K`$ in the experimentally observed range. It has been pointed out recently, that the criticism of Ref. is essentially based on the assumption of electron-hole symmetry in the conduction electron density of states . Electron-hole symmetry, however, is strongly violated in any realistic band structure. As shown in Ref. , $`T_K`$ can be enhanced by orders of magnitude with a relatively small electron-hole symmetry breaking even when all excited states are included, and $`T_K`$ can be in a much broader range than expected previously.
On the other hand, one has very little knowledge about the microstructure of the TLS’s, and in order to make any quantitative prediction it would be extremely important to identify it. We expect that TLS’s with a relatively small effective mass (such as Hydrogen stuck at the surface of the sample or dislocation jogs ) could be able to tunnel over a distance of $`1\AA `$ and probably produce a $`T_K`$ in the experimentally observed range.
Another interesting question is related to the splitting of the two levels, which provides a lower cutoff for the NFL scaling. The presence of splitting and the cutoff of NFL behavior has been observed in several experiments. In particular, measurements on $`Ti`$ point contacts are in perfect agreement with all predictions of the TLS Kondo model. In Ref. , however, the number of TLS’s has been estimated to be about 50, for which concentration already a significant deviation from the NFL scaling should have appeared due to the presence of disorder generated splitting . However, no such deviation has been reported in Ref. . The resolution for this puzzle may also be related to the precise microstructure of the tunneling impurities.
### 3.2 Electron dephasing time $`\tau _\varphi `$
Recent developments in mesoscopic physics raised an interesting question about the electronic dephasing time $`\tau _\varphi `$. This is the time scale for an electron to stay in a given exact one-electron state in the presence of static impurities. The transitions between these states are due to electron-phonon, electron-electron, electron-dynamical defects (e.g. two-level system), or electron - magnetic impurity interactions. At low temperature the electron-phonon interaction freezes out and the electron-electron interaction becomes dominant which has been studied by Altshuler and his collaborators . According to that theory the dephasing time tends to infinity as the temperature is lowered, as the available phase space for the electron-electron scattering gradually vanishes. However, as recently emphasized by Mohanty and Webb , experimentally this is not always the case. In some materials and samples, like $`Ag`$ produced by the Saclay group , $`\tau _\varphi `$ follows very closely the predictions of the electron-electron interaction theory but in other cases there are strong deviations from the predicted behavior and the data indicate some saturations (see Fig. 9).
It has been known since a long time that the saturation-like behavior depends very much on the preparation of the samples and even on the substrates on which the films are deposited . The early suggestion by Lin and Giordano was that the dephasing is either due to the magnetic impurities or some defects which are very sensitive to the metallurgical properties of the films including thickness, annealing etc. The effect of magnetic Kondo impurities on the dephasing rate was carefully studied in those cases where the Kondo temperatures were in the relevant temperature range . There the dephasing rate shows a maximum at the Kondo temperature of the magnetic impurities due to the enhancement of spin-flip scattering, but at lower temperature singlet Kondo ground state is formed and $`\tau _\varphi `$ decreases as the spin-flip rate gradually freezes out (see Table 1). However, the saturation observed has no resemblance to this behavior.
The only possibility for magnetic impurities to produce a saturation of $`\tau _\varphi `$ would be if their Kondo temperatures were much smaller than the temperature range of interest. Then the spin-flip rate is approximately temperature independent. However, in order to have $`T_K10mK`$ the exchange coupling must be very weak $`J\rho _03\times 10^2`$ and to produce the dephasing rate observed an enormous number of unidentified magnetic impurities should be present, which is very unlikely.
Accepting that the low temperature dephasing anomalies are intrinsic properties of the samples and far from a universal behavior it looks reasonable that some local dynamical defects as TLS’s are responsible for them.
Depending on the electron-TLS interaction two different limits must be considered: (i) For weak couplings the electron induced transition in the TLS is treated in second order perturbation theory. In that case to get an almost temperature independent dephasing rate the splittings $`\mathrm{\Delta }`$ (excitation energies) must be smaller than the measured temperature and their distribution must be peaked at very low energies . However, there is no evidence for such anomalous distribution, and linear specific heat measurements on metallic glasses are consistent rather with a uniform distribution . (ii) The other theoretical possibility is given by TLS’s with 2CK behavior . In that case, in contrary to the magnetic Kondo problem, the scattering rate at low temperature is due to processes where the final states contain many electron-hole pairs (see Table 1), and being a dynamical scattering process, this produces dephasing.
In order to get a reasonable dephasing rate less than $`1ppm`$ 2CK defect is required. This explanation has two drawbacks: The questionable existence of such 2CK defects and the required small splitting $`\mathrm{\Delta }`$, even if this latter is renormalized downwards due to the strong interaction by a factor $`\mathrm{\Delta }/T_K`$ if $`\mathrm{\Delta }<T_K`$ . On the other hand in case of 2CK defects the non-universality and metallurgical dependence are quite natural.
Finally it should be mentioned that such saturation like behavior has also been seen in degenerate semiconductor ballistic quantum dots . In degenerate semiconductor as far as we know even the theory of magnetic Kondo effect has not been worked out in detail.
### 3.3 Energy distribution of electrons in short wires
In addition to the dephasing problem another closely related dilemma exists concerning the energy distribution of electrons in short wires with finite bias voltage. In these experiments one measures the energy distribution function of electrons $`f(E)`$ by fabricating a metal-metal oxide-superconductor (M-MO-S) tunnel junction at various positions along the wire (see Fig. 10).
From the $`IV`$ characteristic $`f(E)`$ is determined by deconvolution. The typical length of the wires was $`1.5\mu m`$ and $`5\mu m`$. The measurements were carried out at $`25mK`$, and the samples were in the diffusive limit. In this case the typical energy relaxation processes are slow compared to the time it takes an electron to diffuse through the sample. Therefore the electron distribution exhibits typically two steps corresponding to the Fermi energy of the left and right contacts, thus their difference is proportional to the applied bias $`U`$ . For non-interacting electrons at distance $`x`$ measured from one of the contacts the distribution function is
$$f_x(E,U)=\left(1\frac{x}{L}\right)f^0(E+eU)+\frac{x}{L}f^0(E),$$
(8)
where $`L`$ is the length of the sample and $`f^0`$ is the equilibrium distribution. In the diffusive limit without energy relaxation these two steps are smeared only by the temperature. At low enough temperature this smearing becomes, however, much larger than $`T`$, and the measured smearing gives information on the relaxation processes. In the case of long samples the smearing can be essential as the electrons spend longer time in the sample before leaving into the contacts. Typically the applied voltage $`U0.10.2meV`$ is larger than the temperature $`T30mK`$. The surprising results for short $`Cu`$ wires were that the shape and the amplitude of the smearing could not be explained by the electron-electron interaction. Using a Boltzmann equation approach the line shape could only be reproduced by assuming an anomalously strong electron-electron interaction kernel $`K`$ with an anomalous dependence on the transferred energy $`K1/\epsilon ^2`$ , clearly in disagreement with the predictions of the theory of electron-electron interaction . On the other hand, it has been shown many years ago that such dependence can be due to magnetic impurity mediated inelastic scattering for $`T>T_K`$. The same dependence is valid for two-channel Kondo impurities. For $`Ag`$ samples, however, a good agreement has been found with the electron-electron interaction theory with the expected energy dependence $`K\epsilon ^{3/2}`$ and amplitude.
The close similarity of this problem to the dephasing time dilemma became obvious, when the dephasing time was directly measured on samples prepared in the same way, and it was found that $`\tau _\varphi `$ in the $`Ag`$ samples follows again the standard electron-electron interaction theory and does not saturate.
$`Au`$ wires prepared in Saclay, on the other hand, show also a strong relaxation rate and exhibit an electron distribution that could be fitted using a $`K\epsilon ^2`$ kernel. The strong anomalies in the electron distribution and the dephasing time of $`Cu`$ and $`Au`$ wires and the fact that $`Ag`$ wires behave “regularly” both in dephasing time and electron distribution experiments strongly indicate that the anomalous dephasing time and distribution function are related to the same non-universal and material dependent processes.
It has been realized that assuming a $`K\epsilon ^2`$ electron-electron interaction kernel for energies $`EkT`$ an energy independent electron relaxation rate can be derived making a connection between the distribution function dephasing time experiments. It has been shown that the $`Au`$ and $`Cu`$ experiments can be well explained by 2CK impurities with negligible splitting by treating the 2CK problem in the framework of non-crossing approximation (NCA) and handling the non-equilibrium situation within the quantum Boltzmann equation framework.
That method had been applied earlier to the non-equilibrium transport through nano-point contacts in the presence of 2CK defects . Kroha also showed that the electron-electron interaction mediated by 2CK centers shows a $`K\epsilon ^2`$ energy dependence . That suggests that the classical Boltzmann equation with that kernel may give similar results to those obtained by using quantum Boltzmann equation and the 2CK scattering.
The most challenging feature of the experiments is that for an intermediate range of applied bias $`U`$, $`0.1mV<U<0.5mV`$ the measured distribution function follows a scaling
$$f_x(E,U)=f_x(\frac{E}{U}).$$
(9)
In the framework of the 2CK interpretation the space-dependent non-equilibrium quantum Boltzmann equation is solved, where the collision term is expressed by the non-equilibrium electron self-energy calculated in a self-consistent way. In these calculations $`T_K>1K`$ is assumed which determines the energy scale of the 2CK effect with zero applied bias. The Kondo effect is due to the sharp step of the unperturbed energy distribution function at the Fermi edge and the size of the step determines the Kondo temperature. As far as the bias $`U`$ is not larger than $`T_K`$, from the point of view of the Kondo effect there are no two separate steps and no scaling holds as the presence of $`U`$ influences the shape, width and amplitude of the Kondo resonance. For $`T_K<U`$, however the two steps are separated, and two distinct Kondo resonances are formed at the energy of each step with an effective Kondo temperature which can be essentially smaller than the equilibrium one. That bias region can be described by scaling according to Kroha’s NCA results . At even higher voltage other parameters of the model start to be essential and the scaling breaks down again. This theory is in good agreement with the experimental results for the complete applied voltage range $`U`$ and different positions of the measuring M-MO-S diodes and for samples of various length. For $`Cu`$ and $`Au`$ samples the anomalous smearing can be obtained but for $`Au`$ a much higher concentration of 2CK centers is needed ($`100ppm`$) and for $`Cu`$ $`15ppm`$. These values are in complete accordance with the amplitude of the measured dephasing rate. It should be emphasized that all these data, taken from the measurements of the $`Cu`$, $`Au`$ samples prepared in Saclay, very likely depend a lot on the sample preparation, metallurgy and maybe also on the substrates used.
## 4 Discussions and perspectives
The spin Kondo problem in mesoscopic systems has been thoroughly studied both experimentally and theoretically and is quite well understood. There remain, however, a few further questions to answer: Concerning the surface anisotropy, the crossover from the ballistic to the dirty limits and the effect of disorder on the surface anisotropy should be further clarified. Similarly to the surface anisotropy, local density of states fluctuations decay as $`1/d`$ as a function of the distance from the surface. These fluctuations give probably the dominant effect in very thin films and films with a weak spin-orbit interaction for alloys with a relatively small $`T_K`$. Measurements on a host with weak spin-orbit scattering could help to clarify these issues.
The observation of the Kondo resonance by STM due to a single Co atom on the surface is a very impressive technical achievement. In the future, it would be worthwhile to study magnetic impurities inside the first few surface layers to establish stronger coupling between the spin and the host metals. In order to understand the data or to make predictions further electronic calculations are required for the host metal at the surface, the charge redistribution due to the impurity and the value of spin at the impurity atom.
The anomalous behavior of the zero bias anomalies in point contacts, and the dephasing and the transport in short wires can be well described by dynamical impurities with 2CK behavior. The fact that not all the samples studied (like $`Ag`$ prepared in Saclay) show the anomaly supports that the anomalous NFL behavior is intrinsic, non-uniform and very likely depends on the preparation and treatment of the samples, and even on the substrate on which the sample is deposited. The orbital 2CK interpretation of these experiments is very challenging, however, it is far from being well established.
The possibility of the 2CK effect from a heavy tunneling atom has been suggested longtime ago. The original version of this model, however, cannot explain the large Kondo temperature $`T_K110`$K. However, as realized recently, electron-hole symmetry breaking present in all realistic band structure based density of states, can increase $`T_K`$ substantially and is very promising to resolve the long-standing problem of $`T_K`$. That is certainly the most important question to answer concerning the application of the idea of the 2CK problem.
Finally, we want to remark that a more detailed literature of this topic will be found in the Proceedings of the NATO workshop on ’Size dependent magnetic scattering’ held between 28 May and 1 July 2000 in Pécs, Hungary .
## 5 Acknowledgements
We thank all our collaborators and colleagues for fruitful discussions, especially the participants of the NATO workshop held between 28 May and 3 July 2000 in Pécs and Budapest, Hungary. A. Z. benefited from the hospitality of the Meissner Institute and LMU in Munich where he was supported by the Humboldt Foundation. Our research was supported by the Hungarian Grants OTKA T029813, T024005, T030240, and F29236.
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# Untitled Document
The Nobel Prize in Physics 1999
The last Nobel Prize of the Millenium in Physics has been awarded jointly to Professor Gerardus ’t Hooft of the University of Utrecht in Holland and his thesis advisor Professor Emeritus Martinus J.G. Veltman of Holland. According to the Academy’s citation, the Nobel Prize has been awarded for ’elucidating the quantum structure of electroweak interaction in Physics’. It further goes on to say that they have placed particle physics theory on a firmer mathematical foundation. In this short note, we will try to understand both these aspects of the award. The work for which they have been awarded the Nobel Prize was done in 1971. However, the precise predictions of properties of particles that were made possible as a result of their work, were tested to a very high degree of accuracy only in this last decade. To understand the full significance of this Nobel Prize, we will have to summarise briefly the developement of our current theoretical framework about the basic constituents of matter and the forces which hold them together. In fact the path can be partially traced in a chain of Nobel prizes starting from one in 1965 to S. Tomonaga, J. Schwinger and R. Feynman, to the one to S.L. Glashow, A. Salam and S. Weinberg in 1979, and then to C. Rubia and Simon van der Meer in 1984 ending with the current one.
In the article on ‘Search for a final theory of matter’ in this issue, Prof. Ashoke Sen has described the ‘Standard Model (SM)’ of particle physics, wherein he has listed all the elementary particles according to the SM. These consist of the matter particles: the quarks and leptons along with various vector bosons $`\gamma ,W^\pm ,Z^0`$ and gluons $`g`$ which mediate the various interactions between them. Box 1 summarises them here again for sake of completeness. As explained in that article, one of the conceptual cornerstones of the current description of particle physics is the fact that an interaction (say Coulomb) between two elementary particles (say electrons) can be understood either (i) as the effect of the force field generated by one of them on the other one or equivalently (ii) as arising due to an exchange of the carrier of the force (photon in this case) between them. The photon is the ‘quantum’ of the electromagnetic field. The range as well as the dependence of this force on the relative spins and positions of the particles is correlated with the properties of this quantum and this can be established in a well defined mathematical framework. Box 2 depicts this equivalence in a pictorial manner.
We know that just as Newtonian mechanics is the right mathematical framework to describe the terresterial and celestial motion, quantum mechanics is the right language to describe the motion of molecules, atoms, electrons, neutrons/protons at the molecular/subatomic and subnuclear level. If we want to describe, in addition to these, creations and annihilation of particles, e.g. as happens in the spontaneous transitions of an atom, we need to further extend this mathematical framework to the next higher level of sophistication called ‘Quantum Field Theory’ (QFT). In QFT not only that we employ fields to describe the carriers of interaction, the matter particles are also described by matter fields.
Another cornerstone of our theoretical understanding of the fundamental particles and their interactions is the realization of the important role played by Symmetries / Invariances. The idea of symmetries can be understood in the following way. Laws of physics,let us say $`\stackrel{}{F}=m\frac{d^2\stackrel{}{x}}{dt^2}`$, should be the same no matter which point in the universe do we choose as the origin of our coordinate system. This means that the physics is unchanged under a change of the origin of the coordinate system. This is expressed by saying that the system is invariant under a transformation of coordinates involving translations in space. As per our current understanding, underlying,fundamental invariance principles actually dictate the form of interactions. Let us understand it by taking the example of gravitation. Newton deduced the law of gravitation from the observation of motion. On the other hand Einstein wrote down the general theory of relativity by postulating that the description of motion should be the same for two observers employing two coordinate systems which are related to each other by a general transformation. In particular the transformation can be different at different points in space-time. The non-relativistic limit of this theory (i.e. when objects move at speeds much lower than light) contains Newton’s theory of gravitation. Thus the ‘general co-ordinate invariance’ ‘explains’ the laws of gravitation ‘deduced’ by Newton. So in some sense we have a theoretical udnerstanding of an observed law of nature in terms of a deeper guiding principle. The tenet of current theoretical description of particle physics is that the Quantum Field Theories which have certain invariances are the correct theoretical framework for this description.
The invariance that is most relevant for the discussion here is the so called ‘local gauge invariance’. Without going into the details of the idea, let us just note that this is basically a generalization of the idea that in electrostatics the Electric field and hence the electrostatic force depends only on the difference in potential, and not on the actual values of the potential, i.e. setting of the zero of the potential scale is arbitrary as far as the force is concerned.
Quantum Field Theories , though now the accepted framework for describing particles and their interactions, were in the doghouse for a long time in the 60s because they used to predict nonsensical, infinite results for properties of particles when one tried to compute them accurately. The difficulties arise essentially because of the nontrivial structure that the vacuum has in QFT. This can be visualized by thinking about effect that a medium has on particle properties; e.g. the transport of an electron in a solid can be described more easily by imagining that its mass gets changed to an ‘effective’ mass. Another example is the polarization of the charges, in a dielectric medium, caused by a charged particle. This polarization can cause a ‘change’ of the charge of the particle. In QFT, vacuum acts as a nontrivial medium. The troublesome part, however, is that when one tries to calculate this change in the charge due to the ‘vacuum polarisation’, one gets infinite results. Tomonaga, Schwinger and Feynman (who got the Nobel Prize in 1965) put the Quantum Field Theoretic description of the electron and photon (Quantum Electrodynamics) on a firm mathematical footing They showed how one can use the theory to make sensible, testable predictions for particle properties (such as a small shift in the energy level of an electron in the Hydrogen atom due to the effect of vacuum polarization), in spite of these infinities. If this can be done always in a consistent manner, then the corresponding QFT is said to be renormalizable. The point to note is that the ‘local gauge invariance’ mentioned earlier was absolutely essential for the proof of renormalizability of Quantum Electrodynamics (QED).
The best example where the predictions of this theory were tested to an unprecedented accuracy is the measurement of gyromagnetic ratio of the $`e^{}`$ viz. $`g_e`$. This is predicted to be 2 based on a quantum mechanical equation which is written down with the requirement that the description of the $`e^{}`$ is the same for two observers moving relative to each other with a constant velocity. (Dirac equation for the cognocsenti). The experimentally measured value is close to 2 but differs from it significantly. In QED one can calculate the corrections to the value of $`g_e=2`$ coming from effects of interaction of the electron whereby it emits a $`\gamma `$ and absorbs it again, in a systematic fashion. Box 3 indicates some of these corrections. The measured value agrees with theoretical prediction to 11 significant places as shown in Box 3.
Thus to summarize so far, the electromagnetic interactions between the electron and photons can be described in terms of a QFT. The description has immense predictive power due to the property of renormalizability that the theory has. The theory has this property only because it of its invariance under a set of transformations called U(1) local gauge transformations.
With this, we come to a point in the history in 1971 when particle physicists had a unified description of electromagnetic and weak interaction in terms of exchange of $`\gamma ,W^\pm `$ and $`Z^0`$. S. Weinberg, A. Salam and S. Glashow later got the Nobel Prize in 1979 for putting forward this EW model. Just as unification of electricity and magnetism by Maxwell had predicted the velocity of light ’c’ in terms of the dielectric constant and magnetic permeability $`ϵ_0`$ and $`\mu _0`$ of the vacuum, this unification predicted values of masses $`M_W,M_Z`$ in terms of the ratios of two coupling strengths, called $`\mathrm{sin}^2\theta _W`$. These coupling strengths are the analogue of the electric charge in QED. Details of these relations are displayed in Box 4. C. Rubia and Simon Van der Meer got the Nobel Prize in 1984 for discovering the $`W^\pm `$ and $`Z_0`$ bosons with masses and decays as predicted by the EW model.
The $`W^\pm ,Z^0`$ bosons were found to have nonzero masses $`(M_W=80.33\pm 0.15GeV,M_Z=91.187\pm 0.007GeV`$ where 1 $`GeV`$ is approximately the mass of a proton). As a result the early efforts to cast this electroweak model in the framework of a QFT, by using a more complicated gauge invariance suggested by a generalisation of QED, met with failure. Their nonzero mass makes a QFT incorporating these bosons noninvariant under these gauge transformations. This makes the theory nonrenormalizable. This means calculating corrections to the relation 1 in Box 4 is again riddled with infinities.
At around the same time P. Higgs and others had proposed a way to write a QFT of massive $`W^\pm ,Z^0`$ bosons, where the mass term did not spoil the gauge invariance of the theory. This required existence of an additional particle called the Higgs boson. This is where ’t Hooft and Veltman stepped in. ’t Hooft demonstrated, in his thesis work and the paper published in Nuclear Physics B in 1971, first that the QFT with massless $`W^\pm `$ and $`Z^0`$ was renormalizable and the invariance of the theory under more complex noncommutative local gauge transformations was essential for that. He further showed that a QFT containing $`\mathrm{𝑚𝑎𝑠𝑠𝑖𝑣𝑒}`$ $`W^\pm ,Z^0`$ bosons would be renormalizable (i.e., coefficients of infinite corrections would vanish identically) inspite of nonzero masses as long as the mass was generated through the mechanism suggested by P. Higgs. Together ‘t Hooft and Veltman developed new methods of calculation for the higher order corrections to particle properties, which explicitly preserved this gauge invariance. This work opened the floodgates of the prospects of using the ElectroWeak theory to make accurate predictions and test the theory to a similar degree of accuracy as the QED cf. Box 3. Veltman led the program of calculation of various higher order corrections to EW quantities, having established that the results were guaranteed to be finite. He actually developed a computer program called ‘Schoonship’ to use the computer to do these very complicated analytical calculations specific to Theoretical High Energy Physics. This is the sense in which the work of ’t Hooft and Veltman put the EW theory on a firm mathematical footing. This work was enough to convince the particle theorists that gauge theories with Higgs mechanism was the way to go to describe EW interactions.
In QED the corrections ( e.g. to $`(g2)_e`$ shown in Box 3) depended only on the mass and charge of an $`e^{}`$, wheras in EW theory they depend on the free parameters of this theory viz. the masses of various quarks and leptons. The corrections are dominated by the top quark due to its large mass. Box 4 shows the leading corrections predicted in the EW theory to the ratio $`\rho =\frac{M_W^2}{M_Z^2cos^2\theta _W}`$. The measurement of $`M_W/M_Z`$ and $`\mathrm{sin}^2\theta _W`$ in 1984 were consistent with $`\rho =1`$ which was the analogue of $`g=2`$ prediction of QED. The measurements then were not precise enough to decide what the deviation of experimentlly measured value of $`\rho `$ from 1 was. In the decade since then, a detailed study of the properties of these bosons has been possible using the 10 million $`Z^0`$ bosons created at the Large Electron Positron Collider (LEP) in Geneva and thousands of $`W^\pm `$ bosons at the $`p\overline{p}`$ collider Tevatron at Fermilab in Chicago. By 1993 $`\rho `$ was found to be $`1.011\pm 0.006`$. This implied, as can be seen from the Box 4, that the top quark, which was not discovered till 1995 must have a mass $`M_t180GeV`$. Finding the top quark in 1995 with a mass consistent with this value indeed tested the predictions of the EW theory to high accuracy. The precision of these measurements meant that if one did not use the corrected expressions, the values of $`M_W^2,M_Z^2`$ and $`\mathrm{sin}^2\theta _W`$ would not be consistent with each other within the SM.
Even though not shown in Box 4, the corrections to this ratio also depend on the mass of the only particle in the SM which is as yet undiscovered viz. the Higgs Boson, albeit very weakly. The figure in Box 5 shows the region in the $`M_WM_t`$ plane that is indicated by measurements today. The straight line shows predictions of the SM for different values of the Higgs boson. So just as five years ago, one used these measurements to ‘determine’ values of $`M_t`$ (which was then not measured) now particle physicists are using them to ‘determine’ the mass of the elusive Higgs particle. These precision measurements narrow down the mass range where the Higgs boson is likely to be found if SM is indeed completely correct. Hunt for this will be on at the Large Hadron Collider (LHC) which will go in action in 2006.
The EW theory predicts a slew of measurable quantities in terms of the basic parameters of the theory viz. the couplings and masses of quarks/leptons. Fig. in Box 6 shows a comparison of the predictions of the SM (corrected for these loop effects) with data. The numbers in the third column indicate the difference between the prediction and measurement in units of the standard deviation. It is this agreement, which would be nowhere as excellent if we do not include the higher order corrections, that has proved that the EW interactions are correctly described in terms of a Quantum Field Theory whose renormalizability was established by ’t Hooft and Veltman’s work. Their Nobel prize is also the recognition of the success of QFT and Gauge Principle which are the two cornerstones of the mathematical description and understanding of the electromagnetic, weak and strong interactions among fundamental particles. The only part of this edifice that is as yet not honoured with a Nobel prize is QCD: Gauge theory of strong interactions. Who knows, in a few years we will be reading about the work of D. Gross, H. Politzer and F. Wilczek in a similar article!
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# The Principle of Valence Bond Amplitude Maximization in Cuprates: How it breeds Superconductivity, Spin and Charge Orders
## Abstract
A simple microscopic principle of ‘Valence bond (nearest neighbor singlet) amplitude maximization ’(VBAM) is shown to be present in undoped and optimally doped cuprates and unify the very different orderings such as antiferromagnetism in the Mott insulator and the robust superconductivity accompanied by an enhanced charge and stripe correlations in the optimally doped cuprates. VBAM is nearly synonymous with the energy minimization principle. It is implicit in the RVB theory and thereby makes the predictions of RVB mean field theory of superconductivity qualitatively correct.
The qualitative predictions of the RVB mean field theory of high Tc superconductivity in cuprates,including the symmetry of the order parameter has turned out to be in good agreement with experiments. In spite of its approximate character it has definitely put us in the correct Hilbert space by focusing on the key singlet correlations. On the other hand various orders, that at a superficial level appear to be unusual from the RVB theory point of view, have emerged both experimentally and theoretically. The aim of the present letter is to identify an unique cause behind the AFM order in the Mott insulator, the enhanced charge and spin stripe correlations and low energy spin fluctuations in the optimally doped superconducting cuprates . A simple principle of ‘valence bond (nearest neighbor singlet) amplitude maximization’ (VBAM) is shown to emerge as a fairly simple cause. This maximization respects certain geometrical and dynamic constraint provided by our strongly correlated system. In the case of the Mott insulator, while super exchange interaction respects VBAM, an irreducible amount of local triplet fluctuations is generated in order to satisfy certain geometric constraints. These triplet fluctuations interact and condense at momentum $`(\pi ,..\pi )`$ resulting in long range AFM order for any dimension $`2`$. In the optimally doped Mott insulator the constraints are modified by the charge dynamics. And VBAM continues as an approximate principle and produces new orders such as the robust superconductivity and the enhanced quasi static charge and spin stripe correlations. We also argue that VBAM is responsible for the qualitative success of the RVB mean field theory.
At a more fundamental level the 2 dimensionality, the one band character, strong correlation and closeness to a Mott insulating state are believed to be responsible for various interesting low energy electronic phenomenon in cuprates. The VBAM hypothesis is an emergent consequence of the above and perhaps an useful guiding principle that at least helps us to separate cause from effects. One of the messages of this letter is that it is VBAM, rather than the development of charge or spin stripe correlations that is providing a mechanism for superconductivity: on the contrary whenever stripes are nurtured it is at the expense of superconductivity.
Our VBAM hypothesis is proved for the case of undoped Mott insulator by simply using the principle of energy minimization in the ground state. For the optimally doped Mott insulator various arguments are provided in support of this principle. The existence of the sharp 41 meV triplet resonance at $`(\pi ,\pi )`$ at optimally doped bilayer YBCO and its counterpart at various dopings all the way up to zero doping is presented as an experimental support for our VBAM hypothesis in the optimally doped regime.
The phase coherent resonance of singlet bonds is at the heart of the RVB theory of superconductivity. Anderson’s RVB wave function,
$$|RVB=P_G(a_{ij}b_{ij}^{})^{\frac{N_e}{2}}|0$$
(1)
by construction has enhanced VB amplitude. It is a remarkably universal wave function in the sense, it can describe an antiferromagnetically ordered Mott insulating state and the robust superconducting state of the optimally doped Mott insulator with equal ease for appropriate choices of the pair function $`a_{ij}`$. We believe that this universal feature is a consequence of the VBAM principle. Also without changing the RVB wave function and consequently the superconducting property in a fundamental fashion, the charge stripe and spin stripe order can be incorporated by a modulation of the pair function $`a_{ij}`$.
To begin with we develop a heuristic picture of VBAM. In a free fermi gas, any short range singlet correlation contained in the ground state is a consequence of Pauli principle rather than interactions. In the large U Hubbard model, accepted as a good model for our narrow band conducting cuprates close to half filling, every elementary collision between two electrons tries to establish a nearest neighbor singlet correlations (a valence bond). That is, the virtual transitions to doubly occupied state on a given copper site lowers the energy (compared to the $`U=\mathrm{}`$ case) by the super exchange energy $`J\frac{2t^2}{U}`$ and stabilizes the spin singlet state rather than a triplet state. Elementary collisions, in addition to the Pauli principle induce singlet correlations. Close to half filling elementary two body collisions are more frequent than free hopping of charges. The on site collision induced valence bond proliferation is at the heart of the VBAM principle.
Let us try to understand the development of long range AFM order in the spin half Mott insulators in the light of the principle of VBAM. In an isolated pair of neighboring orbitals, the electron pair has a spin singlet (non-magnetic) ground state, through the super exchange interaction. An important question is how this local non magnetic singlet tendency manifests itself in a 2 or 3 dimensional lattice. The spin half Heisenberg Hamiltonian with nearest neighbor interaction
$$H=J\underset{ij}{}(𝐒_i𝐒_j\frac{1}{4}),$$
(2)
in terms of the underlying electron variables $`c^{}s`$ and the bond singlet operator $`b_{ij}^{}\frac{1}{\sqrt{2}}(c_i^{}c_j^{}c_i^{}c_j^{})`$ takes the form
$$H=J\underset{ij}{}b_{ij}^{}b_{ij},\text{ with }n_i+n_i2$$
(3)
We have used the important identity
$$(𝐒_i𝐒_j\frac{1}{4}n_in_j)b_{ij}^{}b_{ij}$$
(4)
The singlet number operator $`b_{ij}^{}b_{ij}`$. has an eigen value of 0 or 1 in our single occupancy subspace. In view of the above identity, in a translationally invariant ground state, what is maximized consistent with the lattice structure, is the valence bond amplitude. Thus in the nearest neighbor Heisenberg models minimization of the ground state energy is synonymous with maximization of the strength of nearest neighbor singlet bonds. This proves our VBAM hypothesis for the spin half Mott insulator.
Even though diagonal in terms of the number operators $`b_{ij}^{}b_{ij}`$, equation (2) is not really diagonalized, as the number operators themselves do not commute whenever one of the sites coincide:
$$[b_{ij}^{}b_{ij},b_{jk}^{}b_{jk}]=i(𝐒_i\times 𝐒_k)𝐒_j$$
(5)
This non-commutativity propagates an irreducible minimum of bond triplet fluctuations in the ground state by making the ground state average $`b_{ij}^{}b_{ij}_G<1`$. In fact, $`b_{ij}^{}b_{ij}_G0.6391,0.5846\text{and}=0.5`$ respectively for 1d chain, 2d square lattice and infinite d hyper cubic lattice. In the 1d chain the bond triplet fluctuation is finite and manages to produce an algebraic AFM order. In 2d it increases further and is believed to produce a true long range order and in infinite d it increase even further to produce a perfect Neel order. The chiral operator appearing on the right hand side of equation (5) also tells us that chiral fluctuations are also induced in the ground state.
A long range magnetic order, when it occurs in the ground state, is an inevitable consequence of VBAM in the presence of the constraints provided by the lattice structure and the above commutation relation. Thus in a hyper cubic lattice for $`d2`$
$`\text{VBAM + geometrical constraints}\text{ AFM Order}`$ (6)
The nature of magnetic order is strongly lattice dependent. For the non bipartite 2d triangular lattice, what maximizes the bond singlet amplitude is a $`120^o`$ structure with zero point fluctuations. The case of P doped Si in the insulating state is described by a 3d random lattice Heisenberg model. The nature of lattice constraint being very different long range magnetic order is believed to be absent resulting in a kind of singlet bond glass state.
The role of the long range AFM order should not be also overemphasized in this context. As shown by Liang, Docout and Anderson, the energy difference between a disordered spin liquid state with only short range AFM correlations and the the best variational state with long range AFM order is as small as one or two percent of the total energy:
$$\frac{b_{ij}^{}b_{ij}_{SL}b_{ij}^{}b_{ij}_G}{b_{ij}^{}b_{ij}_G}0.01$$
(7)
At the level of variational wave function a small change in the long distance behavior of the pair function $`a_{ij}`$ takes us between a disordered and ordered ground state. Hsu also shows that the AFM order in 2d is a spinon density wave in a robust spin liquid state. Thus the development of long range AFM order in 2d is a result of a small final adjustment of the VB amplitude in a spin liquid state in the maximization procedure.
We wish to argue that VBAM principle continues to be approximately valid for the optimally doped ($`\delta 0.15`$) regime and produce, in the presence of the constraints modified by the hole dynamics, the robust superconducting order and also the enhanced charge and spin stripe correlations. The discussion will be heuristic and also uses known results from the RVB theory. We will not consider the under doped regime, as disorder and long range coulomb interaction play important roles and suppress superconductivity strongly and produce mesoscopic phase separation complication.
The t-J model
$`H_{tJ}=t{\displaystyle \underset{ij}{}}(c_{i\sigma }^{}c_{j\sigma }+h.c.)J{\displaystyle \underset{ij}{}}(𝐒_i𝐒_j{\displaystyle \frac{1}{4}}n_in_j),`$ (8)
takes the form
$$H_{tJ}=t\underset{ij}{}(c_{i\sigma }^{}c_{j\sigma }+h.c.)J\underset{ij}{}b_{ij}^{}b_{ij}$$
(9)
with the usual constraint $`n_i+n_i2`$. We would like to see if VBAM continues to be valid in the presence of the single electron hopping term in the optimally doped case.
From energy minimization point of view the double occupancy constraint limits the single particle kinetic energy gain per site to $`\mathrm{\Delta }_{KE}tz\delta `$ compared to the larger super exchange energy gain $`\mathrm{\Delta }_{SE}Jz(1\delta ^2)`$. For cuprates, with a co-ordination number $`z=4`$ and $`\frac{t}{J}2`$, when $`\delta 0.15`$ the ratio $`\frac{\mathrm{\Delta }_{SE}}{\mathrm{\Delta }_{KE}}3`$. From the energy considerations given above, the VBAM principle continues to be important in the optimally doped regime. The above rough estimate is in agreement with more accurate estimates using variational monte carlo studies.
The nature of constraints on valence bond proliferation has changed now; we will see how it can stabilize new low temperature phases such as the robust superconductivity along with quasi static charge and spin stripe correlations. Maximizing valence bond amplitude is an important step towards establishing a robust superconducting state. The next important step is the development of ‘in phase resonance’ of the valence bonds (or the zero momentum condensation of the charged valence bonds) in the ground state. This is what is precisely achieved in the RVB mean field theory. In this sense the VBAM principle is satisfied by the RVB mean field theory. While the original RVB mean field theory emphasized the extended-s mean field solution, Kotliar’s identification of $`d_{x^2y^2}`$-wave RVB solution as a lower energy mean field solution made the RVB theory closer to experiments in terms of the symmetry of the order parameter.
After the original RVB mean field theory, and some experimental developments of that time, the inclusion of the on site constraint was suspected to produce a large phase fluctuations and in the even remove the finite temperature K-T transition in an isolated $`CuO_2`$ layer. However, recent experiments have strongly supported the one layer d-wave superconductivity with a large $`T_c`$$`95K`$; in addition the RVB Ginzburg Landau functional derived by Anderson and the present author, in the RVB gauge theory approach did not show any singular effect on the GL coefficients due to the on site constraints. Very recently Lee, using RVB gauge theory has provided a rather convincing non perturbative analysis supporting the local stability of the d-wave solution.
What the experiments and also various theoretical clues have been telling us is that the effect of the strong correlation (on site constraint) does affect superconductivity in unexpected fashions by encouraging quasi static charge stripe and spin stripe order. Within the VBAM principle we can understand it in the following heuristic fashion: the super exchange term tries to segregate the holes so that they have more fluctuating Mott insulating region where super exchange gain (or equivalently the bond singlet amplitude) can be maximized. If the hole segregation has the form of a 1d charge stripes, in addition to VBAM, kinetic energy gain also results from coherent charge delocalization along the charged stripes for the following reason. In view of the Brinkman-Rice phenomenon and proliferation of non self retracing paths in 2d, ‘holes’ can not maximize the kinetic energy gain by coherent delocalization in 2d. The quasi 1d fluctuating stripes on the other hand provides maximum of self retracing paths (along the stripes) for the holes thereby gaining coherent charge delocalization energy.
Since the quasi static charge stripe formation also creates coherently fluctuating Mott insulator region, VBAM leads to enhanced quasi static antiferromagnetic correlations (incommensurate order). Thus charge stripe and spin stripes are intimately related. In the optimally doped region, the development of charge and spin stripe correlations arise from finer adjustments of a robust superconducting state in order to satisfy the VBAM principle. These adjustments can be already done at the level of RVB mean field theory in the following fashion both in the normal state and superconducting state: a) development of spontaneous anisotropic valence bond amplitudes, $`|\mathrm{\Delta }_x||\mathrm{\Delta }_y|`$ and b) one dimensional spatial modulation of $`|\mathrm{\Delta }_{ij}|`$ either along the x or y axis.
The first case corresponds to stripes having a orientational order (along the a or b axis) and no spatial order. In the normal state this will be a nematic metal with $`\rho _a\rho _b`$ and other wise a non fermi liquid state. This state will have an enhanced magnetic correlations around $`(\pi ,\pi )`$, in view of the spin localization caused by the quasi static stripe formation, This anisotropy can continue into the low temperature superconducting phase and will result in a real mixing of $`d_{x^2y^2}`$ and extended-S state. Perhaps this phase is already seen in the under doped LSCO - we think that the conducting charge stripe glass state is an anisotropic metallic state (nematic metal). As explained elsewhere coupling to octahedral rotations and displacements in LSCO further enhance the charge stripe correlation at the expense of superconductivity.
The second case corresponds to a spatial ordering of the stripes, which can also start in the normal state, independent of the low temperature superconducting state. Close to the magic filling of $`\frac{1}{8}`$ this is believed to happen in doped LCO.
In both cases there will be domain formation and the transition from the isotropic metallic phase to the above anisotropic phase can be a second order phase transition. In the optimally doped region, even in the most favorable case of LSCO, the stripe correlations are not stabilized into a true long range order. However, they remain as additional correlations in the ground state there by reducing superconductivity.
At the level of wave function modifications the above secondary stripe orders appear as modifications of the pair function $`a_{ij}`$ without changing the RVB superconducting wave function in a fundamental fashion.
We will discuss briefly the question of low energy spin fluctuations in optimally doped cuprates. It is an experimental fact that in the singlet dominated cuprates there are strong low energy spin fluctuation phenomenon in k-space particularly around the $`(\pi ,\pi )`$ region. Among them a remarkably neat and unique phenomenon occurs in the bilayer cuprates: in neutron scattering a resolution limited 41 meV peak corresponding to a triplet excitation with momentum centered around $`(\pi ,\pi )`$ is seen below the superconducting $`T_c`$. These excitations have very little dispersion over a wide momentum interval around $`(\pi ,\pi )`$, suggesting that we can create spin triplet wave packets of size comparable to the lattice parameter as nearly exact eigen excitations. That is, the triplet excitation is predominantly made of nearest neighbor (Eder also has a nearest neighbor triplet bond picture for the triplet resonance) triplet bonds. Availability of large valence bond amplitude is a prerequisite for being able to create this excited state. It tells us that large amplitude valence bond exist in the ground and they can be indeed converted into triplets carrying a momentum of $`(\pi ,\pi )`$ and energy of about 41 meV. And 41 meV is the stabilization of bond singlet energy in the superconducting state relative to the normal state. Indeed various authors have extracted the superconducting condensation energy from the spectral properties of the 41 meV peak. It should be also mentioned that interlayer pair tunneling adds further stability to the valence bonds in bilayer cuprates there by supporting our VBAM principle.
We suggest that the triplet resonance around $`(\pi ,\pi )`$ is what carries the memory of the Mott insulator. It continues to persist as we go to lower dopings in YBCO and the energy of this peak decreases linearly with the corresponding superconducting transition temperature. In a Mott insulator, where the superconducting $`T_c`$vanishes, this triplet excitation around $`(\pi ,\pi )`$ becomes soft (in view of the existing triplet condensate) and becomes a Goldstone mode of the antiferromagnetic order. This also tells us that the maximization of the valence bond amplitude is easier in a doped Mott insulator than in an undoped one. The hole dynamics in some sense decreases the constraints provided by the lattice, by effectively converting the valence bond operators into more of ‘unconstrained’ bosons.
Thus the various antiferromagnetic fluctuations including the narrow resonance peak tells us about the growing antiferromagnetic correlation in the ground state, arising from the stabilization of the bond singlets. It is in this sense the antiferromagnetic fluctuations are effects of a deep and growing bond singlet tendency, forced by some geometrical and dynamical constraints, rather than some thing that provide cooper pairing at low energies. Spin fluctuations are effects rather than causes of the singlet dominated superconductivity phenomenon.
Thus the maximization of the bond singlet amplitude can be thought of as a driving force behind the various anomalous correlations one sees in cuprates including the robust high Tc superconductivity. The emergent quasi static charge and spin stripe correlations and superconductivity mutually adjust themselves when some parameters such as doping or temperature are changed; they are not to be thought of driving each other. The fundamental driving force is VBAM which is implicit in the RVB theory.
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# 1 Introduction
## 1 Introduction
Self-dual Yang-Mills fields are fascinating mathematical objects that play an important role in both Physics and Mathematics. In the Physics literature they emerged through the introduction of the BPST instanton , the minimum action configuration in the sector of topological charge $`Q`$ equal to one. This triggered the joint effort of physicists and mathematicians in the search for multiinstanton configurations. This had its reward with the ADHM formalism , a general set up for the construction of self-dual configurations with vanishing field strength at infinity (which, from the topological point of view, corresponds to fields living on a four-sphere). However, there are certain instances in which other types of boundary conditions are relevant. For example, in considering finite temperature field theory, one is interested in studying configurations which are periodic in thermal time. This periodicity might also be used as a device to study monopole-like objects. In the same fashion, additional periodicities might have other uses and interpretations, as argued in Ref. . This justifies an initial interest in the study of self-dual configurations which are periodic in all euclidean space-time directions. Geometrically, it corresponds to the study of self-dual gauge fields on the torus. It might be considered as the next step after the case of gauge fields on the sphere. The study of gauge fields on the torus (see for a review of this topic) brings in a new topological richness, whose appearance, physical interpretation and usefulness was put forward by ’t Hooft . In addition, periodic self-dual configurations constitute a simple mathematical example of dense multiinstanton configurations, which, as advocated in Ref. , might turn out to be a better description of the confining vacuum than other dilute multiinstanton pictures.
The construction of self-dual configurations on the torus has met very limited success. For particular values of the torus sizes, there is a class of solutions which is known since early times . These configurations have constant field strength, and they are, in some sense, of an abelian nature: all the spatial components of the electric and magnetic fields are parallel in colour space. For other torus sizes, we know, through numerical techniques , that the solutions are quite different. They turn out to be lumpy and very non-abelian: different electric field components are mutually orthogonal in colour space at the center of the lump. Fortunately, there has recently been substantial progress in studying the mixed situation of instanton configurations which are periodic in only a few directions. In particular, in the case of only one compactified direction, corresponding to the physical situation of finite temperature, the most general solution of topological charge one, the caloron, has been found . Considerable progress has also been achieved for the case of gauge fields on $`T^2\times R^2`$ (doubly periodic instantons), by the work of mathematicians and physicists . In these developments a crucial role is played by the Nahm transform, a duality transformation which maps self-dual configurations on tori with dual sizes . For the $`T^3\times R`$ case, the calculation of the abelian Nahm dual of the charge one instanton represents a step forward. The caloron, the doubly periodic instantons and the $`T^3\times R`$ instanton can be viewed as limiting cases of configurations on the torus where either three, two or one directions are taken to be very large with respect to the others.
This paper is a step towards an analytical understanding of non-constant field strength self-dual configurations on the torus. The strategy is to consider torus sizes which depart only slightly from those in which there exists a self-dual constant field strength solution. Our construction is based upon perturbing around the constant solution (at a torus size in which it is not self-dual), and imposing self-duality to the resulting configuration. A systematic perturbative expansion arises that allows the construction of the self-dual solution. We show that the solution exists order by order and is unique up to gauge transformations and space-time translations. This is done in section 2. Our approach is intimately related to the study of van Baal , who considered perturbations around constant field strength solutions.
In section 3 we proceed to compare the lowest non-trivial order results obtained from our expansion with the exact solution obtained by numerical methods. This serves to quantify the rate of convergence of the series for various torus sizes. Finally, in section 4, we investigate the interplay of our perturbative expansion with the Nahm transform. The paper is closed by section 5, where the conclusions and possible extensions are presented.
## 2 The construction
Let us consider $`SU(2)`$ gauge fields living on a torus of size $`l_0\times l_1\times l_2\times l_3`$. Under a translation by one of the periods, the gauge potentials and fields transform by a gauge transformation:
$$A_\nu (x+e_\mu )=[\mathrm{\Omega }_\mu (x)]A_\nu (x),$$
(1)
where $`e_\mu `$ is a 4-vector of length $`l_\mu `$ along the $`\mu `$-th direction, and $`\mathrm{\Omega }_\mu `$ are the twist matrices. The compatibility conditions of the previous equations are:
$$\mathrm{\Omega }_\mu (x+e_\nu )\mathrm{\Omega }_\nu (x)=\mathrm{exp}\{\pi ın_{\mu \nu }\}\mathrm{\Omega }_\nu (x+e_\mu )\mathrm{\Omega }_\mu (x),$$
(2)
where the elements of the antisymmetric twist tensor $`n_{\mu \nu }`$ are integers defined modulo 2. In what follows we will choose $`n_{03}=n_{12}=n_{30}=n_{21}=1`$, with the remaining components being zero. For the twist matrices we will take
$$\mathrm{\Omega }_\mu (x)=\mathrm{exp}\{ı\frac{\pi }{2}n_{\mu \nu }\frac{x_\nu }{l_\nu }\tau _3\},$$
(3)
which is consistent with our choice of twist tensor. The symbols $`\tau _i`$ label the Pauli matrices.
For torus sizes such that $`l_0l_3=l_1l_2`$, there exist self-dual configurations satisfying the previous boundary conditions and having constant field strength. What we will do is to consider a slight deviation from this situation controlled by the parameter:
$$\mathrm{\Delta }=\frac{l_0l_3l_1l_2}{\sqrt{V}},$$
(4)
where $`V=l_0l_1l_2l_3`$ is the torus volume. Without loss of generality we can assume that $`\mathrm{\Delta }`$ is positive. In this case there exists a constant field strength configuration with vector potential:
$$B_\mu (x)=\frac{\pi }{2}n_{\mu \nu }\frac{x_\nu }{l_\mu l_\nu }\tau _3.$$
(5)
This gauge potential gives rise to a field strength of the form $`G_{\mu \nu }\tau _3`$, where:
$$G_{\mu \nu }=\pi \frac{n_{\mu \nu }}{l_\mu l_\nu },$$
(6)
The only non-zero components are $`G_{03}`$ and $`G_{12}`$, which become of equal magnitude at $`\mathrm{\Delta }=0`$, rendering the solution self-dual. With our choice of twist the constant field strength configuration has topological charge $`Q=1/2`$ and, for $`\mathrm{\Delta }=0`$, total action $`4\pi ^2`$.
Now, let us consider perturbing around this gauge potential:
$$A_\mu (x)=B_\mu (x)+S_\mu (x)\tau _3+W_\mu (x)\tau _++W_\mu ^{}(x)\tau _{},$$
(7)
where we have decomposed the additional field into different colour components. The matrices $`\tau _\pm =\frac{1}{2}(\tau _1\pm ı\tau _2)`$ are standard. The boundary conditions on the gauge fields translate into the real function $`S_\mu (x)`$ being periodic on the box, and the complex function $`W_\mu (x)`$ satisfying
$$W_\rho (x+e_\mu )=\mathrm{exp}\{ı\pi n_{\mu \nu }\frac{x_\nu }{l_\nu }\}W_\rho (x).$$
(8)
We will make the following gauge choice, consistent with these boundary conditions (the background field gauge):
$$_\mu A_\mu (x)ı[B_\mu (x),A_\mu (x)]=0.$$
(9)
We will now demand that the resulting gauge field is self-dual
$`𝐅_{\mu \nu }(x)\stackrel{~}{𝐅}_{\mu \nu }(x)=0,`$ (10)
$`\stackrel{~}{𝐅}_{\mu \nu }(x)={\displaystyle \frac{1}{2}}ϵ_{\mu \nu \rho \sigma }𝐅_{\rho \sigma }(x),\mathrm{with}ϵ_{0123}=1,`$
which will be interpreted as equations for the functions $`S_\mu (x)`$ and $`W_\mu (x)`$. The best way to express these equations, together with the gauge fixing condition, is to use the matrices $`\sigma _\mu (𝐈,ı\stackrel{}{\tau })`$ and $`\overline{\sigma }_\mu (𝐈,ı\stackrel{}{\tau })=\sigma _\mu ^{}`$. These matrices satisfy:
$$\overline{\sigma }_\mu \sigma _\nu =\overline{\eta }_{\mu \nu }^\alpha \sigma _\alpha ,$$
(11)
where $`\overline{\eta }_{\mu \nu }^\alpha `$ is the ‘t Hooft symbol, such that $`\overline{\eta }_{\mu \nu }^0=\delta _{\mu \nu }`$ and the $`\overline{\eta }_{\mu \nu }^i`$ are a basis of the antiself-dual tensors. Now contracting $`𝐅_{\mu \nu }(x)`$ with $`\overline{\sigma }_\mu \sigma _\nu `$, we project out the self-dual part. Hence, we might rewrite equation (10) as follows:
$`\overline{}S={\displaystyle \frac{\lambda }{2}}\widehat{G}+{\displaystyle \frac{ı}{2}}(W_c^{}W_cW^{}W)`$ (12)
$`\overline{D}W=ı(S^{}WW_c^{}S),`$ (13)
where $`S=S_\mu (x)\sigma _\mu `$ and $`W=W_\mu (x)\sigma _\mu `$ are $`2\times 2`$ matrices, $`S^{}`$, $`W^{}`$ their adjoints and the parameter $`\lambda `$ is equal to $`1`$. The matrix $`W_c`$ is the charge conjugate:
$$W_c=\tau _2W^{}\tau _2.$$
(14)
The matrix $`\widehat{G}G_{\mu \nu }\overline{\sigma }_\mu \sigma _\nu `$ is given by:
$$\widehat{G}=2\pi ı\frac{\mathrm{\Delta }}{\sqrt{V}}\tau _3.$$
(15)
It vanishes when $`\mathrm{\Delta }=0`$, since then the constant field strength configuration is self-dual. Finally, we define the following differential operators:
$``$ $`=`$ $`\sigma _\mu _\mu ,`$ (16)
$`\overline{}`$ $`=`$ $`\overline{\sigma }_\mu _\mu ,`$ (17)
$`D`$ $`=`$ $`\sigma _\mu D_\mu ,`$ (18)
$`\overline{D}`$ $`=`$ $`\overline{\sigma }_\mu D_\mu ,`$ (19)
$`D_\mu `$ $`=`$ $`_\mu +ı\pi {\displaystyle \frac{x_\nu n_{\mu \nu }}{l_\mu l_\nu }},`$ (20)
Since at $`\mathrm{\Delta }=0`$ the correction terms $`W`$ and $`S`$ vanish, we can think of $`\mathrm{\Delta }`$ as a perturbation parameter. Rigorously speaking this is not quite so, because $`\mathrm{\Delta }`$ depends on the torus sizes, and they enter also in the boundary conditions. To keep a truly perturbative parameter in our expansion we introduced in Eq. (12) the parameter $`\lambda `$, whose interpretation will be clear later. Our goal will be to solve Eqs. (12)-(13) for arbitrary values of $`\lambda `$, as a perturbative expansion in powers of $`\lambda \mathrm{\Delta }`$. For that we have to expand the unknown matrices $`S`$,$`W`$ in powers of $`\sqrt{\lambda \mathrm{\Delta }}`$. We see that the equations are consistent with $`W`$ carrying odd powers and $`S`$ even powers:
$`W`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\lambda \mathrm{\Delta }}}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}(\lambda \mathrm{\Delta })^kW^{(k)}`$ (21)
$`S`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}(\lambda \mathrm{\Delta })^kS^{(k)}.`$ (22)
In the following paragraphs we will show that it is possible to solve the set of equations (12)-(13) order by order in $`\sqrt{\lambda \mathrm{\Delta }}`$. Finally, setting $`\lambda =1`$ one recovers the solution of the self-duality condition. However, even for this value higher orders in the expansion are suppressed by powers of $`\sqrt{\mathrm{\Delta }}`$. Thus, as we will verify later, we expect the first few terms of the expansion to approximate the self-dual solutions for small values of $`\mathrm{\Delta }`$.
On the other hand, the solution for arbitrary $`\lambda `$ can be interpreted as the solution of the following modified self-duality equation:
$$𝐅_{\mu \nu }(x)\stackrel{~}{𝐅}_{\mu \nu }(x)=(1\lambda )(G_{\mu \nu }\stackrel{~}{G}_{\mu \nu })\tau _3.$$
(23)
For $`\lambda =0`$ the constant field strength configuration is a solution, while for $`\lambda =1`$ we recover the self-duality equations.
Now let us address solving the equations order by order in $`\lambda `$. Notice first that $`W^{(1)}`$ satisfies the equation:
$$\overline{D}W^{(1)}=0.$$
(24)
This equation has non-zero regular solutions, as can be deduced from the index theorem. As we will see later the general solution has the form:
$$W^{(1)}=\mathrm{\Psi }(x)\left(\begin{array}{cc}K^{(1)}& Q^{(1)}\\ 0& 0\end{array}\right),$$
(25)
where $`K^{(1)}`$ and $`Q^{(1)}`$ are two arbitrary complex numbers and $`\mathrm{\Psi }(x)`$ is a function whose explicit form will be given later. Having seen that Eq. (13) has a solution for $`k=1`$, let us address the question of whether it also has a solution for all values of $`k`$. For that purpose one has to investigate the adjoint of $`\overline{D}`$, and see whether its kernel is null or not. Indeed, it is easy to see that the kernel of the adjoint vanishes, and hence Eq. (13) has a regular solution no matter what the left-hand side is, provided it is regular. Explicitly, one such solution is given by:
$`W^{(k)}=DU^{(k)}`$ (26)
where $`\overline{D}DU^{(k)}=ı{\displaystyle \underset{l=1}{\overset{k1}{}}}\left(S^{(kl)}W^{(l)}W_c^{(l)}S^{(kl)}\right).`$ (27)
The existence of a unique solution $`U^{(k)}`$ to Eq. (27) follows from the invertibility of the operator $`\overline{D}D`$. In terms of this particular solution, and using (24)-(25), the most general solution of Eq. (13) is given by:
$$W^{(k)}=DU^{(k)}+\mathrm{\Psi }(x)\left(\begin{array}{cc}K^{(k)}& Q^{(k)}\\ 0& 0\end{array}\right),$$
(28)
where $`K^{(k)}`$ and $`Q^{(k)}`$ are complex constants.
Now let us study the solution of equation (12) order by order in $`\lambda `$. It is easy to see that both the left as the right hand sides are periodic functions in the box. They can hence be expanded in Fourier series. The solution can be obtained by equating the corresponding Fourier coefficients. However, notice that the left-hand side has no constant term. Hence, if the right-hand has a non-zero constant Fourier term, the equation has no solution. This can be expressed more formally by saying that the kernel of $``$ is nontrivial. What we will now show is that in solving the equation for $`S^{(k)}`$ it is possible to fix the constants $`K^{(k)}`$ and $`Q^{(k)}`$ appearing in the solution to (13) to order $`k`$, by the condition that the lowest Fourier component of the right-hand side vanishes. For that we have to explore the effect of the replacement (28) in Eq. (12) to order $`(\lambda \mathrm{\Delta })^k`$. The dependence of the right-hand side of Eq. (12) on $`K^{(k)}`$ and $`Q^{(k)}`$ to this order, is contained in the following term:
$$|\mathrm{\Psi }(x)|^2(c_3\tau _3+c_+\tau _++c_+^{}\tau _{}),$$
(29)
where $`c_3=2\mathrm{}(Q^{(k)}Q^{(1)}K^{(k)}K^{(1)})`$ and $`c_+=2(K^{(k)}Q^{(1)}+K^{(1)}Q^{(k)})`$. The symbol $`\mathrm{}`$ denotes the real part of its complex argument. By choosing $`K^{(k)}`$, $`Q^{(k)}`$ appropriately the term within parenthesis can be made equal to an arbitrary hermitian, traceless $`2\times 2`$ matrix. Since the lowest Fourier coefficient (the constant one) of $`|\mathrm{\Psi }(x)|^2`$ is non-zero, the constants can be chosen such that the whole right hand side of Eq. (12) has vanishing constant Fourier term. Actually, this fixes 3 of the 4 real parameters which enter $`K^{(k)}`$, $`Q^{(k)}`$. The remaining one corresponds to the symmetry associated to global colour rotations in the $`12`$ plane, which leave $`B_\mu `$ invariant (we will comment upon this property on the next paragraph).
Having shown that the solution of our set of equations exists order by order in our expansion in $`\sqrt{\lambda \mathrm{\Delta }}`$, we have now to analyse uniqueness. Indeed, on general grounds we know that the solution is non-unique. This fact is associated to the existence of transformations which change one solution into other. We already mentioned one: global gauge transformations of a certain kind. These are the residual gauge transformations that are not gauge fixed by Eq. (9). They are associated with the freedom to multiply any solution matrix $`W`$ by a constant phase. The other transformations are space-time translations (followed by an appropriate gauge transformation to preserve the gauge fixing condition). This latter symmetry manifests itself under the form of a non-uniqueness for the solutions of Eq. (12) to any order $`k`$: notice that we are free to add an arbitrary constant matrix to $`S`$ in the left-hand side of the equation, which would entail the four real parameters associated to a translation.
The best strategy in solving the equations is to fix a unique solution by constraining these transformations. This we will do by imposing the following additional conditions:
$`\mathrm{}(W_{12}(x=0))=|W_{12}(x=0)|`$ (30)
$`{\displaystyle 𝑑xS(x)}=0,`$ (31)
(If $`W_{12}(x=0)=0`$ we take $`\mathrm{}(W_{11}(x=0))=|W_{11}(x=0)|`$). It is now completely clear that the procedure leads to a unique solution order by order in $`\sqrt{\lambda }`$.
It is useful to derive expressions for the field strength tensor itself. Just as we did for the vector potential we might expand in colour components:
$$𝐅_{\mu \nu }(x)=F_{\mu \nu }^{(3)}(x)\tau _3+F_{\mu \nu }^{(+)}(x)\tau _++F_{\mu \nu }^{(+)}(x)\tau _{}.$$
(32)
Now since the field is self-dual, we might contract it with the matrices $`\sigma _\mu \overline{\sigma }_\nu `$ to obtain a traceless hermitian matrix combining the three spatial directions:
$`^{(+)}`$ $``$ $`{\displaystyle \frac{ı}{4}}F_{\mu \nu }^{(+)}(x)\sigma _\mu \overline{\sigma }_\nu `$ (33)
$`=`$ $`{\displaystyle \frac{ı}{2}}DW_c^{}{\displaystyle \frac{1}{2}}(SW_c^{}WS^{})`$
$`^{(3)}`$ $``$ $`{\displaystyle \frac{ı}{4}}F_{\mu \nu }^{(3)}(x)\sigma _\mu \overline{\sigma }_\nu `$ (34)
$`=`$ $`{\displaystyle \frac{\pi }{2}}\left({\displaystyle \frac{l_0l_3+l_1l_2}{V}}\right)\tau _3{\displaystyle \frac{ı}{2}}S^{}{\displaystyle \frac{1}{4}}(WW^{}W_cW_c^{}).`$
Having set up the full procedure for calculating the potentials and fields in powers of $`\lambda `$, let us now exemplify it by computing the first terms in this expansion. These results will be used in the next section. The starting point is the equation for $`W^{(1)}`$ (Eq. (24)). We mentioned previously what is the form of the solution. Let us for the moment skip the proof and also the determination of $`\mathrm{\Psi }(x)`$ and proceed. The next step is to look at the equation for $`S^{(1)}`$. As mentioned in the general case, both the left and right-hand sides can be expanded in Fourier coefficients. The condition that the constant coefficient of the right hand side vanishes imposes a constraint on $`K^{(1)}`$ and $`Q^{(1)}`$:
$`Q^{(1)}K^{(1)}=0`$ (35)
$`|Q^{(1)}|^2|K^{(1)}|^2={\displaystyle \frac{2\pi }{\sqrt{V}}},`$ (36)
where we have fixed the normalisation of $`\mathrm{\Psi }(x)`$, such that its constant Fourier coefficient is equal to one. The previous equations lead to
$$K^{(1)}=0;Q^{(1)}=\frac{\sqrt{2\pi }}{V^{\frac{1}{4}}},$$
(37)
where we have used (30).
Now the equation for $`S^{(1)}`$ reads:
$$\overline{}S^{(1)}=\frac{ı\pi }{\sqrt{V}}(|\mathrm{\Psi }(x)|^21)\tau _3.$$
(38)
This can be solved together with Eq. (31) to give:
$$S^{(1)}=\frac{ı\pi }{\sqrt{V}}(h)\tau _3,$$
(39)
where $`h(x)`$ is a periodic function on the box, solution of the equation:
$$\mathrm{}h(x)=|\mathrm{\Psi }(x)|^21$$
(40)
and $`\mathrm{}`$ is the 4-dimensional Laplacian. The previous equation can be solved by expanding both sides in Fourier series and equating.
Let us now work out the details of the solution to Eq. (24). For future purposes we will consider a more general equation:
$$\overline{D_q}\phi (x)=0,$$
(41)
where $`\phi `$ is a two component vector. The operator $`\overline{D_q}`$ is given by:
$$\overline{D_q}=\overline{\sigma }_\mu (_\mu +ı\pi q\frac{x_\nu n_{\mu \nu }}{l_\mu l_\nu }),$$
(42)
where $`q`$ is a constant. The vector of functions $`\phi (x)`$ is required to satisfy the boundary condition:
$$\phi (x+e_\mu )=\mathrm{exp}\{ı\pi qn_{\mu \nu }\frac{x_\nu }{l_\nu }\}\phi (x)$$
(43)
which is only consistent for integer $`q`$. It is easy to see that the operator $`\overline{D_q}`$ preserves this boundary condition.
Now it is seen that locally a solution of Eq. (41) takes the form:
$$\left(\begin{array}{c}\stackrel{~}{\phi }_q(x)\kappa _+(u_0,u_1)\\ (\stackrel{~}{\phi }_q(x))^1\kappa _{}(u_0^{},u_1^{})\end{array}\right),$$
(44)
where we have introduced complex coordinates:
$$u_\mu =\frac{1}{l_\mu }(x_\mu +ın_{\mu \nu }x_\nu )$$
(45)
and $`u_\mu ^{}`$ are the complex conjugates. These coordinates are not independent, and satisfy:
$$u_\mu =\frac{ı}{l_\mu }n_{\mu \nu }u_\nu l_\nu .$$
(46)
We might for future benefit introduce the complex constants:
$$\tau _\mu =\frac{ı}{l_\mu }|n_{\mu \nu }l_\nu |.$$
(47)
The function $`\stackrel{~}{\phi }_q`$ is given by:
$$\stackrel{~}{\phi }_q(x)=\mathrm{exp}\{\frac{\pi q}{2l_0l_3}(x_0^2+x_3^2)\frac{\pi q}{2l_1l_2}(x_1^2+x_2^2)\}.$$
(48)
The boundary conditions Eq. (43) impose constraints on the value on the holomorphic and anti-holomorphic functions $`\kappa _\pm `$:
$`\kappa _+(x+e_\mu )=\mathrm{exp}\{\pi qı{\displaystyle \frac{(u_\mu +\frac{1}{2})}{\tau _\mu }}\}\kappa _+(x)`$ (49)
$`\kappa _{}(x+e_\mu )=\mathrm{exp}\{\pi qı{\displaystyle \frac{(u_\mu ^{}+\frac{1}{2})}{\tau _\mu }}\}\kappa _{}(x).`$ (50)
Choosing $`u_0`$ and $`u_1`$ as our two independent complex variables, we might write:
$`\kappa _+(u_0,u_1)=\mathrm{exp}\{\pi qı({\displaystyle \frac{u_0^2}{2\tau _0}}+{\displaystyle \frac{u_1^2}{2\tau _1}})\}\stackrel{~}{\kappa }_+(u_0,u_1)`$ (51)
$`\kappa _{}(u_0^{},u_1^{})=\mathrm{exp}\{\pi qı({\displaystyle \frac{u_0^2}{2\tau _0}}+{\displaystyle \frac{u_1^2}{2\tau _1}})\}\stackrel{~}{\kappa }_{}(u_0^{},u_1^{})`$ (52)
The functions $`\stackrel{~}{\kappa }_\pm `$ are periodic in their arguments with period $`1`$ and satisfy:
$`\stackrel{~}{\kappa }_+(u_0+\tau _0,u_1)=\mathrm{exp}\{2\pi qıu_0\pi qı\tau _0\}\stackrel{~}{\kappa }_+(u_0,u_1)`$ (53)
$`\stackrel{~}{\kappa }_+(u_0,u_1+\tau _1)=\mathrm{exp}\{2\pi qıu_1\pi qı\tau _1\}\stackrel{~}{\kappa }_+(u_0,u_1).`$ (54)
For $`q=1`$ these are precisely the conditions satisfied by the Riemann $`\theta `$ function . Actually, up to a multiplicative constant, this function is the only (regular) holomorphic function satisfying these boundary conditions. Similarly, one obtains that the equation for $`\stackrel{~}{\kappa }_{}`$ has no regular solutions. Hence, for the $`q=1`$ case we have arrived at the solution given in Eq. (25), and determined the expression for the function $`\mathrm{\Psi }(x)`$:
$$\mathrm{\Psi }(x)=\sqrt{\frac{4l_3l_2}{l_0l_1}}\mathrm{exp}\{\frac{\pi }{l_0l_3}(x_3^2ıx_3x_0)\frac{\pi }{l_1l_2}(x_2^2ıx_2x_1)\}\theta (u_0,\tau _0)\theta (u_1,\tau _1)$$
(55)
The multiplicative factor preceding the right hand side of the previous expression is determined by the condition that the lowest Fourier coefficient of $`|\mathrm{\Psi }(x)|^2`$ is unity.
## 3 Comparison with numerical results
In the previous section we have set up a general procedure to construct the form of the gauge potentials and field strengths for $`SU(2)`$ self-dual solutions on the torus with twist tensor $`n_{03}=n_{12}=1`$. The result is an expansion in powers of $`\sqrt{\mathrm{\Delta }}`$, where $`\mathrm{\Delta }`$ is defined in Eq. (4). However, we do not have an analytical estimate of the size of the coefficient. Our purpose in this section is to test the rate of convergence of the expansion by comparing the results obtained from the first non-trivial order with the exact result as obtained by numerical methods on the lattice.
We will restrict to the analysis of the gauge invariant traces $`\text{Tr}(𝐄_i𝐄_j)`$, where $`𝐄_i=𝐅_{0i}`$ are the electric fields. To lowest order in our expansion, we fall into the constant field strength configuration, and the only non-zero gauge invariant trace is $`\text{Tr}(𝐄_3^2)`$. The next correction is order $`\mathrm{\Delta }`$ and vanishes for $`\text{Tr}(𝐄_1𝐄_2)`$. It also predicts that $`\text{Tr}(𝐄_1^2)=\text{Tr}(𝐄_2^2)`$. Making use of the general formulas (33), (34) and substituting the explicit form of $`S^{(1)}`$ and $`W^{(0)}`$ (Eqs. (39), (25)) one arrives at the following result, valid to order $`\mathrm{\Delta }`$:
$`\text{Tr}(𝐄_1^2(x))`$ $`=`$ $`\text{Tr}(𝐄_2^2(x))=\mathrm{\Delta }{\displaystyle \frac{\pi }{\sqrt{V}}}\left|D_0\mathrm{\Psi }(x)\right|^2,`$ (56)
$`\text{Tr}(𝐄_3^2(x))`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{2}}\left({\displaystyle \frac{l_0l_3+l_1l_2}{V}}\right)^2\times `$ (57)
$`\left\{1\mathrm{\Delta }\left({\displaystyle \frac{2\sqrt{V}}{l_0l_3+l_1l_2}}\right)\left(1+2(_0^2+_3^2)h(x)\right)\right\},`$
$`\text{Tr}(𝐄_1(x)𝐄_3(x))`$ $`=`$ $`\mathrm{\Delta }\pi ^2\left({\displaystyle \frac{l_0l_3+l_1l_2}{V^{3/2}}}\right)\left(_0_2+_1_3\right)h(x),`$ (58)
$`\text{Tr}(𝐄_2(x)𝐄_3(x))`$ $`=`$ $`\mathrm{\Delta }\pi ^2\left({\displaystyle \frac{l_0l_3+l_1l_2}{V^{3/2}}}\right)\left(_0_1_2_3\right)h(x),`$ (59)
where $`\mathrm{\Psi }`$ and $`h`$ have been defined in Section 2 (Eqs. (55) and (40)). Using the standard representation of Riemann’s theta function :
$$\theta (u,\tau )=\underset{n𝐙}{}\mathrm{exp}\{2\pi ınu+\pi ın^2\tau \}$$
(60)
one can easily obtain the Fourier coefficients of all the functions appearing in Eqs. (56-59). As for the numerical comparison, summing the first few hundred terms of the Fourier expansion allows to compute these functions with negligible errors. It is also extremely simple to use these Fourier coefficients to integrate analytically over some of the four real coordinates, to arrive at a quantity better suited for graphically displaying the comparison.
A numerical approximation to the exact solutions of the self-duality equations, with which the results coming from the perturbative approximation that we have developed are to be compared, can be constructed by means of standard lattice gauge theory techniques . We will use for this purpose an $`\epsilon =0`$ overimproved cooling procedure, that was found in previous works to be able to produce very accurate approximants to continuum self-dual fields. In particular, it allows to extract the exact values of the gauge invariant densities $`\text{Tr}(𝐄_i𝐄_j)`$ under concern up to $`𝒪(a^4)`$ corrections, $`a`$ being the lattice spacing (whose precise definition we will discuss below).
To explore the accuracy of the next to leading term in the perturbative expansion with varying values of the remaining parameters, we will consider tori of lengths $`(l_0=l_t(1+ϵ),l_1=l_t,l_2=l_s,l_3=l_s)`$. The results will then depend on the perturbative parameter $`\mathrm{\Delta }=ϵ/\sqrt{1+ϵ}`$, and on the ratio $`l_s/l_t`$ measuring the degree of spatial asymmetry of the torus. Thus, by keeping $`0ϵ1`$ we will remain within the perturbative regime. On the other hand, we will vary the ratio $`l_s/l_t`$ between $`1`$ (the more symmetrical case) and $`0`$, were the $`\theta `$ functions and their derivatives entering into the analytical expressions are well described by polynomials times gaussians.
In each case, a numerical solution is obtained in a lattice with a number of points $`L_\mu =l_\mu /a`$ along direction $`\mu `$. To define the lattice spacing $`a`$, we need to fix a unit. In our case we take $`l_1l_2=1`$. This is justified by noticing from the expression for $`\mathrm{\Psi }`$ in Eq. (55) that the region having nontrivial structure in the action density at nonzero $`\mathrm{\Delta }`$ is of size $`\sqrt{l_1l_2}`$.
We will present the results of the comparison of the analytical results Eqs. (56\- 59) with their numerical counterparts, for three different configurations, having different values of $`\mathrm{\Delta }`$ and $`l_s/l_t`$. The lattice sizes that we will use, together with their associated $`\mathrm{\Delta }`$ and $`l_s/l_t`$ values, are detailed in Table 1.
In Figures 12 and 3 we show the numerical and perturbative results for the integrated electric field densities $`\mathrm{\Phi }_{33}^{(2)}(x_0,x_1)\text{d}x_2\text{d}x_3\text{Tr}(𝐄_3^2(x))`$ and $`\mathrm{\Phi }_{11}^{(2)}(x_0,x_1)=\mathrm{\Phi }_{22}^{(2)}(x_0,x_1)\text{d}x_2\text{d}x_3\text{Tr}(𝐄_1^2(x))`$, and $`\mathrm{\Phi }_{23}^{(2)}(x_0,x_1)\text{d}x_2\text{d}x_3\text{Tr}(𝐄_2(x)𝐄_3(x))`$, respectively (notice that $`\mathrm{\Phi }_{13}^{(2)}`$ vanishes at the present perturbative order, despite the fact that $`\text{Tr}(𝐄_1(x)𝐄_3(x))`$ does not, because of the particular form of the expression for this latter quantity). The qualitative agreement is clearly good. The main features of the exact solution are present in the analytical expression. It is possible to obtain a graphical quantitative measure of the comparison by integrating the previous densities over an additional coordinate to yield the time profiles $`\mathrm{\Phi }_{33}^{(1)}(x_0)\text{d}x_1\text{d}x_2\text{d}x_3\text{Tr}(𝐄_3^2(x))`$ and $`\mathrm{\Phi }_{11}^{(1)}(x_0)=\mathrm{\Phi }_{22}^{(1)}(x_0)\text{d}x_1\text{d}x_2\text{d}x_3\text{Tr}(𝐄_1^2(x))`$ (similarly to what happened with $`\mathrm{\Phi }_{13}^{(2)}`$, $`\mathrm{\Phi }_{23}^{(1)}`$ vanishes despite $`\mathrm{\Phi }_{23}^{(2)}`$ does not). The comparison for these quantities is displayed in Figure 4.
Let us briefly comment some salient features of the solution. As for the density of the component $`𝐄_3`$ of the electric field, a hole appears overlaying the flat background supplied by the zero-order constant abelian field. The width of this structure is, as we mentioned above, proportional to $`\sqrt{l_sl_t}`$, as can be derived from an analysis of the perturbative expression for the potential, and its contribution to the total action is of order $`\mathrm{\Delta }`$ (cf. Eq. (57)). Meanwhile, the action density associated to the other components of the electric field, whose contribution to the total action is as well of order $`\mathrm{\Delta }`$ in the perturbative approach, exhibits in the asymmetric torus case (configurations B and C) a double lump structure, again of size $`\sqrt{l_sl_t}`$, with the maxima aligned in Euclidian time. We see that the perturbative approximation is most accurate for the configuration A, despite the relatively large value of $`\mathrm{\Delta }`$ associated to it. This fact indicates that the convergence behaviour of the perturbative series depends on the value of the asymmetry parameter $`l_s/l_t`$, in such a way that it is worse the more asymmetric the torus is chosen. The overall conclusion is, anyway, that for values of $`\mathrm{\Delta }`$ in the probed range, between $`0.02`$ and $`0.09`$, the NLO perturbative result constitutes a good approximation to the exact solution.
Once the convergence behaviour of the perturbative series for small values of $`\mathrm{\Delta }`$ has been checked to be good, it would be also interesting to study to what extent the NLO result at hand remains useful to describe solutions occurring at larger values of $`\mathrm{\Delta }`$. This possibility is tempting because it would open the door to apply our results to improve the analytical control over some particularly interesting fields. For instance, it is known that in a torus of geometry $`l_t\times l_s^3`$ with $`l_s/l_t1`$ the solutions approach self-dual fields on $`T^3\times R`$, the approximation being already remarkably good for $`l_t3l_s`$. For the considered twist and $`SU(2)`$ gauge group a $`Q=1/2`$ solution is obtained, whose action density displays a single lump exponentially decaying in the large direction of length $`l_t`$, and whose width is controlled by $`l_s`$. In this geometry, and setting $`l_tl_s(1+\delta )`$, one has $`\mathrm{\Delta }=\delta /\sqrt{1+\delta }`$, and the analysis would proceed by moving from the case $`\delta =0`$, where the torus is symmetric and the solution is the abelian one, to values of $`\delta 1`$, where the features of the $`T^3\times R`$ solution would start to arise. Having performed computations on lattices of sizes $`L_t\times L_s^3`$, with $`L_s=12`$ and $`L_t`$ ranging from $`13`$ to $`48`$, we have found that the self-dual configuration evolves smoothly with changing $`\delta `$. Unfortunately, the NLO perturbative approximation begins to deviate substantially from the exact result before the interesting regime $`\delta 1`$ is reached.
In the same spirit one could investigate other torus geometries, e.g. $`l_t^2\times l_s^2`$ with $`l_s/l_t1`$, which in some cases is known to lead to limiting $`T^2\times R^2`$ solutions with a vortex-like structure , or $`l_t\times l_s^3`$ with $`l_s/l_t1`$, which leads to the $`R^3\times S^1`$ caloron solutions . In these cases, we would expect a similar behaviour to that found for the $`T^3\times R`$ case.
## 4 Nahm transform
Nahm’s transformation maps self-dual configurations on the torus onto other self-dual configurations. The modifications necessary to cope with twisted boundary conditions have only been worked out recently . In general, the transformation changes the twist tensor and torus sizes and maps the rank of the group ($`N`$) and the topological charge ($`Q`$) onto each other through the formula:
$`QQ^{}=N/N_0`$ (61)
$`NN^{}=QN_0,`$ (62)
which preserves the dimensionality of the moduli spaces ($`QN=Q^{}N^{}`$). The integer constant $`N_0`$ depends on the twist. This transformation provides an interesting tool for studying (anti-)self-dual gauge fields on the torus.
In previous sections we have expressed certain self-dual potentials as an expansion in the parameter $`\mathrm{\Delta }`$. It is henceforth interesting to analyse the interplay of this result with the Nahm transform. First of all we should find out the general properties of the Nahm transform for the configurations in question. In our case the group is $`SU(2)`$ and the configuration has nontrivial twist tensor $`n_{03}=n_{12}=1`$. This implies that the parameter $`N_0=4`$. Furthermore, the topological charge $`Q`$ of these configurations is determined by the twist matrices and equals $`Q=\frac{1}{2}`$, as for the corresponding constant field strength configuration. Hence, according to the formulas given above, the Nahm dual is again an $`SU(2)`$ solution with topological charge $`Q^{}=\frac{1}{2}`$. Now we can make use of the results of Ref. to determine the twist tensor and torus size of the Nahm transformed field. Indeed, the Nahm dual twist tensor is equal to the original one, and the torus size is given by: $`\frac{1}{2l_0}\times \frac{1}{2l_1}\times \frac{1}{2l_2}\times \frac{1}{2l_3}`$. Thus, except for the different torus size, the Nahm transformed field is of the same type as the original one. Furthermore, the size parameter $`\mathrm{\Delta }^{}`$ of the Nahm transformed field is given by:
$$\mathrm{\Delta }^{}=\mathrm{\Delta }.$$
(63)
Therefore, the Nahm transform provides a nonlinear relation for our perturbative expansion. A full analysis of this point is difficult and lengthy and will be left out from this paper, however it is instructive to look at the first few terms of this connection.
In order to construct the Nahm transform one has to study the zero modes of the Weyl equation in the fundamental representation of the group:
$$(\overline{𝒟}2\pi ı\overline{z})\chi (x;z)=0,$$
(64)
where $`z_\mu `$ represent the coordinates of a point in the Nahm dual torus. The Weyl operator $`\overline{𝒟}`$ contains the self-dual potential Eq. (7) which can be expanded in powers of $`\sqrt{\lambda \mathrm{\Delta }}`$. Although, for $`\lambda 1`$ the configuration is not self-dual, it is still possible to define a Nahm transform (which will not be self-dual). Thus, we can expand $`\chi (x;z)`$ in the same way and equate to zero all of the powers of the equation separately. From Eqs. (21),(22) it is easy to see that the upper and lower components in colour space only mix for odd-even or even-odd powers of the expansion parameter. Thus, we might write:
$$\chi (x;z)=\left(\begin{array}{c}\varphi _++\sqrt{\lambda \mathrm{\Delta }}\varphi _+^{}\\ \varphi _{}+\sqrt{\lambda \mathrm{\Delta }}\varphi _{}^{}\end{array}\right).$$
(65)
In the previous formula, the explicit vector is in colour space, while the quantities $`\varphi _\pm (x;z)`$, $`\varphi _\pm ^{}(x;z)`$ are bi-spinors, which can be expanded in power series in $`\lambda \mathrm{\Delta }`$. Eq. (64) amounts for $`\varphi _+`$ and $`\varphi _{}^{}`$ to the equations:
$`(\overline{D}_{\frac{1}{2}}ıS^{}2\pi ı\overline{z})\varphi _+=ı\sqrt{\lambda \mathrm{\Delta }}W_c^{}\varphi _{}^{}`$ (66)
$`(\overline{D}_{\frac{1}{2}}+ıS^{}2\pi ı\overline{z})\varphi _{}^{}=ı{\displaystyle \frac{1}{\sqrt{\lambda \mathrm{\Delta }}}}W^{}\varphi _+`$ (67)
and a similar equation holds for the remaining components. The symbol $`\overline{D}_{\frac{1}{2}}`$ is defined in (42). We see that in this way we get two independent solutions of (64) as predicted by the index theorem. The Nahm transformed $`SU(2)`$ vector potential is then given by the formula:
$$\widehat{A}_\mu ^{ij}(z)=ıd^4x\chi ^i(x;z)\frac{}{z_\mu }\chi ^j(x;z),$$
(68)
where the indices $`i`$,$`j\{1,2\}`$ label the two linearly independent and orthonormal solutions.
For the whole construction, the question of the boundary conditions satisfied by the spinors is crucial. Indeed, the naive periodicity requirement:
$$\varphi _\pm (x+e_\mu ;z)=\mathrm{exp}\{\pm ı\frac{\pi }{2}n_{\mu \nu }\frac{x_\nu }{l_\nu }\}\varphi _\pm (x;z)$$
(69)
is inconsistent. How to remedy this situation is what is studied in Ref. . In the case at hand the easiest way out is to impose the periodicity requirement only for the $`x_0`$ and $`x_1`$ direction, while requiring only double period conditions on the other two. In short, this is just replicating the torus in the $`x_3`$ and $`x_2`$ directions. Consistently the integration in (68) has to be performed in this larger torus.
To illustrate the procedure we will explicitly work out the lowest order term $`\varphi _+^{(0)}`$, which satisfies:
$$(\overline{D}_{\frac{1}{2}}2\pi ı\overline{z})\varphi _+^{(0)}=0$$
(70)
This is a modification of the general equation studied in the previous chapter for $`q=\frac{1}{2}`$. Hence, following the same steps as before and imposing the new boundary conditions, we arrive at a unique solution (up to a multiplicative constant):
$$\varphi _+^{(0)}(x;z)=\mathrm{exp}\{\pi ız_\mu y_\mu )\}\widehat{\mathrm{\Psi }}(y)(\begin{array}{c}K^{(0)}\\ 0\end{array})$$
(71)
where we have defined the auxiliary variable:
$$y_\mu =x_\mu +2l_\mu l_\nu n_{\mu \nu }z_\nu ,$$
(72)
and the function $`\widehat{\mathrm{\Psi }}`$ is the same one that appears in expression (55), with the replacement of $`l_{2,3}`$ by $`2l_{2,3}`$ and $`\tau _{0,1}`$ by $`2\tau _{0,1}`$. The constant $`K^{(0)}`$ is fixed by the normalisation condition. In the same way one derives for $`\varphi _{}^{(0)}`$ the expression:
$$\varphi _{}^{(0)}(x;z)=i\tau _2\varphi _+^{(0)}(x;z).$$
(73)
Now we might compute $`\widehat{A}_\mu ^{11}(x)`$ by replacing Eq. (71) in (68). Now introducing the complex variables:
$$v_\mu =\frac{1}{l_\mu }(y_\mu +ın_{\mu \nu }y_\nu )$$
(74)
we can express all derivatives with respect to $`z_\mu `$ in terms of derivatives with respect to $`v_0`$, $`v_1`$, $`v_0^{}`$ and $`v_1^{}`$. For example:
$$\frac{}{z_0}\varphi _+^{(0)}=\left(\pi ıy_02ıl_3\left(\frac{}{v_0}\frac{}{v_0^{}}\right)\right)\varphi _+^{(0)}.$$
(75)
$`\widehat{\mathrm{\Psi }}(y)`$ has a very simple dependence on $`v_0^{}`$ and hence one has:
$$\frac{}{v_0^{}}\varphi _+^{(0)}=l_0\left(\frac{\pi }{4l_3}v_0+\frac{ı\pi }{2}\left(z_0+ız_3\right)\right)\varphi _+^{(0)}.$$
(76)
The result of differentiating with respect to $`v_0`$ is much more complicated, involving derivatives of Riemann’s theta function. However, in the expression for $`\widehat{A}`$ it is possible to integrate by parts and make the derivatives with respect to $`v_0`$ act onto the complex conjugate of $`\varphi _+^{(0)}`$, for which the complex conjugate of Eq. (76) allows us to obtain a simple expression. We end up with:
$`(\widehat{A}_0^{(0)})^{11}(z)`$ $`=`$ $`ı{\displaystyle d^4x\varphi _+^{(0)}\left(\pi ıy_0+4ıl_3l_0\mathrm{}\left\{\frac{\pi }{4l_3}v_0+\frac{ı\pi }{2}(z_0+ız_3)\right\}\right)\varphi _+^{(0)}}`$ (77)
$`=`$ $`2\pi l_0l_3z_3.`$
In the right hand side of the first equality we have displayed the contribution of the two terms entering the right hand side of (75). One can compute the other components of the Nahm transformed field in the same fashion arriving at:
$$\widehat{A}_\mu ^{(0)}(z)=B_\mu ^{}(z)$$
(78)
where $`B_\mu ^{}(z)`$ is given by the same expression (5) as in section 2, but with the lengths of the torus $`l_\mu `$, replaced by those of the Nahm dual one $`l_\mu ^{}=\frac{1}{2l_\mu }`$. The previous result implies that the Nahm dual of constant field strength configuration is a constant field strength configuration, even if they are not self-dual, generalising the result of Ref. .
## 5 Conclusions
In the previous sections we have presented a systematic expansion which allows the construction of self-dual $`SU(2)`$ Yang-Mills solutions with twist tensor $`n_{03}=n_{12}=1`$ on the torus. The size of higher order corrections depends on the lengths of the torus. For certain torus sizes the solution becomes equal to the the well-known constant field-strength ones. The magnitude of higher order corrections grows as we move away from these torus sizes. We have also compared the landscape of the solution, as obtained from our analytical expressions to leading non-trivial order in the expansion, to the numerical result obtained through standard techniques. The result is quite satisfactory at a qualitative and quantitative level. Finally, the expansion is used to obtain the Nahm transform of the self-dual configuration. Curiously, the perturbative construction of the Nahm transform has the same structure as the direct perturbative construction of the self-dual configuration on the Nahm-self-dual torus. We have not been able to equate these expansions order by order, but have shown this to be the case to lowest order.
Let us now comment about the usefulness of our programme. Up to a proof, which we do not give, of convergence of our expansion, our method gives a direct proof of the existence of the solutions and for them to have the correct number of degrees of freedom. Even if calculating higher orders of the expansion turns out to be a difficult task, the expansion can be of theoretical interest for different reasons. For example, it might allow to investigate some general properties of the solutions. The case of the interplay with the Nahm transform is interesting, and should be pursued. Furthermore, one can study certain extreme limits of the torus sizes, which might allow to obtain exact solutions.
Finally, we comment that our choice of $`SU(2)`$ and of the aforementioned twist tensor has been dictated by simplicity. There is, however, no a priori essential difficulty in generalising the construction given here to other $`SU(N)`$ groups and different twist tensors.
## Acknowledgements
This work was financed by the CICYT under grant AEN97-1678. We thank the Centro de Computación Científica (UAM) for the use of computing resources.
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# 1 Introduction
## 1 Introduction
It is a long established fact that the negative sign of the one-loop beta function of QCD can be beautifully explained as the result of a competition between the paramagnetic and diamagnetic contributions to the bare coupling constant obtained from integrating out very high-energy degrees of freedom. Indeed, the paramagnetic contribution, which anti-screens the charge, prevails over the diamagnetic contribution, which screens the charge . This dominance can be explicitly computed by evaluating the high-energy quantum fluctuations around a classical field configuration , which in turn can be done in a very elegant way using the background field method .
The vacuum-to-vacuum amplitude in the presence of an $`SU(N)`$ background field $`B_\mu `$ is given by the partition function $`Z[B]`$, which up to one loop, reads in the Feynman-‘t Hooft background field gauge
$$Z[B]=e^{iS_{\mathrm{cl}}[B]}𝒟Q𝒟c𝒟\overline{c}\mathrm{exp}\left\{\frac{i}{2g_0^2}\left[S^{\mathrm{diam}}+S^{\mathrm{param}}\right]\right\},$$
(1.1)
where
$$\begin{array}{c}S^{\mathrm{diam}}=d^Dx\mathrm{Tr}\left\{\left(D_\mu [B]Q_\nu \right)\left(D^\mu [B]Q^\nu \right)2\overline{c}D_\mu [B]D^\mu [B]c\right\}\\ S^{\mathrm{param}}=2d^Dx\mathrm{Tr}\left\{F^{\mu \nu }[B]Q^\sigma (S_{\mu \nu })_{\sigma \rho }Q^\rho \right\}.\end{array}$$
(1.2)
Here $`D[B]_\mu =_\mu i[B_\mu ,]`$ and $`F_{\mu \nu }[B]`$ denote the covariant derivative and the field strength for the background field $`B_\mu `$, the $`SU(N)`$ gauge field $`Q_\mu `$ describes the quantum fluctuations about the background and $`(S_{\mu \nu })_{\sigma \rho }`$ are the generators of the Lorentz group in the spin one representation, i.e.
$$(S_{\mu \nu })_{\sigma \rho }=i\left(\eta _{\mu \sigma }\eta _{\nu \rho }\eta _{\mu \rho }\eta _{\nu \sigma }\right).$$
The high-momentum modes contributing to the path integral $`Z[B]`$ yield a logarithmic UV divergence. This divergence shows in dimensional regularization as a pole in $`\epsilon =0`$, with $`D=4+2\epsilon `$. As a matter of, introducing dimensional regularization and integrating over $`Q_\mu `$, one obtains
$$\begin{array}{c}\mathrm{\Gamma }[B]=\frac{1}{i}\mathrm{ln}Z[B]\hfill \\ =\frac{1}{2}\left(\frac{1}{g_0^2}+\frac{a^{\mathrm{diam}}}{\epsilon }+\frac{a^{\mathrm{param}}}{\epsilon }\right)d^Dx\mathrm{Tr}\left(F^{\mu \nu }[B]F_{\mu \nu }[B]\right)+\mathrm{\Gamma }_{\mathrm{finite}}[B],\hfill \end{array}$$
(1.3)
where<sup>3</sup><sup>3</sup>3In this paper, to emphasize the value of diamagnetic and paramagnetic contributions, we will enclose the relevant factors in square brackets, as in eq. (1.4).
$$a^{\mathrm{diam}}=\frac{N}{16\pi ^2}\left[\frac{1}{3}\right]a^{\mathrm{param}}=\frac{N}{16\pi ^2}[4]$$
(1.4)
and $`\mathrm{\Gamma }_{\mathrm{finite}}[B]`$ collects all finite contributions as $`D4`$. The origin of the coefficients $`a^{\mathrm{diam}}`$ and $`a^{\mathrm{param}}`$ can be explained as follows. The orbital motion of the charged quanta with only two polarizations $`(D2`$ in $`D`$ space-time dimensions) of the gluon field $`Q_\mu `$ in the background $`B_\mu `$ yields “diamagnetism”. This motion is described by $`S^{\mathrm{diam}}`$ above and its contribution to the r.h.s of eq. (1.3) has a singular part coming from high-energy quanta which is given by $`a^{\mathrm{diam}}/\epsilon `$. On the other hand, $`S^{\mathrm{paran}}`$ involves the spin current $`Q^\sigma (S_{\mu \nu })_{\sigma \rho }Q^\rho `$ and describes the coupling between the spin of the gluon field $`Q_\mu `$ and the background $`B_\mu `$. This coupling gives rise to “paramagnetism” and the contribution to $`\mathrm{\Gamma }[B]`$ of the high-momentum quanta involved in it is $`a^{\mathrm{param}}/\epsilon `$. One next defines the renormalized coupling constant or renormalized charge $`g(\mu )`$ in the MS scheme at one loop as usual:
$$\frac{1}{g_0^2}+\frac{a^{\mathrm{diam}}}{\epsilon }+\frac{a^{\mathrm{param}}}{\epsilon }=\frac{1}{g^2(\mu )\mu ^{2\epsilon }}.$$
(1.5)
Hence, the one-loop beta function $`\beta (g^2)`$ reads for $`D=4`$
$$\beta (g^2)=\mu \frac{dg^2}{d\mu }=2g^2(\mu )\left(a^{\mathrm{diam}}+a^{\mathrm{param}}\right)=\frac{N}{16\pi ^2}g^4(\mu )\left[\frac{22}{3}\right].$$
From this equation for the beta function one can draw the following conclusions. First, the coefficient $`a^{\mathrm{diam}}`$ being positive implies that its effect is to make the charge $`g^2(\mu =\mathrm{\Lambda })`$ decrease with the momentum scale $`\mathrm{\Lambda }`$. This is the charge screening effect due to the orbital motion of the two physical polarizations of the quanta $`Q_\mu `$ in the field $`B_\mu `$. Secondly, since $`a^{\mathrm{param}}`$ is negative, its effect is to make $`g^2(\mu =\mathrm{\Lambda })`$ grow as $`\mathrm{\Lambda }`$ decreases. This is called anti-screening of the charge and it is due to the interaction of the spin with the field $`B_\mu `$. And thirdly, the inequalities $`|a^{\mathrm{param}}|>a^{\mathrm{diam}}`$ and $`a^{\mathrm{param}}<0`$ explain quantitatively the negative sign of the beta function of the theory, hence, that the charge goes to zero as $`\mathrm{\Lambda }`$ goes to infinity (asymptotic freedom).
It is already a year since it has been shown that $`U(1)`$ gauge theory on non-commutative $`\mathrm{I}\mathrm{R}^4`$ has an UV divergent behaviour at one loop very similar to that of conventional Yang-Mills theory \[see for the $`U(N)`$ case\]. In fact, the one-loop beta function is also negative , which leads to asymptotic freedom. By contrast, as discovered in ref. , the IR behaviour of non-commutative $`U(1)`$ gauge theory presents completely novel features . Indeed, the renormalization of UV divergences induce IR divergences, thus yielding a relation between UV and IR divergences which has been interpreted as a sort of UV/IR duality. This duality seems not to be an artifact of perturbation theory , since it has been re-obtained by defining the field theory as the infinite tension limit of the appropriate open bosonic string theory on a magnetic $`B`$-field .
The purpose of this paper is two-fold. First, to investigate if, in analogy with conventional Yang-Mills theory, the one-loop beta function of $`U(1)`$ gauge theory on non-commutative Minkowski space can be understood as a dominance of paramagnetism over diamagnetism. And secondly, to study through explicit computations UV/IR mixing from the point of view of paramagnetism versus diamagnetism. The paper is organized as follows. In section 2, we will explicitly compute the two-point function and show that the one-loop beta function of $`U(1)`$ gauge theory on non-commutative Minkowski space has a paramagnetic contribution, producing anti-screening of the charge, and a diamagnetic contribution, giving rise to screening of the charge. We will see that the paramagnetic contribution dominates, thus providing a negative beta function. Furthermore, we will take advantage of the computations in ref. to show that these paramagnetic and diamagnetic contributions can be given a stringy interpretation as tachyon magnification and zero mode contributions, in the same sense as for conventional Yang-Mills theory . In sections, 3 and 4 we will calculate the UV divergent terms and the leading non-commutative IR terms of the three and four-point functions. It will turn out that the UV/IR mixing occurs for paramagnetic and diamagnetic logarithmic contributions separately. Section 5 collects or conclusions. We will argue there how our results lead to an IR renormalization of the coupling constant.
## 2 The vacuum polarization tensor
Non-commutative Minkowski space-time is defined by the algebra generated by the operators $`X^\mu (\mu =0,\mathrm{},D1)`$ subject to the commutation relations
$$[X^\mu ,X^\nu ]=i\theta ^{\mu \nu },$$
where $`\theta ^{\mu \nu }`$ is an anti-symmetric real matrix and contraction of indices is performed with the Minkowski metric. We shall take $`\theta ^{\mu \nu }`$ to be “magnetic”, i.e. $`\theta ^{0i}=0`$ for $`i=1,2,3`$, since in this case the field theory exists as the zero slope limit of a string theory in a magnetic background. If $`\theta ^{\mu \nu }`$ is “electric”, the field theory does not exist as the zero slope limit of a string theory and does not make sense on its own since it does not preserve unitarity . Without loss of generality, we will take $`\theta ^{\mu \nu }`$ to vanish for all $`\mu `$ and $`\nu `$, except for $`\mu ,\nu =1,2`$, for which we write
$$\theta ^{12}=\theta ^{21}\theta .$$
The classical action of $`U(1)`$ gauge theory on non-commutative Minkowski space-time is given by
$$S_{\mathrm{class}}[A]=\frac{1}{4g^2}d^Dx\left(F^{\mu \nu }F_{\mu \nu }\right)(x),$$
(2.6)
where
$$\begin{array}{c}F_{\mu \nu }(x)=_\mu A_\nu (x)_\nu A_\mu (x)i[A_\mu ,A_\nu ](x)\\ [A_\mu ,A_\nu ]=(A_\mu A_\nu A_\nu A_\mu )(x)\end{array}$$
is the field strength and the symbol $``$ stands for the Moyal product of functions
$$\left(fg\right)(x)=f(x)e^{\frac{i}{2}\theta ^{\mu \nu }\stackrel{}{_\mu }\stackrel{}{_\nu }}g(x).$$
The classical theory is invariant under non-commutative $`U(1)`$ gauge transformations, which in infinitesimal form have the form
$$\delta _\omega A_\mu (x)=D_\mu [A]\omega =_\mu \omega (x)i[A_\mu ,\omega ](x),$$
the commutator being defined with regard to the Moyal product.
To quantize the theory around a background field configuration, say $`B_\mu (x)`$, we shall use the background field method in the Feynman-‘t Hooft gauge. To this end, we write the gauge field $`A_\mu `$ as the sum of the the background $`B_\mu `$ and the quantum fluctuation $`Q_\mu `$ about it, $`A_\mu =B_\mu +Q_\mu `$. The tree-level action $`S`$ then becomes
$$S=S_{\mathrm{class}}[B+Q]+S_{\mathrm{gf}}+S_{\mathrm{gh}},$$
(2.7)
where the gauge fixing and ghost terms have the form
$$\begin{array}{c}S_{\mathrm{gf}}=\frac{1}{2}d^DxD_\mu [B]Q^\mu D_\nu [B]Q^\nu \\ S_{\mathrm{gh}}=d^Dx\overline{c}D_\mu [B+Q]D^\mu [B]c.\end{array}$$
The fields $`c`$ and $`\overline{c}`$ are the ghost fields and the commutators in the covariant derivatives
$$D_\mu [B+Q]=_\mu i[B_\mu +Q_\mu ,]D_\mu (B)=_\mu i[B_\mu ,]$$
are defined with regard to the Moyal product. Some straightforward manipulations give for the partition function $`Z[B]`$ of the theory up to one loop in the $`U(1)`$ background field $`B_\mu `$ the expression
$$Z[B]=e^{iS_{\mathrm{cl}}[B]}𝒟Q𝒟c𝒟\overline{c}\mathrm{exp}\left[\frac{i}{2g_0^2}\left(S^{\mathrm{diam}}+S^{\mathrm{param}}\right)\right],$$
(2.8)
with
$$\begin{array}{c}S^{\mathrm{diam}}=d^Dx\left(D_\mu [B]Q_\nu D^\mu [B]Q^\nu 2\overline{c}D_\mu [B]D^\mu [B]c\right)\\ S^{\mathrm{param}}=2d^DxF^{\mu \nu }[B]Q^\sigma (S_{\mu \nu })_{\sigma \rho }Q^\rho ,\end{array}$$
(2.9)
in complete analogy with eq. (1.1). The term $`S^{\mathrm{param}}`$ involves the spin one non-commutative current $`Q^\sigma (S_{\mu \nu })_{\sigma \rho }Q^\rho `$ and describes the coupling of the spin to the background field $`B_\mu `$. This term thus gives rise to non-commutative Pauli “paramagnetism”. In turn, the functional $`S^{\mathrm{diam}}`$ is the classical action for the motion in the field $`B_\mu `$ of the $`D2`$ physical degrees of freedom of the field $`Q_\mu `$. Indeed,
$$𝒟Q𝒟c𝒟\overline{c}\mathrm{exp}\left[\frac{i}{2g_0^2}S^{\mathrm{diam}}\right]=\left(\mathrm{det}^{1/2}D^2[B]\right)^{D2}.$$
(2.10)
We shall then say that $`S^{\mathrm{diam}}`$ gives rise to non-commutative Landau “diamagnetism”. From $`S^{\mathrm{diam}}`$ and $`S^{\mathrm{param}}`$ one readily extracts the Feynman rules needed for one-loop perturbative computations. We have collected them in fig. 1, where we have used the notation
$$qp=\theta ^{\mu \nu }q_\mu p_\nu .$$
Vertices coming from $`S^{\mathrm{diam}}`$ and $`S^{\mathrm{param}}`$ will be called diamagnetic and paramagnetic respectively.
In this section we compute up to one loop the vacuum polarization tensor $`\mathrm{\Pi }_{\mu \nu }(p)`$, defined as
$$i\mathrm{\Pi }_{\mu \nu }(x,y)=\frac{\delta ^2\mathrm{\Gamma }[B]}{\delta B_\mu (x)\delta B_\nu (y)}|_{B=0}=\frac{d^Dp}{(2\pi )^D}e^{ip(xy)}i\mathrm{\Pi }_{\mu \nu }(p),$$
where $`i\mathrm{\Gamma }[B]=\mathrm{ln}Z[B]`$. According to the nature of their vertices, the one-loop Feynman diagrams contributing to $`\mathrm{\Pi }_{\mu \nu }(p)`$ fall into three categories: diagrams with only diamagnetic vertices, diagrams with only paramagnetic vertices and diagrams with both diamagnetic and paramagnetic vertices. The contributions to the vacuum polarization tensor coming from these three categories will be denoted by $`\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam}}(p)`$, $`\mathrm{\Pi }_{\mu \nu }^{\mathrm{param}}(p)`$ and $`\mathrm{\Pi }_{\mu \nu }^{\mathrm{mixed}}(p)`$, so we write
$$i\mathrm{\Pi }_{\mu \nu }(p)=i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam}}(p)+i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param}}(p)+i\mathrm{\Pi }_{\mu \nu }^{\mathrm{mixed}}(p).$$
(2.11)
It is very easy to see that there is only one diagram contributing to $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{mixed}}(p)`$, namely that depicted in fig 2. Upon performing some algebra in the integrand, the corresponding Feynman integral takes the form
$$2i\frac{d^Dq}{(2\pi )^D}(\eta _{\mu \nu }\eta _{\mu \nu })\mathrm{sin}^2\left(\frac{qp}{2}\right)\frac{(2q+p)_\rho }{q^2(q+p)^2}=0.$$
Hence
$$\mathrm{\Pi }_{\mu \nu }^{\mathrm{mixed}}(p)=0.$$
(2.12)
We are thus left with the diamagnetic and paramagnetic contributions.
The one-loop contribution to $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam}}(p)`$ is given by the sum of the diagrams in fig. 3, which using the Feynman rules reads
$$i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam}}(p)=[D2]\frac{d^Dq}{(2\pi )^D}\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}\mathrm{sin}^2\left(\frac{qp}{2}\right)\left[\frac{(p+2q)_\mu (p+2q)_\nu }{q^2(p+q)^2}2\frac{\eta _{\mu \nu }}{q^2}\right].$$
(2.13)
The $`D`$ in the factor $`[D2]`$ in front of the integral comes from the diagrams with photons flowing around the loop, whereas the diagrams with ghost loops yield the contribution $`2`$. Hence the effective one-loop contribution to $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam}}(p)`$ corresponds to $`D2`$ scalar fields (transforming in the adjoint representation of the gauge group) moving in the field $`B_\mu `$, in agreement with the discussion following eqs. (2.8) and (2.9) above. To keep to a minimum the number of diagrams to draw, we have not considered planar and non-planar diagrams separately. In fact, each of our diagrams is the sum of a planar and a non-planar contribution. To disentangle these contributions from one another, it is enough to use the identity $`\mathrm{\hspace{0.17em}4}\mathrm{sin}^2\left(qp/2\right)=2e^{iqp/2}e^{iqp/2}`$. This identity gives, upon substitution in (2.13), a $`\theta ^{\mu \nu }`$-independent integral, which defines the planar contribution $`i\mathrm{\Pi }^{\mathrm{diam},\mathrm{P}}(p)`$, and a $`\theta ^{\mu \nu }`$-dependent integral, which constitutes the non-planar contribution $`i\mathrm{\Pi }^{\mathrm{diam},\mathrm{NP}}(p)`$:
$$i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam}}(p)=i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam},\mathrm{P}}(p)+i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam},\mathrm{NP}}(p).$$
(2.14)
After introducing Schwinger parameters, performing the momenta integrals in dimensional regularization and integrating by parts to factorize out the transverse tensor $`p^2\eta _{\mu \nu }p_\mu p_\nu `$, we obtain
$$i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam},\mathrm{P}}(p)=i\frac{[D2]}{(4\pi )^{D/2}}(p^2\eta _{\mu \nu }p_\mu p_\nu )_0^1𝑑x(12x)^2_0^{\mathrm{}}𝑑tt^{1D/2}e^{p^2tx(1x)}$$
(2.15)
and
$$\begin{array}{cc}\hfill i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam},\mathrm{NP}}(p)=i\frac{[D2]}{(4\pi )^{D/2}}& [(p^2\eta _{\mu \nu }p_\mu p_\nu )_0^1dx(12x)^2_0^{\mathrm{}}dtt^{1D/2}e^{p^2tx(1x)\stackrel{~}{p}^2/4t}\hfill \\ & +\stackrel{~}{p}_\mu \stackrel{~}{p}_\nu _0^1dx_0^{\mathrm{}}dtt^{1D/2}e^{p^2tx(1x)\stackrel{~}{p}^2/4t}],\hfill \end{array}$$
(2.16)
where
$$\stackrel{~}{p}^\mu =\theta ^{\mu \nu }p_\nu \stackrel{~}{p}^2=p_\mu \theta ^{\mu \rho }\theta _\rho {}_{}{}^{\nu }p_{\nu }^{}.$$
It is well known that when a Feynman amplitude is expressed in terms of Schwinger parameters (the so-called parametric representation) the contribution to the amplitude coming from virtual quanta carrying arbitrarily high momenta is given by the contribution to the corresponding parametric integral coming from regions of the integration domain where all the Schwinger parameters are arbitrarily small. If we apply this reasoning to $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam},\mathrm{P}}(p)`$, we conclude that, for $`D4`$, the non-integrable singularity in eq. (2.15) at $`t=0`$ shows that virtual quanta $`Q_\mu `$ carrying arbitrarily high momenta yield an UV divergent contribution. This divergence needs renormalization. On the other hand, if $`\stackrel{~}{p}^20`$, the integrand of any of the integrals in eq. (2.16) is non-singular at $`t=0`$. Hence the contribution to $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam},\mathrm{NP}}(p)`$ coming from arbitrary high momenta quanta $`Q_\mu `$ is curbed by the non-commutativity of space through the exponential $`\mathrm{exp}(\stackrel{~}{p}^2/4t)`$, acting $`1/\stackrel{~}{p}^2`$ as a regulator. Of course, if the regulator is removed, divergences spring back. This explains the UV origin of the IR divergences that occur in $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam},\mathrm{NP}}(p)`$ at $`\stackrel{~}{p}^2=0`$. Furthermore, as it will come up shortly, the fact that the tree-level vertices vanish as $`\theta ^{\mu \nu }0`$, so that loops formally vanish in this limit, makes the coefficient of the logarithmic UV divergence in $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam},\mathrm{P}}(p)`$ at $`D=4`$ to be opposite to the coefficient of its logarithmic IR divergence at $`\stackrel{~}{p}^2=0`$. All this is the UV/IR mixing at work .
The one-loop paramagnetic contribution $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param}}(p)`$ to the vacuum polarization tensor is given by the contribution involving only two background fields which comes from the diagram in fig. 4. For the diagram itself the Feynman rules yield
$$F_{\mu \nu }(p)F_{\rho \sigma }(p)_{1\mathrm{P}\mathrm{I}}=4(\eta _{\mu \rho }\eta _{\nu \sigma }\eta _{\mu \sigma }\eta _{\nu \rho })J(p),$$
(2.17)
where $`J(p)`$ is the integral
$$J(p)=\frac{d^Dq}{(2\pi )^D}\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}\mathrm{sin}^2\left(\frac{qp}{2}\right)\frac{1}{q^2(p+q)^2}.$$
(2.18)
The contribution $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param}}(p)`$ is then given by
$$i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param}}(p)=[\mathrm{\hspace{0.17em}8}]i(p^2\eta _{\mu \nu }p_\mu p_\nu )J(p),$$
(2.19)
which has a planar and a non-planar part
$$i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param}}(p)=i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param},\mathrm{P}}(p)+i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param},\mathrm{NP}}(p).$$
Using the same arguments as for the diamagnetic part, we obtain
$$\begin{array}{c}i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param},\mathrm{P}}(p)=[\mathrm{\hspace{0.17em}8}]i(p^2\eta _{\mu \nu }p_\mu p_\nu )J^\mathrm{P}(p)\\ J^\mathrm{P}(p)=\frac{1}{(4\pi )^{D/2}}_0^1𝑑x_0^{\mathrm{}}𝑑tt^{1D/2}e^{p^2tx(1x)}\end{array}$$
(2.20)
for the planar part, and
$$\begin{array}{c}i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param},\mathrm{NP}}(p)=[\mathrm{\hspace{0.17em}8}]i(p^2\eta _{\mu \nu }p_\mu p_\nu )J^{\mathrm{NP}}(p)\\ J^{\mathrm{NP}}(p)=\frac{1}{(4\pi )^{D/2}}_0^1𝑑x_0^{\mathrm{}}𝑑tt^{1D/2}e^{p^2tx(1x)\stackrel{~}{p}^2/4t}\end{array}$$
(2.21)
for the non-planar part. A similar analysis as for the diamagnetic contributions is in order. Indeed, the high momenta modes going around the loop of the paramagnetic diagram in fig. 4 give, for $`D4`$, an UV divergent contribution to $`J^\mathrm{P}(p)`$, hence to $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param},\mathrm{P}}(p)`$, which corresponds to the non-integrable singularity at $`t=0`$ in eq. (2.20). The contribution of these modes to $`J^{\mathrm{NP}}(p)`$, hence to $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{NP}}(p)`$ is finite provided $`\stackrel{~}{p}^20`$. We may then say that the non-commutative character of space makes $`1/\stackrel{~}{p}^2`$ to play the rôle of an UV regulator for the high-momentum modes contributing to the non-planar part of $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param}}(p)`$. If we remove the regulator, i.e. take the limit $`1/\stackrel{~}{p}^2\mathrm{}`$ ($`\stackrel{~}{p}^20`$), the divergence is recovered, although this time under the guise of an IR divergence.
We are now ready to compute the contributions to the one-loop beta function of the theory at $`D=4`$ coming from the diamagnetic and paramagnetic parts of the vacuum polarization tensor. Recall that the mixed part vanishes. If we set $`D=4+2\epsilon `$ in eqs. (2.15) and (2.20) and make a Laurent expansion around $`\epsilon =0`$, we obtain
$$\begin{array}{c}i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam},\mathrm{P}}(p)=i(p^2\eta _{\mu \nu }p_\mu p_\nu )\left[\frac{a^{\mathrm{diam}}}{\epsilon }+a^{\mathrm{diam}}\mathrm{ln}\left(\frac{p^2}{4\pi }\right)\frac{5}{72\pi ^2}+O(\epsilon )\right]\\ i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param},\mathrm{P}}(p)=i(p^2\eta _{\mu \nu }p_\mu p_\nu )\left[\frac{a^{\mathrm{param}}}{\epsilon }+a^{\mathrm{param}}\mathrm{ln}\left(\frac{p^2}{4\pi }\right)+\frac{1}{\pi ^2}+O(\epsilon )\right],\end{array}$$
(2.22)
where
$$a^{\mathrm{diam}}=\frac{1}{16\pi ^2}\left[\frac{2}{3}\right]a^{\mathrm{param}}=\frac{1}{16\pi ^2}[8].$$
(2.23)
Note that the coefficients $`a^{\mathrm{diam}}`$ and $`a^{\mathrm{param}}`$ can be obtained from the corresponding coefficients in eq. (1.4) by replacing in the latter $`N`$ with $`2`$. To subtract UV divergences, we work in the MS renormalization scheme and add to the classical action $`S_{\mathrm{class}}[B]`$ the counterterm
$$\delta S[B]=\frac{1}{4\epsilon }\left(a^{\mathrm{diam}}+a^{\mathrm{param}}\right)d^DxF^{\mu \nu }[B]F_{\mu \nu }[B],$$
(2.24)
since its term quadratic in the background field $`B_\mu `$ cancels the UV divergences in eqs. (2.22). In sections 3 and 4 we will see that the UV divergences in the three and four-point function are also subtracted by the counterterm (2.24). Taking now into account that the tree-level part of the vacuum polarization tensor in terms of the bare coupling constant $`g_0^2`$ reads
$$\frac{i}{g_0^2}(p^2\eta _{\mu \nu }p_\mu p_\nu )$$
and using eqs. (2.22), we conclude that up to one loop the renormalized vacuum polarization tensor in the MS scheme takes the form
$$\begin{array}{c}i\mathrm{\Pi }_{\mu \nu }^{\mathrm{ren}}(p)=i(p^2\eta _{\mu \nu }p_\mu p_\nu )\left[\frac{1}{g^2}+\left(a^{\mathrm{diamag}}+a^{\mathrm{param}}\right)\mathrm{ln}\left(\frac{p^2}{4\pi \mu ^2}\right)+\frac{67}{72\pi ^2}\right]\hfill \\ +i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam},\mathrm{NP}}(p)+i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param},\mathrm{NP}}(p),\hfill \end{array}$$
(2.25)
where the non-planar contributions $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam},\mathrm{NP}}(p)`$ and $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param},\mathrm{NP}}(p)`$ are obtained by setting $`D=4`$ in eqs. (2.16) and (2.21) respectively. The letter $`g`$ in eq. (2.25) is the renormalized coupling constant in the MS scheme, given by eq. (1.5) for $`a^{\mathrm{diam}}`$ and $`a^{\mathrm{param}}`$ as in eq. (2.23). The beta function of the theory then reads
$$\beta (g^2)=\mu \frac{dg^2}{d\mu }=2g^2(\mu )\left(a^{\mathrm{diam}}+a^{\mathrm{param}}\right)=\frac{1}{4\pi ^2}g^4(\mu )\left[\frac{11}{3}\right].$$
(2.26)
We see that the negative sign of the beta function comes about because the high-momentum paramagnetic contributions, which yield the coefficient $`a^{\mathrm{param}}`$, overcome the high-momentum diamagnetic contributions, which originate the coefficient $`a^{\mathrm{diam}}`$. The analogy with $`SU(N)`$ theory on commutative Minkowski space-time is clear.
Let us next show that the dominance at high energies of paramagnetic contributions over diamagnetic ones has, due to the UV/IR mixing, a cut effect on the IR behaviour<sup>4</sup><sup>4</sup>4The IR divergences that occur when certain linear combinations $`l_i(p)`$ of the external momenta $`p_j`$ satisfy $`\stackrel{~}{l_i}=0`$, with $`l_i0`$, will be called non-commutative IR divergences throughout this paper. The behaviour of the Green functions in the neibourhood of those momentum configurations will be called non-commutative IR behaviour. of the theory at $`\stackrel{~}{p}=0`$. We start by noting that the divergent IR behaviour at $`\stackrel{~}{p}=0`$ of the renormalized vacuum polarization tensor in eq. (2.25) is only caused by $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam})\mathrm{NP}}(p)`$ and $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param})\mathrm{NP}}(p)`$. The $`\stackrel{~}{p}=0`$ divergent terms of the latter are easily computed from eqs. (2.16) and (2.21). For them we obtain
$$\begin{array}{c}i\mathrm{\Pi }_{\mu \nu }^{\mathrm{diam},\mathrm{NP}}(p)ia^{\mathrm{diam}}\mathrm{ln}(p^2\stackrel{~}{p}^2)(p^2\eta _{\mu \nu }p_\mu p_\nu )+\frac{2i}{\pi ^2}\frac{\stackrel{~}{p}_\mu \stackrel{~}{p}_\nu }{(\stackrel{~}{p}^2)^2}\\ i\mathrm{\Pi }_{\mu \nu }^{\mathrm{param},\mathrm{NP}}(p)ia^{\mathrm{param}}\mathrm{ln}(p^2\stackrel{~}{p}^2)(p^2\eta _{\mu \nu }p_\mu p_\nu ),\end{array}$$
(2.27)
where $``$ means that we have dropped any term which is finite at $`D=4`$ as $`\stackrel{~}{p}0`$. It is plain that the IR divergent logarithmic contributions in eqs. (2.27) are dual to the UV divergent logarithmic terms in eqs. (2.22), in the sense that in the momenta region
$$|\stackrel{~}{p}|\theta \mathrm{\Lambda }_{\mathrm{IR}}$$
(2.28)
the former can be obtained from the latter by using the replacements
$$\frac{1}{\epsilon }\mathrm{ln}\mathrm{\Lambda }_{\mathrm{UV}}^2\mathrm{\Lambda }_{\mathrm{UV}}\frac{1}{\theta \mathrm{\Lambda }_{\mathrm{IR}}}.$$
(2.29)
By contrast, the quadratic IR divergent term $`\stackrel{~}{p}_\mu \stackrel{~}{p}_\nu /(\stackrel{~}{p}^2)^2`$ in eq. (2.27) has no dual counterpart as a singular UV contribution in eq. (2.22). This is due to the fact that, if gauge invariance is preserved, no local and quadratic UV divergent contribution occurs in the vacuum polarization tensor. Yet, as can be seen from eq. (2.16), the origin of both logarithmic and quadratic divergences at $`\stackrel{~}{p}=0`$ is the same: the curbing by the non-commutative character of the space, through the exponential $`\mathrm{exp}(\stackrel{~}{p}^2/4t)`$, of the non-planar contribution coming from the high-momentum quanta flowing along the loop of the diagrams in fig. 3. Note that both the diamagnetic and paramagnetic functionals $`S^{\mathrm{diam}}`$ and $`S^{\mathrm{param}}`$ in eq. (2.9) contribute to the non-commutative IR logarithmic divergence of the vacuum polarization tensor, whereas the quadratic divergence only receives contributions from the diamagnetic functional, describing the orbital motion of $`D2`$ scalar quanta in the field $`B_\mu `$. We shall see in the next section that the non-logarithmic IR divergences at $`\stackrel{~}{p}=0`$ of the one-loop three-point function also have a purely diamagnetic origin. Finally, adding the contributions in eq. (2.27), we conclude that
$$i\mathrm{\Pi }_{\mu \nu }(p)\frac{i}{16\pi ^2}\left[\frac{22}{3}\right]\mathrm{ln}(p^2\stackrel{~}{p}^2)(p^2\eta _{\mu \nu }p_\mu p_\nu )+\frac{2i}{\pi ^2}\frac{\stackrel{~}{p}_\mu \stackrel{~}{p}_\nu }{\stackrel{~}{p}^4}i\mathrm{\Pi }_{\mu \nu }^{\mathrm{IR}}(p).$$
(2.30)
We have thus shown that the phenomenon of paramagnetic dominance at high energies, $`|a^{\mathrm{param}}|>a^{\mathrm{diam}}`$, which explains the sign of the one-loop beta function, also renders, via UV/IR mixing, the coefficient of the logarithmic divergence $`\mathrm{ln}(p^2\stackrel{~}{p}^2)`$ in the vacuum polarization tensor equal to $`22/3`$. This coefficient and the coefficient of $`\mathrm{ln}(p^2/4\pi \mu ^2)`$ in eq. (2.25) have the same absolute value but opposite sign, yet another manifestation of the mixing.
We shall close this section by showing that the right hand side of eqs. (2.15), (2.16), (2.20) and (2.21) also arise in the infinite tension limit of open string theory in a constant magnetic field. We will follow ref. and obtain the non-commutative $`U(1)`$ theory as the infinite tension limit of an open bosonic string theory on a D3-brane stuck at an $`\mathrm{I}\mathrm{R}^{22}/Z_2`$ orbifold singularity with a constant magnetic field along the worldvolume directions of the brane. The reader is referred to ref. for a comprehensive discussion. Let us first compute the one-loop non-planar contribution to the string amplitude with two photon vertex insertions of polarizations $`e_\mu ^{(1)}`$ and $`e_\nu ^{(2)}`$. This is given by
$$𝒜^{\mathrm{NP}}=e_\mu ^{(1)}i\mathrm{\Pi }^{\mathrm{NP}}{}_{}{}^{\mu \nu }e_{\nu }^{(2)},$$
where
$$\begin{array}{c}i\mathrm{\Pi }^{\mathrm{NP}}{}_{}{}^{\mu \nu }=g^2_0^{\mathrm{}}\frac{dt}{2t}_0^{2\pi t}dy_1_0^{2\pi t}dy_2Z(t)\hfill \\ \times \left(p_\rho p_\sigma _{y_1}\stackrel{~}{𝒢}^{\rho \sigma }_{y_2}\stackrel{~}{𝒢}^{\mu \nu }+p_\rho p_\sigma _{y_1}\stackrel{~}{𝒢}^{\mu \rho }_{y_2}\stackrel{~}{𝒢}^{\nu \sigma }\right)e^{p_\rho p_\sigma \stackrel{~}{𝒢}^{\rho \sigma }},\hfill \end{array}$$
(2.31)
$`Z(t)`$ is the open string partition function of an open bosonic string ending on a D$`p`$-brane glued at an orbifold singularity in the presence of a constant magnetic field along the directions of the brane, and $`\stackrel{~}{𝒢}^{\mu \nu }`$ is related to the string propagator and reads
$$\stackrel{~}{𝒢}^{\mu \nu }=\alpha ^{}\left[G^{\mu \nu }\mathrm{\Gamma }^{\mathrm{NP}}(y|t)+\frac{i\theta ^{\mu \nu }}{2\pi \alpha ^{}}\frac{y}{t}+\frac{(\theta G\theta )^{\mu \nu }}{(2\pi \alpha ^{})^2}\frac{\pi }{2t}\right].$$
(2.32)
Here $`G^{\mu \nu }`$ is the inverse of the effective open string metric $`G_{\mu \nu }`$ , $`y=y_1y_2`$ and $`\mathrm{\Gamma }^{\mathrm{NP}}`$ has the form
$$\mathrm{\Gamma }^{\mathrm{NP}}(y|t)=\mathrm{ln}\left|2\pi \frac{\vartheta _2(iy/2\pi ,it)}{\vartheta _1^{}(0,it)}\right|^2\frac{y^2}{2\pi t}.$$
(2.33)
From eqs. (2.32) and (2.31), one readily has
$$\begin{array}{c}i\mathrm{\Pi }^{\mathrm{NP}}{}_{}{}^{\mu \nu }=g^2_0^{\mathrm{}}\frac{dt}{2t}_0^{2\pi t}dy_1_0^{2\pi t}dy_2Z(t)\alpha _{}^{}{}_{}{}^{2}\{(p^2G^{\mu \nu }p^\mu p^\nu )\left[_y\mathrm{\Gamma }^{\mathrm{NP}}(y|t)\right]^2\hfill \\ +\frac{i(p^2\theta ^{\mu \nu }p^\mu \stackrel{~}{p}^\nu p^\nu \stackrel{~}{p}^\mu )}{2\pi \alpha ^{}}\frac{1}{t}_y\mathrm{\Gamma }^{\mathrm{NP}}(y|t)+\frac{\stackrel{~}{p}^\mu \stackrel{~}{p}^\nu }{(2\pi \alpha ^{})^2}\frac{1}{t^2}\}e^{\alpha ^{}p^2\mathrm{\Gamma }^{\mathrm{NP}}(y|t)}e^{\stackrel{~}{p}^2/8\pi \alpha ^{}t},\hfill \end{array}$$
(2.34)
with $`p^2=G_{\mu \nu }p^\mu p^\nu `$, $`p_\nu =G_{\nu \mu }p^\mu `$ and $`\stackrel{~}{p}^\mu =\theta ^{\mu \nu }p_\nu `$. To obtain from this expression the corresponding field theory result, one needs the large $`t`$ expansion of the functions in the integrand upon performing the change of variables $`y=2\pi xt`$. The relevant terms of these expansions are
$`Z(t){\displaystyle \frac{V_d}{(8\pi ^2\alpha ^{}t)^{d/2}}}\left(e^{2\pi t}+d2\right)`$ (2.35)
$`\mathrm{\Gamma }^{\mathrm{NP}}(2\pi xt|t)2\pi tx(1x)`$ (2.36)
$`_y\mathrm{\Gamma }^{\mathrm{NP}}(y|t)|_{y=2\pi xt}12x+2\left(e^{2\pi t}e^{2\pi xt}e^{2\pi xt}\right),`$ (2.37)
where we have taken a $`\mathrm{D}(d1)`$-brane to compute $`Z(t)`$. Substituting these expansions in eq. (2.34) and making the changes of scale $`tt/\alpha ^{}`$ and $`yy/\alpha ^{}`$, one obtains after some calculations<sup>5</sup><sup>5</sup>5This includes dropping divergent contributions as $`t\mathrm{}`$ for arbitrary $`x`$. the non-planar contribution to the vacuum polarization tensor as the sum
$$i\mathrm{\Pi }_{\mu \nu }^{\mathrm{NP}}(p)=i\mathrm{\Pi }_{\mu \nu }^{\mathrm{NP},\mathrm{tach}}(p)+i\mathrm{\Pi }_{\mu \nu }^{\mathrm{NP},\mathrm{\hspace{0.17em}0}\mathrm{mode}}(p)$$
of two terms:
$$\begin{array}{cc}\hfill i\mathrm{\Pi }_{\mu \nu }^{\mathrm{NP},\mathrm{\hspace{0.17em}0}\mathrm{mode}}(p)=ig_0^2\frac{[d2]}{(4\pi )^{d/2}}& [(p^2G_{\mu \nu }p_\mu p_\nu )_0^1dx(12x)^2_0^{\mathrm{}}dtt^{1d/2}e^{p^2tx(1x)\stackrel{~}{p}^2/4t}\hfill \\ & +\stackrel{~}{p}_\mu \stackrel{~}{p}_\nu _0^1dx_0^{\mathrm{}}dtt^{1d/2}e^{p^2tx(1x)\stackrel{~}{p}^2/4t}]\hfill \end{array}$$
(2.38)
and
$$i\mathrm{\Pi }_{\mu \nu }^{\mathrm{NP},\mathrm{tach}}(p)=ig_0^2\frac{[\mathrm{\hspace{0.17em}8}]}{(4\pi )^{d/2}}\left(p^2G_{\mu \nu }p_\mu p_\nu \right)_0^1𝑑x_0^{\mathrm{}}𝑑tt^{1d/2}e^{p^2tx(1x)\stackrel{~}{p}^2/4t}.$$
(2.39)
The same arguments as those used in ref. for Yang-Mills theory on commutative Minkowski space-time show that $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{NP},\mathrm{\hspace{0.17em}0}\mathrm{mode}}(p)`$ originates from the product of the one-loop photon contribution to the partition function and the zero-mode contribution to the string Green function. In turn, $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{NP},\mathrm{tach}}(p)`$ arises from the combination of the tachyonic contribution to the partition function and the exponentially vanishing part of the string Green function, an effect called tachyon magnification in ref. . This can be understood as follows.
Indeed, the contribution proportional to $`p^2G_{\mu \nu }p_\mu p_\nu `$ in $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{NP},\mathrm{\hspace{0.17em}0}\mathrm{mode}}(p)`$ arises from combining the term $`d2`$ in eq. (2.35), carrying the one-loop photon contribution to $`Z(t)`$, with the square of the zero-mode term $`12x`$ in eq. (2.37). And the term proportional to $`\stackrel{~}{p}_\mu \stackrel{~}{p}_\nu `$ in $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{NP},\mathrm{\hspace{0.17em}0}\mathrm{mode}}(p)`$, note that it comes with the extra zero-mode factor $`1/t^2`$. As concerns $`i\mathrm{\Pi }_{\mu \nu }^{\mathrm{NP},\mathrm{tach}}(p)`$ , we note that it originates from the term in eq. (2.34) proportional to $`p^2G_{\mu \nu }p_\mu p_\nu `$ that goes with the product
$$e^{2\pi t}\times 4\left(e^{2\pi t}e^{2\pi xt}e^{2\pi xt}\right)^2.$$
The first factor in this product, $`e^{2\pi t}`$, is supplied by the large-$`t`$ expansion of $`Z(t)`$ and is due to the open string tachyon. The second factor arises from squaring the term $`2(e^{2\pi t}e^{2\pi xt}e^{2\pi xt})`$ in eq. (2.37) and has its origin in the open string propagator. It is precisely $`8e^{2\pi t}`$, the evanescent large-$`t`$ and $`x`$-independent contribution to this second factor, which is magnified when combined with the $`t`$-divergent tachyon contribution to $`Z(t)`$, all in all yielding a finite non-zero result. Note that eqs. (2.38) and (2.39) agree with eqs. (2.16) and (2.21), if we identify $`G_{\mu \nu }`$ with $`\eta _{\mu \nu }`$ and $`d`$ with $`D`$. Hence we have identified the diamagnetic and paramagnetic non-planar contributions in field theory with the zero-mode and tachyon magnification contributions in string theory, an observation first made for Yang-Mills theory in commutative Minkowski space time in ref. .
A similar analysis can be performed for the planar contribution to the string amplitude with two photon vertex insertions. It is this contribution that develops an UV divergence when the field theory limit is taken. Furthermore, in this limit, this contribution can be decomposed in the sum of a zero-mode part and a tachyon magnification part, whose expressions agree with the r.h.s. of eqs. (2.15) and (2.20) respectively. Thus, the charge screening contribution to the beta function –given by the coefficient $`a_1^{diam}`$– is concocted by the diamagnetic contributions in the field theory setting and has a zero mode origin in string theory framework. On the other hand, the charge anti-screening contribution –provided by $`a_1^{\mathrm{param}}`$– has a paramagnetic origin in field theory and is the result of tachyon magnification in string theory.
## 3 The three-point function
In this section we study, in the light of high-energy paramagnetic dominance, the UV and non-commutative IR divergences of the three-point function and their mixing. The 1PI three-point function $`\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}(x_1,x_2,x_3)`$ is defined as
$$i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}(x_1,x_2,x_3)=\frac{\delta ^3\mathrm{\Gamma }[B]}{\delta B_{\mu _1}(x_1)\delta B_{\mu _2}(x_2)\delta B_{\mu _3}(x_3)}|_{B=0},$$
where $`i\mathrm{\Gamma }[B]=\mathrm{ln}Z[B]`$ and $`Z[B]`$ is given in eq. (2.8). In momentum space we will write $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}(p_1,p_2,p_3)`$, with $`p_1+p_2+p_3=0`$. The three-point function can be expressed as the sum of a diamagnetic contribution, which sums over diagrams with all vertices of diamagnetic type, a paramagnetic contribution, given by diagrams with only paramagnetic vertices, and a mixed contribution, made of diagrams with at least one vertex of each type:
$$i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}(p_1,p_2,p_3)=i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{diam}}(p_1,p_2,p_3)+i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{param}}(p_1,p_2,p_3)+i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{mixed}}(p_1,p_2,p_3).$$
We start computing the diamagnetic contribution $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{diam}}(p_1,p_2,p_3)`$ to the three-point function. The one-loop diamagnetic Feynman diagrams that contribute to this function are shown in fig. 5. The Feynman diagrams in the first row will be called triangle diagrams, while the remaining diagrams will be referred to as swordfish diagrams. The sum $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{triang}}(p_1,p_2,p_3)`$ of the triangle diagrams is formally given in $`D`$ dimensions by the integral
$$\begin{array}{c}(\mathrm{a})+(\mathrm{b})+(\mathrm{c})=i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{triang}}(p_1,p_2,p_3)\hfill \\ =\mathrm{\hspace{0.17em}8}i[D2]\frac{d^Dq}{(2\pi )^D}\frac{(2q+p_1)_{\mu _1}(2q+2p_1+p_2)_{\mu _2}(2q+p_1+p_2)_{\mu _3}}{q^2(q+p_1)^2(q+p_1+p_2)^2}\hfill \\ \times \mathrm{sin}\left[\frac{qp_1}{2}\right]\mathrm{sin}\left[\frac{(q+p_1)p_2}{2}\right]\mathrm{sin}\left[\frac{q(p_1+p_2)}{2}\right].\hfill \end{array}$$
(3.40)
The sum $`i\mathrm{I}_{\mu _1\mu _2\mu _3}^{\mathrm{sfish}}(p_1,p_2,p_3)`$ of the swordfish diagrams (d) and (g), also in $`D`$ dimensions, is given by
$$\begin{array}{c}(\mathrm{d})+(\mathrm{g})=i\mathrm{I}_{\mu _1\mu _2\mu _3}^{\mathrm{sfish}}(p_1,p_2,p_3)\hfill \\ =\mathrm{\hspace{0.17em}4}i[D2]\eta _{\mu _1\mu _2}\frac{d^Dq}{(2\pi )^D}\frac{(2q+p_1+p_2)_{\mu _3}}{q^2(q+p_1)^2(q+p_1+p_2)^2}\mathrm{sin}\left[\frac{q(p_1+p_2)}{2}\right]\hfill \\ \times \left\{\mathrm{sin}\left[\frac{(q+p_1)p_2}{2}\right]\mathrm{sin}\left[\frac{qp_1}{2}\right]+\mathrm{sin}\left[\frac{(q+p_2)p_1}{2}\right]\mathrm{sin}\left[\frac{qp_2}{2}\right]\right\},\hfill \end{array}$$
(3.41)
and the total swordfish contribution $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{sfish})}(p_1,p_2,p_3)`$ reads
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{sfish}}(p_1,p_2,p_3)=(\mathrm{d})+(\mathrm{g})+(\mathrm{e})+(\mathrm{h})+(\mathrm{f})+(\mathrm{i})\hfill \\ =i\mathrm{I}_{\mu _1\mu _2\mu _3}^{\mathrm{sfish}}(p_1,p_2,p_3)+i\mathrm{I}_{\mu _3\mu _2\mu _1}^{\mathrm{sfish}}(p_3,p_2,p_1)+i\mathrm{I}_{\mu _1\mu _3\mu _2}^{\mathrm{sfish}}(p_1,p_3,p_2).\hfill \end{array}$$
As in the vacuum polarization tensor case, the overall factors $`[D2]`$ in eqs. (3.40) and (3.41) account for the fact that the one-loop diamagnetic contribution is due to the orbital motion of $`D2`$ real scalar quanta in the background field $`B_\mu `$. The $`D`$ comes from the diagrams with fields $`Q_\mu `$ $`(\mu =1,\mathrm{},D)`$ going around the loop, while the $`2`$ is provided by the corresponding diagrams with ghost loops. Each diagram in fig. 5 can be expressed as the sum of its planar part and its non-planar part by repeated use of the identity $`2i\mathrm{sin}\left(\frac{1}{2}qp\right)=e^{iqp/2}e^{iqp/2}`$. The planar part is the contribution whose $`\theta ^{\mu \nu }`$-dependent complex exponential factors do not depend on the loop momenta, while the non-planar part collects all terms whose $`\theta ^{\mu \nu }`$-dependent complex exponentials depend on the loop momenta. At $`D=4`$, the planar parts of the diagrams in fig. 5 are UV divergent by power counting. This UV divergence is logarithmic and shows in dimensional regularization as a simple pole at $`D=4`$. By contrast, the non-planar parts are finite since the UV divergence is tamed by the the exponentials $`e^{iqp_i}`$, provided $`p_i0`$. In other words, the contribution to the integrals in eqs. (3.40) and (3.41) of the quantum modes flowing along the loop with arbitrarily high momenta is UV divergent when one considers their planar part, and it is turned into an IR divergence by the non-commutativity of space when one looks at their non-planar parts. This IR divergence occurs whenever any of the conditions
$$\stackrel{~}{p}_1=0\stackrel{~}{p}_2=0\stackrel{~}{p}_3=0$$
(3.42)
is met. This provides a qualitative explanation of the UV/IR for the diamagnetic part of the three-point function. To make it quantitative, we need to compute the UV divergent terms in the planar contribution and the non-commutative IR divergent parts of the non-planar contribution.
To calculate the UV divergence of the planar contribution, we make a Laurent expansion about $`D=4`$ of the corresponding dimensionally regularized integrals and retain singular contributions. Proceeding in this way and defining $`D=4+2\epsilon `$, we obtain for non-exceptional external momenta
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{trian},\mathrm{P}}(p_1,p_2,p_3)=\frac{i}{16\pi ^2\epsilon }\frac{2}{3}\mathrm{\hspace{0.17em}2}i\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\hfill \\ \times \left[\eta _{\mu _1\mu _2}(p_2p_1)_{\mu _3}+\eta _{\mu _1\mu _3}(p_1p_3)_{\mu _2}+\eta _{\mu _2\mu _3}(p_3p_2)_{\mu _1}\right]+O(\epsilon ^0),\hfill \end{array}$$
for the sum of the triangle diagrams, and
$$i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{sfish},\mathrm{P}}(p_1,p_2,p_3)=O(\epsilon ^0)$$
for the sum of the swordfish diagrams. From this we conclude that the UV divergent contribution $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{diam},\mathrm{UV}}(p_1,p_2,p_3)`$ of the diamagnetic part to the three-point function is given in dimensional regularization by
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{diam},\mathrm{UV}}(p_1,p_2,p_3)=\frac{ia_1^{\mathrm{diam}}}{\epsilon }\mathrm{\hspace{0.17em}2}i\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\hfill \\ \times \left[\eta _{\mu _1\mu _2}(p_2p_1)_{\mu _3}+\eta _{\mu _1\mu _3}(p_1p_3)_{\mu _2}+\eta _{\mu _2\mu _3}(p_3p_2)_{\mu _1}\right],\hfill \end{array}$$
(3.43)
where $`a_1^{\mathrm{diam}}`$ is as in eq. (2.23). To compute the IR divergences that appear in the non-planar contribution when one or more of the $`\stackrel{~}{p}_i`$ vanish, we set $`D=4`$ and use the following basic results for $`\stackrel{~}{k}0`$:
$`{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{e^{iqk}}{q^2(q+\mathrm{}_1)^2}}{\displaystyle \frac{i}{16\pi ^2}}\mathrm{ln}\stackrel{~}{k}^2`$ (3.44)
$`{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{q_{\mu _1}q_{\mu _2}e^{iqk}}{q^2(q+\mathrm{}_1)^2(q+\mathrm{}_2)^2}}{\displaystyle \frac{i}{16\pi ^2}}{\displaystyle \frac{1}{4}}\mathrm{ln}\stackrel{~}{k}^2\eta _{\mu _1\mu _2}`$ (3.45)
$$\begin{array}{c}\frac{d^4q}{(2\pi )^4}\frac{q_{\mu _1}q_{\mu _2}q_{\mu _3}e^{iqk}}{q^2(q+\mathrm{}_1)^2(q+\mathrm{}_2)^2}\frac{i}{16\pi ^2}\hfill \\ \times \{\frac{i}{2\stackrel{~}{k}^2}[\eta _{\mu _1\mu _2}\stackrel{~}{k}_{\mu _3}+\eta _{\mu _2\mu _3}\stackrel{~}{k}_{\mu _1}+\eta _{\mu _1\mu _3}\stackrel{~}{k}_{\mu _2}2\frac{\stackrel{~}{k}_{\mu _1}\stackrel{~}{k}_{\mu _2}\stackrel{~}{k}_{\mu _3}}{\stackrel{~}{k}^2}]\hfill \\ +\frac{1}{12}\mathrm{ln}\stackrel{~}{k}^2[\eta _{\mu _1\mu _2}(\mathrm{}_1+\mathrm{}_2)_{\mu _3}+\eta _{\mu _2\mu _3}(\mathrm{}_1+\mathrm{}_2)_{\mu _1}+\eta _{\mu _1\mu _3}(\mathrm{}_1+\mathrm{}_2)_{\mu _2}]\},\hfill \end{array}$$
(3.46)
where $`\mathrm{}_1`$ and $`\mathrm{}_2`$ never vanish and $``$ indicates that we have dropped any contribution which remains finite as $`\stackrel{~}{k}0`$. These expressions, together with some algebra, lead to the following divergences for the non-planar parts of the diagrams in fig. 5 when one or more $`\stackrel{~}{p}_i`$ approach zero:
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{triang},\mathrm{NP}}(p_1,p_2,p_3)\frac{i}{16\pi ^2}\frac{2}{9}\mathrm{\hspace{0.17em}2}i\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\left(\mathrm{ln}\stackrel{~}{p}_1^2+\mathrm{ln}\stackrel{~}{p}_2^2+\mathrm{ln}\stackrel{~}{p}_3^2\right)\hfill \\ \times \left[\eta _{\mu _1\mu _2}(p_2p_1)_{\mu _3}+\eta _{\mu _1\mu _3}(p_1p_3)_{\mu _2}+\eta _{\mu _2\mu _3}(p_3p_2)_{\mu _1}\right]\hfill \\ +\frac{1}{\pi ^2}\mathrm{cos}\left(\frac{p_1p_2}{2}\right)\underset{i=1}{\overset{3}{}}\left[\frac{2(\stackrel{~}{p}_i)_{\mu _1}(\stackrel{~}{p}_i)_{\mu _2}(\stackrel{~}{p}_i)_{\mu _3}}{(\stackrel{~}{p}_i^2)^2}\eta _{\mu _1\mu _2}\frac{(\stackrel{~}{p}_i)_{\mu _3}}{\stackrel{~}{p}_i^2}\eta _{\mu _2\mu _3}\frac{(\stackrel{~}{p}_i)_{\mu _1}}{\stackrel{~}{p}_i^2}\eta _{\mu _1\mu _3}\frac{(\stackrel{~}{p}_i)_{\mu _2}}{\stackrel{~}{p}_i^2}\right]\hfill \end{array}$$
and
$$i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{sfish},\mathrm{NP}}(p_1,p_2,p_3)\frac{1}{\pi ^2}\mathrm{cos}\left(\frac{p_1p_2}{2}\right)\underset{i=1}{\overset{3}{}}\left[\eta _{\mu _1\mu _2}\frac{(\stackrel{~}{p}_i)_{\mu _3}}{\stackrel{~}{p}_i^2}+\eta _{\mu _2\mu _3}\frac{(\stackrel{~}{p}_i)_{\mu _1}}{\stackrel{~}{p}_i^2}+\eta _{\mu _1\mu _3}\frac{(\stackrel{~}{p}_i)_{\mu _2}}{\stackrel{~}{p}_i^2}\right].$$
Adding these two expressions we obtain for the non-commutative IR divergent part of the diamagnetic contribution to the three-point function
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{diam},\mathrm{IR}}(p_1,p_2,p_3)=\frac{2}{\pi ^2}\mathrm{cos}\left(\frac{p_1p_2}{2}\right)\underset{i=1}{\overset{3}{}}\frac{(\stackrel{~}{p}_i)_{\mu _1}(\stackrel{~}{p}_i)_{\mu _2}(\stackrel{~}{p}_i)_{\mu _3}}{(\stackrel{~}{p}_i^2)^2}\hfill \\ +\frac{i}{3}a^{\mathrm{diam}}\mathrm{\hspace{0.17em}2}i\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\left(\mathrm{ln}\stackrel{~}{p}_1^2+\mathrm{ln}\stackrel{~}{p}_2^2+\mathrm{ln}\stackrel{~}{p}_3^2\right)\hfill \\ \times \left[\eta _{\mu _1\mu _2}(p_2p_1)_{\mu _3}+\eta _{\mu _1\mu _3}(p_1p_3)_{\mu _2}+\eta _{\mu _2\mu _3}(p_3p_2)_{\mu _1}\right],\hfill \end{array}$$
(3.47)
where $`a_1^{\mathrm{diam}}`$ is given in eq. (2.23). Notice that the logarithmic term in this expression is dual to the the UV divergent contribution given in eq. (3.43), in the sense that in the momenta region
$$|\stackrel{~}{p}_1||\stackrel{~}{p}_2||\stackrel{~}{p}_3|\theta \mathrm{\Lambda }_{\mathrm{IR}}0$$
(3.48)
the former can be obtained from the latter by using the definitions (2.29). Note also that the factor $`\mathrm{sin}(\frac{1}{2}(p_1p_2)`$, together with momentum conservation $`\stackrel{~}{p}_1+\stackrel{~}{p}_2+\stackrel{~}{p}_3=0`$, renders the second term in eq. (3.47) finite when $`\stackrel{~}{p}_i0`$ $`(i=1,2,3)`$. This is so in spite of the logarithmic non-commutative IR divergences of the integrals (3.44)-(3.46) that enter the computation of the non-planar part of the diamagnetic contribution, and can be thought of as a consequence of gauge invariance –which determines the momentum structure of the UV divergent part of the three-point function in eq. (3.43)– and the UV/IR mixing mechanism. Indeed, UV/IR mixing in the form of eq. (2.29) allows to retrieve from the UV divergent contribution $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{diam},\mathrm{UV}}(p_1,p_2,p_3)`$ the terms in $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{diam},\mathrm{NP}}(p_1,p_2,p_3)`$ which involve $`\mathrm{ln}\stackrel{~}{p}_i^2`$, with $`\stackrel{~}{p}_i`$ verifying eq. (3.48).
Let us compute now the paramagnetic part $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{param}}(p_1,p_2,p_3)`$ of the the three-point function. This is given by the contribution involving only three background fields that comes from the diagram in fig. 4, which in turn is given by eqs. (2.17) and (2.18). Extracting from eq. (2.17) the contribution cubic in the background fields and using the results of section 2 for the integral $`J(p)`$, we obtain, after some calculations,
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{param},\mathrm{UV}}(p_1,p_2,p_3)=\frac{ia^{\mathrm{param}}}{\epsilon }\mathrm{\hspace{0.17em}2}i\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\hfill \\ \times \left[\eta _{\mu _1\mu _2}(p_2p_1)_{\mu _3}+\eta _{\mu _1\mu _3}(p_1p_3)_{\mu _2}+\eta _{\mu _2\mu _3}(p_3p_2)_{\mu _1}\right]\hfill \end{array}$$
(3.49)
for the UV divergent part, and
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{param},\mathrm{IR}}(p_1,p_2,p_3)=\frac{i}{3}a^{\mathrm{param}}\mathrm{\hspace{0.33em}2}i\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\left(\mathrm{ln}\stackrel{~}{p}_1^2+\mathrm{ln}\stackrel{~}{p}_2^2+\mathrm{ln}\stackrel{~}{p}_3^2\right)\hfill \\ \times \left[\eta _{\mu _1\mu _2}(p_2p_1)_{\mu _3}+\eta _{\mu _1\mu _3}(p_1p_3)_{\mu _2}+\eta _{\mu _2\mu _3}(p_3p_2)_{\mu _1}\right]\hfill \end{array}$$
(3.50)
for the IR divergences that appear as $`\stackrel{~}{p}_i0`$. Here $`a_1^{\mathrm{param}}`$ is defined in eq. (2.23). As in the case of the diamagnetic contribution, the IR part $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{param},\mathrm{IR}}(p_1,p_2,p_3)`$ is dual to the UV part $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{param},\mathrm{UV}}(p_1,p_2,p_3)`$, in the sense that one can go from one to another by using the identifications (2.29) in the momentum region (3.48).
We finally look at the mixed part $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{mixed}}(p_1,p_2,p_3)`$. It receives contributions from the diagrams in figs. 2 and 6. As explained in section 2, the diagram of fig. 2 vanishes. The very same arguments as for the diagram in fig. 2 show that the diagram in fig. 6(a) also vanishes. By contrast, the diagrams in figs, 6(b) and 6(c) are different from zero; yet their planar parts are finite by power counting and their non-planar parts are IR finite for vanishing $`\stackrel{~}{p}_i`$. Altogether, we conclude
$$\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{mixed},\mathrm{UV}}(p_1,p_2,p_3)=\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{mixed},\mathrm{IR}}(p_1,p_2,p_3)=0.$$
(3.51)
The UV divergent part $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{UV}}(p_1,p_2,p_3)`$ of the three-point function will be the sum of the diamagnetic, paramagnetic and mixed UV divergent parts. Using the results obtained, we have
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{UV}}(p_1,p_2,p_3)=\frac{i}{\epsilon }\left(a^{\mathrm{diam}}+a^{\mathrm{param}}\right)\mathrm{\hspace{0.17em}2}i\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\hfill \\ \times \left[\eta _{\mu _1\mu _2}(p_2p_1)_{\mu _3}+\eta _{\mu _1\mu _3}(p_1p_3)_{\mu _2}+\eta _{\mu _2\mu _3}(p_3p_2)_{\mu _1}\right],\hfill \end{array}$$
(3.52)
with $`a^{\mathrm{diam}}`$ and $`a^{\mathrm{param}}`$ as in eq. (2.23). It is clear that the counterterm in eq. (2.24) subtracts the UV divergences in the three-point function and that the beta function computed from the three-point function is also given by eq. (2.26), all this in accordance with gauge invariance. We again see that the negative sign of the one-loop beta function is due to the dominance of the paramagnetic high-momentum contributions over the diamagnetic ones. We close this section displaying the leading non-commutative IR behaviour of the three-point function, obtained by summing the contributions in eqs. (3.47) and (3.50):
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{IR}}(p_1,p_2,p_3)=\frac{2}{\pi ^2}\mathrm{cos}\left(\frac{p_1p_2}{2}\right)\underset{i=1}{\overset{3}{}}\frac{(\stackrel{~}{p}_i)_{\mu _1}(\stackrel{~}{p}_i)_{\mu _2}(\stackrel{~}{p}_i)_{\mu _3}}{(\stackrel{~}{p}_i^2)^2}\hfill \\ +\frac{i}{3}\left(a^{\mathrm{diam}}+a^{\mathrm{param}}\right)\mathrm{\hspace{0.17em}2}i\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\left(\mathrm{ln}\stackrel{~}{p}_1^2+\mathrm{ln}\stackrel{~}{p}_2^2+\mathrm{ln}\stackrel{~}{p}_3^2\right)\hfill \\ \times \left[\eta _{\mu _1\mu _2}(p_2p_1)_{\mu _3}+\eta _{\mu _1\mu _3}(p_1p_3)_{\mu _2}+\eta _{\mu _2\mu _3}(p_3p_2)_{\mu _1}\right].\hfill \end{array}$$
(3.53)
Notice that the non-logarithmic non-commutative IR contributions to the three-point are only supplied by its diamagnetic part \[see eqs. (3.47) and (3.50)\]; an effect that was met for the first time when we studied the vacuum polarization tensor.
## 4 The four-point function
In this section we compute the UV divergent and the leading non-commutative IR part of the four-point function and exhibit how UV/IR connects them. The four-point function $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}(x_1,x_2,x_3,x_4)`$ is defined by
$$i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}(x_1,x_2,x_3,x_4)=\frac{\delta ^4\mathrm{\Gamma }[B]}{\delta B_{\mu _1}(x_1)\delta B_{\mu _2}(x_2)\delta B_{\mu _3}(x_3)\delta B_{\mu _4}(x_4)}|_{B=0},$$
(4.54)
where $`i\mathrm{\Gamma }[B]=\mathrm{ln}Z[B]`$ and $`Z[B]`$ is as in eq. (2.8). In momentum space we will write $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}(p_1,p_2,p_3,p_4)`$, with $`p_1+p_2+p_3+p_4=0`$, and will be the sum of a diamagnetic, a paramagnetic and a mixed contribution:
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}(p_1,p_2,p_3,p_4)=i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{diam}}(p_1,p_2,p_3,p_4)\hfill \\ +i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{param}}(p_1,p_2,p_3,p_4)+i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{mixed}}(p_1,p_2,p_3,p_4).\hfill \end{array}$$
Let us consider each one of these contributions separately.
The diamagnetic contribution is given by the sum of the diagrams in fig. 7. The diagrams in the first, second and third row of the figure have different topologies and will be called box, lynx and bubble diagrams respectively. Each one of these diagrams is a sum of a planar part and a non-planar part. The same arguments as in sections 2 and 3 show that the planar parts are UV divergent by power counting at $`D=4`$, the divergence being logarithmic and being characterized as a pole in dimensional regularization. As regards the non-planar parts, these are finite for $`\stackrel{~}{p}_i0`$ and develop divergences for any of the configurations
$$\stackrel{~}{p}_1=0\stackrel{~}{p}_2=0\stackrel{~}{p}_3=0\stackrel{~}{p}_4=0\stackrel{~}{p}_1+\stackrel{~}{p}_2=0\stackrel{~}{p}_1+\stackrel{~}{p}_3=0\stackrel{~}{p}_2+\stackrel{~}{p}_3=0.$$
(4.55)
To calculate the UV divergences, we make a Laurent expansion about $`D=4`$ of the dimensionally regularized integrals defining the planar parts and obtain, for non-exceptional momenta:
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{box},\mathrm{P}}(p_1,p_2,p_3,p_4)=\frac{i}{16\pi ^2\epsilon }\frac{8}{3}\left(\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}+\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}+\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3}\right)\hfill \\ \times [\mathrm{cos}\left(\frac{p_1p_2+p_1p_3+p_2p_3}{2}\right)+\mathrm{cos}\left(\frac{p_1p_2+p_1p_3p_2p_3}{2}\right)\hfill \\ +\mathrm{cos}\left(\frac{p_1p_2p_1p_3p_2p_3}{2}\right)]+O(\epsilon ^0)\hfill \end{array}$$
(4.56)
for the sums of the of the box diagrams, and
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{lynx},\mathrm{P}}(p_1,p_2,p_3,p_4)=\frac{1}{2}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{bub},\mathrm{P}}(p_1,p_2,p_3,p_4)+O(\epsilon ^0)\hfill \\ =\frac{i}{16\pi ^2\epsilon }\mathrm{\hspace{0.33em}8}[(\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}+\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3})\mathrm{cos}\left(\frac{p_1p_2+p_1p_3+p_2p_3}{2}\right)\hfill \\ +\left(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}+\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3}\right)\mathrm{cos}\left(\frac{p_1p_2+p_1p_3p_2p_3}{2}\right)\hfill \\ +(\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}+\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4})\mathrm{cos}\left(\frac{p_1p_2p_1p_3p_2p_3}{2}\right)]+O(\epsilon ^0).\hfill \end{array}$$
(4.57)
for the sum of the lynx and bubble diagrams. Here, as in previous sections, $`D=4+2\epsilon `$. Summing these three contributions and doing some simple trigonometry, we arrive at the following UV divergences for the diamagnetic part of the four-point function:
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{diam},\mathrm{UV}}(p_1,p_2,p_3,p_4)=\frac{ia^{\mathrm{diam}}}{\epsilon }[(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3})\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\mathrm{sin}\left(\frac{p_3p_4}{2}\right)\hfill \\ +\left(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}\right)\mathrm{sin}\left(\frac{p_1p_4}{2}\right)\mathrm{sin}\left(\frac{p_3p_2}{2}\right)\hfill \\ +(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4})\mathrm{sin}\left(\frac{p_4p_2}{2}\right)\mathrm{sin}\left(\frac{p_3p_1}{2}\right)],\hfill \end{array}$$
(4.58)
with $`a^{\mathrm{diam}}`$ given in eq. (2.23). Let us stress that, unlike for the three-point function, the trigonometric and tensor structure of the box, lynx and bubble contributions is different from that in the classical action $`S_{\mathrm{class}}[B]`$ for the quartic vertex $`B^4`$, and that it is precisely the sum over diagrams with different topologies what produces the correct structure, thus making non-trivial the answer given by eq. (4.58). Of course, this is so because gauge invariance is at work and brings the necessary simplifications.
We next compute the leading non-commutative IR behaviour of the non-planar parts of the diagrams in fig. 7. To do this, we use eqs. (3.44) and (3.45) and the result
$$\begin{array}{c}\frac{d^4q}{(2\pi )^4}\frac{q_{\mu _1}q_{\mu _2}q_{\mu _3}q_{\mu _4}e^{iqk}}{q^2(q+\mathrm{}_1)^2(q+\mathrm{}_2)^2(q+\mathrm{}_3)^2}\hfill \\ \frac{i}{16\pi ^2}\frac{1}{24}\mathrm{ln}\stackrel{~}{k}^2\left(\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}+\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}+\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3}\right)\hfill \end{array}$$
for $`\stackrel{~}{k}0`$, where $`\mathrm{}_1,\mathrm{}_2`$ and $`\mathrm{}_3`$ never vanish and $``$ indicates that we have dropped contributions which remain finite as $`\stackrel{~}{k}0`$. After some algebra, we obtain for the three sets of diagrams under consideration:
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{box},\mathrm{NP}}(p_1,p_2,p_3,p_4)\frac{i}{16\pi ^2}\left(\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}+\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}+\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3}\right)\hfill \\ \times \mathrm{\hspace{0.17em}4}\{\mathrm{cos}\left(\frac{p_1p_2+p_1p_3+p_2p_3}{2}\right)[\mathrm{ln}(\stackrel{~}{p}_1+\stackrel{~}{p}_2)^2+\mathrm{ln}(\stackrel{~}{p}_2+\stackrel{~}{p}_3)^2\frac{2}{3}\underset{i=1}{\overset{4}{}}\mathrm{ln}\stackrel{~}{p}_i^2]\hfill \\ +\mathrm{cos}\left(\frac{p_1p_2+p_1p_3p_2p_3}{2}\right)[\mathrm{ln}(\stackrel{~}{p}_2+\stackrel{~}{p}_3)^2+\mathrm{ln}(\stackrel{~}{p}_1+\stackrel{~}{p}_3)^2\frac{2}{3}\underset{i=1}{\overset{4}{}}\mathrm{ln}\stackrel{~}{p}_i^2]\hfill \\ +\mathrm{cos}\left(\frac{p_1p_2+p_1p_3p_2p_3}{2}\right)[\mathrm{ln}(\stackrel{~}{p}_1+\stackrel{~}{p}_2)^2+\mathrm{ln}(\stackrel{~}{p}_1+\stackrel{~}{p}_3)^2\frac{2}{3}\underset{i=1}{\overset{4}{}}\mathrm{ln}\stackrel{~}{p}_i^2]\}\hfill \end{array}$$
(4.59)
and
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{lynx},\mathrm{NP}}(p_1,p_2,p_3,p_4)\frac{1}{2}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{bub},\mathrm{NP}}(p_1,p_2,p_3,p_4)\hfill \\ \frac{i}{16\pi ^2}\mathrm{\hspace{0.33em}8}\{[(\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}+\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3})\mathrm{cos}\left(\frac{p_1p_2+p_1p_3+p_2p_3}{2}\right)\hfill \\ +\left(\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}+\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\right)\mathrm{cos}\left(\frac{p_1p_2p_1p_3p_2p_3}{2}\right)\hfill \\ +(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}+\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3})\mathrm{cos}(\frac{p_1p_2+p_1p_3p_2p_3}{2})]_{i=1}^4\mathrm{ln}\stackrel{~}{p}_i^2\hfill \\ [(\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}+2\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3})\mathrm{cos}\left(\frac{p_1p_2+p_1p_3+p_2p_3}{2}\right)\hfill \\ +(\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}+2\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4})\mathrm{cos}(\frac{p_1p_2p_1p_3p_2p_3}{2})]\mathrm{ln}(\stackrel{~}{p}_1+\stackrel{~}{p}_2)^2\hfill \\ [(2\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}+\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3})\mathrm{cos}\left(\frac{p_1p_2+p_1p_3+p_2p_3}{2}\right)\hfill \\ +(2\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}+\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3})\mathrm{cos}(\frac{p_1p_2+p_1p_3p_2p_3}{2})]\mathrm{ln}(\stackrel{~}{p}_2+\stackrel{~}{p}_3)^2\hfill \\ [(2\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}+\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4})\mathrm{cos}\left(\frac{p_1p_2p_1p_3p_2p_3}{2}\right)\hfill \\ +(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}+2\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3})\mathrm{cos}(\frac{p_1p_2p_1p_3p_2p_3}{2})]\mathrm{ln}(\stackrel{~}{p}_1+\stackrel{~}{p}_3)^2\}.\hfill \end{array}$$
(4.60)
Summing these three results we obtain for the leading non-commutative IR diamagnetic contribution to the four-point function
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{diam},\mathrm{IR}}(p_1,p_2,p_3,p_4)=4ia^{\mathrm{diam}}\hfill \\ \times \{(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3})\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\mathrm{sin}\left(\frac{p_3p_4}{2}\right)[\underset{i=1}{\overset{4}{}}\mathrm{ln}\stackrel{~}{p}_i^23\mathrm{ln}(\stackrel{~}{p}_1+\stackrel{~}{p}_2)^2]\hfill \\ +(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4})\mathrm{sin}\left(\frac{p_1p_4}{2}\right)\mathrm{sin}\left(\frac{p_3p_2}{2}\right)[\underset{i=1}{\overset{4}{}}\mathrm{ln}\stackrel{~}{p}_i^23\mathrm{ln}(\stackrel{~}{p}_2+\stackrel{~}{p}_3)^2]\hfill \\ +(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4})\mathrm{sin}\left(\frac{p_4p_2}{2}\right)\mathrm{sin}\left(\frac{p_3p_1}{2}\right)[\underset{i=1}{\overset{4}{}}\mathrm{ln}\stackrel{~}{p}_i^23\mathrm{ln}(\stackrel{~}{p}_1+\stackrel{~}{p}_3)^2]\},\hfill \end{array}$$
(4.61)
$`a^{\mathrm{diam}}`$ being given in eq. (2.23). Two comments are now in order. First, eqs. (4.59) and (4.60) show that the non-planar parts of the four-point diamagnetic diagrams exhibit IR divergences for configurations (4.55). Yet, the full diamagnetic contribution, given by eq. (4.61), is free of such singularities. There is, therefore, a cancellation mechanism at work among the non-commutative IR divergent contributions coming from the box, lynx and bubble diagrams. This mechanism is a consequence of gauge invariance being preserved at one-loop. Let us substantiate this statement. As we shall see in the next section, gauge invariance leads to the following Ward identity involving the three and four-point functions:
$$\begin{array}{c}p_4^{\mu _4}\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu }(p_1,p_2,p_3,p_4)=2\mathrm{sin}\left(\frac{p_1p_4}{2}\right)\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}(p_2p_3,p_2,p_3)\hfill \\ +\mathrm{\hspace{0.17em}2}\mathrm{sin}\left(\frac{p_2p_4}{2}\right)\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}(p_1p_3,p_3,p_1)\hfill \\ +\mathrm{\hspace{0.17em}2}\mathrm{sin}\left(\frac{p_3p_4}{2}\right)\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}(p_1p_2,p_1,p_2).\hfill \end{array}$$
(4.62)
Now, the r.h.s of this equation is finite for any of the configurations (4.55), as can be readily shown by substituting eq. (3.53) in it. Thus, if $`\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}(p_1,p_2,p_3,p_4)`$ had an IR diamagnetic divergence for momenta (4.55), it would have to be transverse with respect to $`p_4^{\mu _4}`$. However, there is no tensor transverse to $`p_4^{\mu _4}`$ that can be built with the tensors and functions that occur in eqs. (4.59) and (4.60). Secondly, the contribution to $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{diam},\mathrm{NP}}(p_1,p_2,p_3,p_4)`$ coming from the momentum region
$$|\stackrel{~}{p}_1||\stackrel{~}{p}_2||\stackrel{~}{p}_3||\stackrel{~}{p}_4||\stackrel{~}{p}_1+\stackrel{~}{p}_2||\stackrel{~}{p}_1+\stackrel{~}{p}_3||\stackrel{~}{p}_2+\stackrel{~}{p}_3|\theta \mathrm{\Lambda }_{\mathrm{IR}}0$$
(4.63)
is
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{diam},\mathrm{NP}}(p_1,p_2,p_3,p_4)4ia^{\mathrm{diam}}\mathrm{ln}(\theta \mathrm{\Lambda }_{\mathrm{IR}})^2\hfill \\ \times [(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3})\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\mathrm{sin}\left(\frac{p_3p_4}{2}\right)\hfill \\ +\left(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}\right)\mathrm{sin}\left(\frac{p_1p_4}{2}\right)\mathrm{sin}\left(\frac{p_3p_2}{2}\right)\hfill \\ +(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4})\mathrm{sin}\left(\frac{p_4p_2}{2}\right)\mathrm{sin}\left(\frac{p_3p_1}{2}\right)].\hfill \end{array}$$
(4.64)
As happened for the two-point and three-point functions, this logarithmic IR contribution can be retrieved from the UV divergent contribution $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{diam},\mathrm{UV}}(p_1,p_2,p_3,p_4)`$ by performing the identifications (2.29). This provides an explicit realization for the four-point function of the UV/IR mixing characteristic of quantum field theories on non-commutative Minkowski space-time. This UV/IR mixing works for each single Feynman diagram: note that one can apply the identifications (2.29) to eqs. (4.56) and (4.57) to obtain eqs. (4.59) and (4.60) for momenta verifying the conditions (4.63).
We now come to the paramagnetic part of the four-point function. It is given by the contribution quartic in the background field that comes from the diagram in fig. 4, the latter being given by eqs. (2.17) and (2.18). As already explained, the planar part of the diagram contains all UV divergences, while the non-planar part collects all integrals with non-commutative IR singularities. Projecting out from eq. (2.17) the contribution quartic in $`B_\mu `$ and using the results of section 2 for $`J(p)`$, we obtain the following UV divergences and leading non-commutative IR terms in the paramagnetic contribution to the four-point function:
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{param},\mathrm{UV}}(p_1,p_2,p_3,p_4)=\frac{ia^{\mathrm{param}}}{\epsilon }\hfill \\ \times [(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3})\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\mathrm{sin}\left(\frac{p_3p_4}{2}\right)\hfill \\ +\left(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}\right)\mathrm{sin}\left(\frac{p_1p_4}{2}\right)\mathrm{sin}\left(\frac{p_3p_2}{2}\right)\hfill \\ +(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4})\mathrm{sin}\left(\frac{p_4p_2}{2}\right)\mathrm{sin}\left(\frac{p_3p_1}{2}\right)]\hfill \end{array}$$
(4.65)
and
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{param},\mathrm{IR}}(p_1,p_2,p_3,p_4)=\mathrm{\hspace{0.17em}4}ia^{\mathrm{param}}\hfill \\ \times [(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3})\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\mathrm{sin}\left(\frac{p_3p_4}{2}\right)\mathrm{ln}(\stackrel{~}{p}_1+\stackrel{~}{p}_2)^2\hfill \\ +(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4})\mathrm{sin}\left(\frac{p_1p_4}{2}\right)\mathrm{sin}\left(\frac{p_3p_2}{2}\right)\mathrm{ln}(\stackrel{~}{p}_2+\stackrel{~}{p}_3)^2\hfill \\ +(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4})\mathrm{sin}\left(\frac{p_4p_2}{2}\right)\mathrm{sin}\left(\frac{p_3p_1}{2}\right)\mathrm{ln}(\stackrel{~}{p}_1+\stackrel{~}{p}_3)^2],\hfill \end{array}$$
(4.66)
$`a^{\mathrm{param}}`$ being as in eq. (2.23). Note that, despite the fact that the paramagnetic contribution to the four-point function is made of integrals $`J(p)`$, with $`p=p_1+p_2,p_2+p_3,p_1+p_3`$, whose non-planar parts become singular for $`\stackrel{~}{p}=0`$, $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{param},\mathrm{IR}}(p_1,p_2,p_3,p_4)`$ is finite for any configuration (4.55). This finiteness can be explained as a consequence of gauge invariance following the same lines as for the diamagnetic contribution: see paragraph below eq. (4.61). Also note that in the momentum region (4.63) the value of $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{param},\mathrm{IR}}(p_1,p_2,p_3,p_4)`$ can be obtained from the UV divergent contribution $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{param},\mathrm{UV}}(p_1,p_2,p_3,p_4)`$ by using the identifications (2.29), in accordance with UV/IR mixing.
It remains to study the mixed contribution to the four-point function. There are quite a few diagrams producing contributions quartic in the background field with both diamagnetic and paramagnetic vertices. It is straightforward to see, however, that their planar parts are all UV finite by power counting and that their non-planar parts are free of singularities for momenta (4.55), so we conclude that
$$\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{mixed},\mathrm{UV}}(p_1,p_2,p_3,p_4)=\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{mixed},\mathrm{IR}}(p_1,p_2,p_3,p_4)=0.$$
(4.67)
We are now ready to give the complete UV divergent contribution $`i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{UV}}(p_1,p_2,p_3,p_4)`$ to the four-point function. It is the sum of the UV divergent diamagnetic and paramagnetic contributions and takes the form
$$\begin{array}{c}i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{UV}}(p_1,p_2,p_3,p_4)=\frac{i}{\epsilon }\left(a^{\mathrm{diam}}+a^{\mathrm{param}}\right)\hfill \\ \times [(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _4}\eta _{\mu _2\mu _3})\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\mathrm{sin}\left(\frac{p_3p_4}{2}\right)\hfill \\ +\left(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4}\right)\mathrm{sin}\left(\frac{p_1p_4}{2}\right)\mathrm{sin}\left(\frac{p_3p_2}{2}\right)\hfill \\ +(\eta _{\mu _1\mu _3}\eta _{\mu _2\mu _4}\eta _{\mu _1\mu _2}\eta _{\mu _3\mu _4})\mathrm{sin}\left(\frac{p_4p_2}{2}\right)\mathrm{sin}\left(\frac{p_3p_1}{2}\right)]\hfill \end{array}$$
(4.68)
It follows that the MS counterterm (2.24) subtracts the UV divergences in the four-point function and that the beta function computed from the resulting renormalized four-point is as in eq. (2.26). As concerns the leading non-commutative IR contribution to the four-point function, it is finite after summing over diagrams and is given by the sum of eqs. (4.61) and (4.66):
$$i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{IR}}(p_1,p_2,p_3,p_4)=i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{diam},\mathrm{IR}}(p_1,p_2,p_3,p_4)+i\mathrm{\Gamma }_{\mu _1\mu _2\mu _3\mu _4}^{\mathrm{param},\mathrm{IR}}(p_1,p_2,p_3,p_4).$$
(4.69)
## 5 Gauge invariance
In this section we shall prove that gauge invariance for the one-loop renormalized theory as mathematically expressed by the equation
$$d^Dx\omega (x)D_\mu [B]\frac{\delta \mathrm{\Gamma }[B]}{\delta B_\mu (x)}=0$$
(5.70)
holds, where $`i\mathrm{\Gamma }[B]=\mathrm{ln}Z[B]`$ and $`Z[B]`$ is given in eq. (2.8). In this regard it is worth noting that, to the best of our knowledge, it remains an open question whether non-logarithmic non-commutative IR divergences are compatible with gauge invariance and finiteness. Here we see that non-logarithmic IR divergences preserve gauge invariance.
We have shown in sections 2, 3 and 4 that, to subtract the UV divergences in the two, three and four-point functions, it is enough to add the counterterm (2.24)
$$\delta S[B]=\frac{1}{4\epsilon }\left(a^{\mathrm{diam}}+a^{\mathrm{param}}\right)d^DxF^{\mu \nu }[B]F_{\mu \nu }[B].$$
(5.71)
This counterterm satisfies eq. (5.70) and does not change the non-commutative IR behaviour of the Green functions. Hence, the renormalized two, three and four-point functions have the same non-commutative IR behaviour as the regularized ones, given in eqs. (2.30), (3.53) and (4.69). We must check that these expressions are consistent with the Ward identity (5.70). To this end, we note that eq. (5.70) implies
$`p^\nu \mathrm{\Pi }_{\mu \nu }(p)=0`$ (5.72)
$`p_3^{\mu _3}\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}(p_1,p_2,p_3)=2i\mathrm{sin}\left({\displaystyle \frac{p_1p_2}{2}}\right)\left[\mathrm{\Pi }_{\mu _1\mu _2}(p_1)\mathrm{\Pi }_{\mu _1\mu _2}(p_2)\right],`$ (5.73)
together with eq. (4.62). It is apparent that the non-commutative IR divergent contribution $`\mathrm{\Pi }_{\mu \nu }^{\mathrm{IR}}(p)`$ to the vacuum polarization tensor in eq. (2.30) verifies the identity (5.72). It is also clear that, neglecting terms which remain finite for some $`i`$ $`(i=1,2,3)`$ in the limit $`\stackrel{~}{p}_i0`$, we have
$$\begin{array}{c}p_3^{\mu _3}\mathrm{\Gamma }_{\mu _1\mu _2\mu _3}^{\mathrm{IR}}(p_1,p_2,p_3)\frac{2}{\pi ^2}(p_1p_2)\left[\frac{(\stackrel{~}{p}_1)_{\mu _1}(\stackrel{~}{p}_1)_{\mu _2}}{\stackrel{~}{p}_1^2}\frac{(\stackrel{~}{p}_2)_{\mu _1}(\stackrel{~}{p}_2)_{\mu _2}}{\stackrel{~}{p}_2^2}\right]\hfill \\ 2i\mathrm{sin}\left(\frac{p_1p_2}{2}\right)\left[\mathrm{\Pi }_{\mu _1\mu _2}^{\mathrm{IR}}(p_1)\mathrm{\Pi }_{\mu _1\mu _2}^{\mathrm{IR}}(p_2)\right],\hfill \end{array}$$
in agreement with the Ward identity (5.73) above. As concerns the leading non-commutative IR behaviour of the four-point function, its consistency with gauge invariance has already been thoroughly checked in section 4.
We end this section writing the effective action which reproduces the logarithmic contributions in the non-commutative IR momentum regions (2.28), (3.48) and (4.63) to the vacuum polarization tensor, three-point function and four-point function. From the results derived in sections 2, 3 and 4 it is clear that it has the form
$$\frac{1}{4}(a^{\mathrm{diam}}+a^{\mathrm{param}})\mathrm{ln}(\theta \mathrm{\Lambda }_{\mathrm{IR}})^2d^4x(F^{\mu \nu }[B]F_{\mu \nu }[B],$$
(5.74)
with the already known result (2.23) for $`a^{\mathrm{diam}}`$ and $`a^{\mathrm{param}}`$. In accordance with ref. , it is dual to the UV counterterm (5.71) under the identifications (2.29).
## 6 Conclusions and suggestions
In this paper we have computed the complete UV and non-commutative IR divergent contributions to the one-loop effective action of $`U(1)`$ gauge theory on non-commutative Minkowski space-time in a background field. The UV divergences arise from the planar parts of the one-loop diagrams, and come from virtual quanta with arbitrarily high momenta moving around the loop of the diagrams. The non-commutative IR divergent contributions are also produced by high-momentum virtual quanta, but occur in the non-planar parts of the diagrams. The vacuum polarization has quadratic and logarithmic non-commutative IR divergent contributions, while the three-point function has linear and logarithmic contributions, and the four-point function only has logarithmic contributions . Logarithmic contributions add to an overall IR finite contribution for the three and four-point functions, while this is not so for the vacuum polarization tensor. The overall logarithmic non-commutative IR behaviour of the three Green functions is dual to the overall UV divergent behaviour in the sense that the identifications (2.29) transform one into another.
The main result presented in this article is the following. We have shown that the paramagnetic contributions to the effective action dominate over the diamagnetic contributions at very high momentum. Here we say that we are in the high-momentum or high-energy region if the momentum is large as measured against the non-commutative energy scale $`|\theta |^{1/2}`$. It turns out that, in the high-momentum regime, paramagnetic contributions give rise to anti-screening of the coupling constant, whereas diamagnetic contributions produce screening. The combined effect explains the negative sign of the beta function and its actual value. We recall that, while for a Lorentz invariant theory paramagnetism always comes with charge anti-screening and diamagnetism with charge screening , this is not necessarily so if the theory lacks Lorentz invariance, a circumstance that occurs if space is non-commutative. Yet the net result is analogous to that for Yang-Mills theory on commutative Minkowski space-time. The reason for this is that the beta function can be computed from the UV divergences in the vacuum polarization tensor, provided gauge invariance holds. These divergences occurring only in the non-planar part implies that they are independent of $`\theta ^{\mu \nu }`$ and do not break Lorentz invariance.
We have also shown that the quadratic and linear non-commutative IR divergences that exhibit the vacuum polarization tensor and three-point function have a purely diamagnetic origin. By contrast, the leading non-commutative IR logarithmic contributions to the vacuum polarization tensor, the three-point and the four-point functions receive both paramagnetic an diamagnetic contributions. The duality of logarithmic contributions to UV divergences discussed above does not need to sum over diamagnetic and paramagnetic contributions but holds for each of them separately.
Let us now look at the renormalized theory and consider the non-commutative IR regime (4.63). As shown in section 5, the leading non-commutative IR logarithmic contribution to the effective action in this region has the form (5.74). This contribution prevails over logarithmic contributions $`\mathrm{ln}(p^2/\mu ^2)`$, if $`p^2`$ is of the same order of magnitude as the squared renormalization scale $`\mu ^2`$. This leads us to claim that, in this regime, eq. (5.74) can be thought of as giving rise to an effective IR renormalization of the renormalized coupling constant of the theory $`g^2(\mu )`$, thus introducing an IR effective coupling $`g_{\mathrm{IR}}^2`$ given by
$$\frac{1}{g_{\mathrm{IR}}^2}=\frac{1}{g^2(\mu )}(a^{\mathrm{diam}}+a^{\mathrm{param}})\mathrm{ln}(\theta \mathrm{\Lambda }_{\mathrm{IR}}\mu )^2,$$
with $`a^{\mathrm{diam}}`$ and $`a^{\mathrm{param}}`$ as in eq. (2.23). In other words, in the non-commutative IR domain (4.63), the dominant logarithmic contribution to the 1PI one-loop Green functions of the theory are given by their tree-level expressions at finite $`\theta `$ but with $`g^2(\mu )`$ replaced with $`g_{\mathrm{IR}}^2`$. Note also that, in this non-commutative IR region, the paramagnetic contributions produce a screening of the charge, whereas the diamagnetic contributions anti-screen the charge. This does not come as a surprise since we know that when Lorentz invariance is lost there are media where paramagnetism comes with screening and diamagnetism with anti-screening , the dpeed of light being not equal to 1. This is precisely what happens in the case at hand. As pointed out in ref. , the quadratic non-commutative IR divergent behaviour of the vacuum polarization tensor gives rise to a modification of the dispersion relation for photon polarizations parallel to $`\stackrel{~}{p}^\mu `$. Indeed, let $`p^\mu =(E,\stackrel{}{p})`$, with $`\stackrel{}{p}=(p^1,p^2,p^3)`$, and take a photon polarized along the direction defined by $`\stackrel{~}{p}^\mu `$. We then have
$$E^2=(p^3)^2+\stackrel{}{P}^2\left(1\frac{2}{\pi ^2}\frac{g^2}{\theta ^2|\stackrel{}{P}|^4}\right),$$
(6.75)
where $`\stackrel{}{P}=(p^1,p^2,0)`$. This equation leads to a photon speed $`\stackrel{}{v}_g=\stackrel{}{}_\stackrel{}{p}E`$, with modulus greater than 1 in the region $`2g^2\pi ^2\theta ^2|\stackrel{}{P}|^4`$ where perturbation theory is valid. Whether this signals a true instability of the theory or is an artifact of perturbation theory remains an open question .
One may wonder if the results presented here for $`U(1)`$ gauge theory generalize to the $`U(N)`$ case. In this regard we note that, although in the non-commutative IR domain no new divergent dependence on the external momenta appears, the colour structures of the UV and the non-commutative IR divergent parts are different . This indicates that UV/IR duality in this case is more involved and deserves further investigation. Finally, it will be worth exploring whether there exists a connection between the results presented here and the dipole structures that occur in non-commutative quantum field theories .
## Acknowledgement
The authors are grateful to E. López for pointing them out that the relative sign in the dispersion relation (6.75) in a previous version of this paper was wrong. They also thank CICyT, Spain for financial support through grant No. PB98-0842.
Figure 1: Feynman rules of $`U(1)`$ gauge theory on non-commutative Minkowski space-time in the Feynman background gauge. The last vertex is the only one of paramagnetic type.
Figure 2: One-loop Feynman diagram mixing diamagnetic and paramagnetic vertices contributing to the vacuum polarization tensor.
Figure 3: One-loop diamagnetic Feynman diagrams contributing to the vacuum polarization tensor. The diagrams include planar and non-planar terms.
Figure 4: One-loop paramagnetic diagram contributing to the 1PI Green function $`F_{\mu \nu }(p)F_{\lambda \rho }(p)`$. The diagram includes planar and non-planar terms.
Figure 5: One-loop 1PI diamagnetic diagrams contributing to the three-point function.
Figure 6: One-loop mixed diagrams contributing to the three-point function.
Figure 7: One-loop 1PI diamagnetic diagrams contributing to the four-point function.
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# FUTURE HIGH 𝐐^𝟐 DEEP INELASTIC SCATTERING AT HERA
## 1 Introduction
The electron-proton collider HERA started operation in the summer of 1992. The proton beam energy was 820 GeV while the electron beam energy was 26.7 GeV and was later raised to 27.5 GeV. In the years 1992-1997 H1 and ZEUS each collected a luminosity of $``$1 pb<sup>-1</sup> using electron beams and $``$50 pb<sup>-1</sup> using positron beams. These data have extended the kinematic range covered by deep inelastic scattering, DIS, measurements by two orders of magnitude in both $`Q^2`$, the four-momentum transfer squared, and $`x`$, the fraction of the proton four-momentum carried by the struck quark. These data have been used to dertermine the proton structure function, make measurements which test the electroweak Standard Model, SM, and the theory of the strong interactions, QCD, in both neutral current, NC, and charged current, CC, DIS. Jet analyses in DIS and photoproduction have been used to address fundamental issues in QCD. The observation of diffraction in DIS has led to a careful investigation of the transition from the kinematic region in which perturbative QCD is valid to the region where phenomenological models based on Regge theory must be applied (see for example $`^\mathrm{?}`$ and $`^\mathrm{?}`$ and references therein).
During the running period August 1998 to April 1999 $``$20 pb<sup>-1</sup> of $`e^{}p`$ data were delivered with a proton beam energy of 920 GeV. HERA is now delivering $`e^+p`$ collisions with a proton beam energy of 920 GeV. By the end of running in September 2000 H1 and ZEUS will each have an $`e^+p`$ data set of $`100`$ pb<sup>-1</sup>. The data collected by H1 and ZEUS will be used to study the dependence of the NC and CC DIS cross sections on the charge of the lepton beam.
The HERA experiments will continue to take data until September 2000 when a long, 9 month, shutdown is scheduled. The shutdown will be used to upgrade the HERA accelerator and the collider detectors. The HERA luminosity will be increased by a factor of five and longitudinal lepton beam polarisation ($`70\%`$) will be provided for ZEUS and H1. Over a six year running period it is anticipated that a total luminosity of 1000 pb<sup>-1</sup> will be delivered $`^\mathrm{?}`$. The physics motivation for this major upgrade programme is discussed in detail in reference $`^\mathrm{?}`$.
## 2 Physics at HERA after the Upgrade
Following the HERA upgrade the proton will be probed using each of the four possible combinations of lepton beam charge and polarisation. The combination of high luminosity and polarisation will lead to a rich and diverse programme of measurements which can only be sketched below using a few examples.
### 2.1 Proton Structure
The large data volume will allow $`F_2^{\mathrm{NC}}`$ to be extracted with an accuracy of $``$3% over the kinematic range $`2\times 10^5<x<0.7`$ and $`2\times 10^5<Q^2<5\times 10^4\mathrm{GeV}^2`$ $`^\mathrm{?}`$. If QCD evolution codes which go beyond next to leading order become available and a careful study of the dependence of the systematic errors on the kinematic variables is made it will be possible to determine $`\alpha _\mathrm{S}`$ from the scaling violations of $`F_2^{\mathrm{NC}}`$ with a precision of $`0.003`$. The gluon distribution will also be determined from such a fit with a precision of $``$ 3% for $`x=10^4`$ and $`Q^2=20\mathrm{GeV}^2`$.
The combination of high luminosity and high charm tagging efficiency transforms the measurement of the charm contribution to $`F_2^{\mathrm{NC}},F_2^{cc}`$ $`^\mathrm{?}`$. The precision will be sufficient to allow a detailed study of the charm production cross section to be made. The lifetime tag provided by the silicon micro-vertex detector allows the tagging of $`b`$-quarks and the determination of the ratio of the beauty contribution to $`F_2^{\mathrm{NC}}`$, $`F_2^{bb}`$, to $`F_2^{cc}`$.
In the quark parton model CC DIS is sensitive to specific quark flavours. The $`e^+p`$ CC DIS cross section is sensitive to the $`d`$\- and $`s`$-quark parton densities and the $`\overline{u}`$\- and $`\overline{c}`$-anti-quark densities, while the $`e^{}p`$ CC DIS cross section is sensitive to the $`u,c,\overline{d}`$ and $`\overline{s}`$ parton density functions. With the large CC data sets expected following the upgrade it will be possible to use $`e^\pm p`$ CC data to determine the $`u`$\- and $`d`$\- quark densities. Further, by identifying charm in CC DIS it will be possible to determine the strange quark contribution to the proton structure function $`F_2^{\mathrm{NC}}`$ with an accuracy of between 15% and 30% $`^\mathrm{?}`$.
### 2.2 Tests of the Electroweak Standard Model
The high luminosity provided by the upgrade will allow access to low cross section phenomena such as the production of real $`W`$-bosons. The SM cross section for the process $`epeWX`$ is $``$ 1 pb $`^\mathrm{?}`$ which, combined with an acceptance of $`30\%`$, gives a sizeable data sample for a luminosity of 1000 pb<sup>-1</sup>. The production of the $`W`$-boson at HERA is sensitive to the non-abelian coupling $`WW\gamma `$ $`^\mathrm{?}`$. The sensitivity of HERA to non-SM couplings is comparable to the sensitivities obtained at LEP and at the Tevatron and complementary in that at HERA is predominantly sensitive to the $`WW\gamma `$ vertex, independent of assumptions about the nature of the $`WWZ`$ vertex.
The full potential of electroweak tests at HERA will be realised through measurements using polarised lepton beams $`^\mathrm{?}`$. Within the SM NC and CC DIS cross sections may be written in terms of $`\alpha `$, $`M_W`$ and $`m_t`$ together with the mass of the $`Z`$ boson, $`M_Z`$, and the mass of the Higgs boson, $`M_H`$. In order to test the consistency of the theory we may fix the values of $`\alpha `$ and $`M_Z`$ to those obtained at LEP or elsewhere and use measurements of the NC and CC DIS cross sections to place constraints in the $`M_W`$, $`m_t`$ plane for fixed values of $`M_H`$. The SM is consistent if the values of the parameters $`M_W`$ and $`m_t`$ obtained agree with the values determined in other experiments. Combining NC and CC data corresponding to a luminosity of 1000 pb<sup>-1</sup>, recorded with a lepton beam polarisation of 70%, with a top mass measurement from the Tevatron with a precision of $`\pm 5`$ GeV yields a measurement of $`M_W`$ with an error of $`60`$ MeV $`^\mathrm{?}`$.
The sensitivity of NC DIS to lepton beam polarisation is shown in figure 1(a). The figure shows the ratio of the full NC cross section to the cross section obtained in the single photon exchange approximation. The strong polarisation dependence of the NC cross section can be used to extract the NC couplings of the light quarks. In such an analysis the CC cross section may be used to reduce the sensitivity of the results to uncertainties in the PDFs $`^\mathrm{?}`$. The precision of the results obtained depends strongly on the degree of polarisation of the lepton beam as shown in figure 1(b). The figure shows the anticipated error on the vector and axial-vector couplings of the $`u`$-quark, $`v_u`$ and $`a_u`$ respectively, obtained in a fit in which $`v_u`$ and $`a_u`$ are allowed to vary while all other couplings are fixed at their SM values. With a luminosity of 250 pb<sup>-1</sup> per charge, polarisation combination and taking the vector and axial-vector couplings of the $`u`$\- and $`d`$-quarks as free parameters gives a precision of 13%, 6%, 17% and 17% for $`v_u`$, $`a_u`$, $`v_d`$ and $`a_d`$ respectively. By comparing these results with the NC couplings of the $`c`$\- and $`b`$-quarks obtained at LEP a stringent test of the universality of the NC couplings of the quarks will be made.
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# Quantum mechanical histories and the Berry phase
## 1 Introduction
The consistent histories formulation of quantum mechanics focuses on the temporally ordered properties of physical systems: these are known as histories. There are two important structural features of quantum mechanical history theories. The first is that there does not exist a probability measure in the space of all histories: there exists interference between pairs of histories.
The second is their non-trivial temporal structure, that allows a differentiation between the kinematical and the dynamical aspects of time. This is present in the quantum temporal logic formulation of consistent histories ; this allows a description of continuous-time histories , in which there exist distinct generators of time translation according to whether they refer to dynamical or kinematical features of the histories .
In this paper, we establish that this distinction is mirrored in a differentiation that is well known in standard quantum thory: the one between the dynamical phase due to Hamiltonian evolution and the geometric phase of Berry . What is more, the geometric phase manifests itself strongly in the probabilistic structure of histories: it is the basic building block of the interference phase between pairs of histories. These results are established in an elementary fashion by the study of fine-grained, continuous-time histories; they can be then suitably generalised.
Now, the appearance (and measurability) of the geometric phase in the time evolution of quantum systems is arguably one of the most important structural features of quantum theory. Berry showed that when a system undergoes a cyclic evolution, due to an adiabatic change of parameters in the Hamiltonian, a contribution in the phase appears, that is purely geometric. In particular, the phase contribution does not depend on the details of the dynamics but only on the loop that was transversed by the system in the parameter space.
It was soon realised that the Berry phase is the holonomy of a U(1) connection on the parameter space. In fact it can be generalised for any kind of unitary evolution on the Hilbert space, since it arises by a natural U(1) connection on the projective Hilbert space.
The geometric phase is a measurable quantity, that does not formally correspond to a self-adjoint operator. Furthermore it provides a paradigm and a motivation for investigations of topological phenomena in quantum theory, as it highlights the natural appearance of gauge structures in the quantum formalism.
The key point of Berry’s result however, that sets the subject in the foundations of quantum theory, is the following: the Berry phase has no analogue in the language of probability theory.
A probabilistic theory for a physical system—either classical or quantum—has as basic notions observables, propositions and states that are represented by suitable mathematical objects. In classical probability theory observables are functions on a space $`\mathrm{\Omega }`$, propositions correspond to measurable subsets of $`\mathrm{\Omega }`$ and states to probability distributions. In quantum theory these probabilistic concepts are also fundamental: they are represented by Hilbert space objects. These are self-adjoint operators for observables, projection operators for propositions and density matrices for states.
However, the standard quantum mechanical formalism refers to properties of the system at a single moment of time: it assigns probabilities to possible instantaneous events and studies the evolution of these single-time probabilities. In this context, the phase of a Hilbert space vector is not physically relevant, as it does not enter the single-time probability assignment. When this phase is ignored, quantum theory is only a generalisation of probability theory, with main difference the non-distributivity of the lattice of propositions or equivalently the non-commutativity of the algebra of observables.
This is the attitude taken by approaches to quantum theory that attempt to write an axiomatic framework without assuming a priori the existence of a Hilbert space, for example, the $`C^{}`$–approach, quantum logic schemes or the operational approach to quantum theory.
The existence of the Berry phase, as a measurable quantity, shows that the single-time probabilistic description does not exhaust the physical content of quantum theory. The Berry phase appears in distinction to the well-known phases of unitary evolution that are generated by a Hamiltonian. Its nature is purely kinematical and it is a manifestation of the non-trivial topology of the space of pure quantum states, since it appears naturally when we view the Hilbert space $`H`$ of quantum theory as a complex line bundle over the projective Hilbert space $`PH`$ . The Berry phase is then the holonomy of the natural connection of this bundle (i.e. the connection induced by the inner product).
It needs to be emphasised that this bundle structure is irrelevant to any probabilistic aspects of quantum theory. In other words, in the unitary time evolution of quantum theory there appears an extra phase due to the topological structure of the theory; it has no intuitive physical explanation and it has no classical analogue—either in classical mechanics <sup>1</sup><sup>1</sup>1The Hannay angle is an analogue in classical mechanics. But this appears whenever certain degrees of freedom can be ignored due a symmetry, whereas the wave function is assumed to give a complete description of the quantum system. or in classical probability theory.
In the single-time description of a quantum system the geometric phase is lost. Hence it is rather difficult to understand its physical meaning in standard quantum theory. However, the importance of geometric phase is more clear in a quantum theory that is based on histories. A history is defined at different moments of time, in distinction to standard single-time quantum theory. Such a formulation is provided by the consistent histories approach to quantum theory.
This approach was developed as a realist interpretational scheme for quantum theory . As such, it suffers from the generic problems of such schemes, i.e., contextuality of predictions about properties of the physical system . Nonetheless, it provides a new insight in understanding the appearance of the Berry phase in quantum theory, in a manner independent of any particular interpretational scheme one may choose to employ.
The basic object of the histories formalism is a history, i.e., a sequence of time-ordered propositions about properties of the physical system. It corresponds to different possible scenaria of the system. The main feature, that distinguishes quantum mechanical histories from the ones appearing in the theory of stochastic processes (classical probability theory), is that the probabilities for histories do not satisfy the additivity condition.
$$p(\alpha \beta )=p(\alpha )+p(\beta ),$$
(1)
where $`\alpha `$ and $`\beta `$ are mutually exclusive scenaria. This is due to the fact that quantum theory is based on amplitudes rather than probability measures, and it further implies the existence of interference between histories.
The corresponding information, together with the probabilities, is encoded in an object called the decoherence functional. This object incorporates the kinematics, the dynamics and the initial condition of the physical system.
In our effort to identify the role of the Berry phase in the histories scheme, we arrived at a surprisingly simple result: the geometric phase is the main building block of the decoherence functional. Hence, interference between histories is ultimately to be attributed to the presence of the geometric phase. Moreover, we showed that the distinction between geometric phase and the dynamical phase of canonical quantum theory –i.e the one appearing due to Hamiltonian time evolution– is a manifestation of the temporal structure of history theories: the existence of two laws of time transformation each corresponding to the causal/kinematical and dynamical notions of time .
## 2 The geometric phase
The simplest way to demonstrate the origin of the Berry phase is in the context of differential geometry. To this end, let us take the complex Hilbert space $`H`$ to be finite dimensional ($`H=𝐂^{n+1}`$). The inner product $`<z|w>=\overline{z}_aw^a`$ gives a metric $`ds^2=d\overline{z}_adz^a`$ (where $`a`$ runs from $`0`$ to $`n`$ and $`z`$ refers to coordinates with respect to a basis), from its real part, and a symplectic form $`\omega =d\overline{z}_adz^a`$ on $`H`$, from its imaginary part.
The metric induces the standard metric to the unit sphere $`S^{2n+1}`$ of all normalised vectors. The unit sphere is a $`U(1)`$ principal bundle over the projective Hilbert space $`PH`$, the space of rays; this structure is known as the Hopf bundle. An element of $`PH`$ is represented by $`[\psi ]`$, the equivalence class of all normalised vectors that differ from the normalised vector $`|\psi `$ only with respect to a phase. The metric on $`S^{2n+1}`$ induces a metric on $`PH`$, defined as
$$ds^2(PH)=\frac{1}{1+\overline{w}_aw^a}d\overline{w}_adw^a$$
(2)
and an one-form $`A=i\overline{w}_adw^a`$. Here, we have defined coordinates such that for $`1an`$, $`w^a=z^a/z^0`$.
In particular, the one-form $`A`$ is a $`U(1)`$ connection form for the Hopf bundle, and it is called the Berry connection; its curvature is equal to the projection of the symplectic form in $`PH`$, modulo $`i`$. It may be written in a coordinate independent way as $`A=i\psi |d|\psi `$.
We assume an arbitrary unitary time evolution $`U(s)`$ on the Hilbert space $`H`$, and we take an initial vector $`|\psi _0`$ at time $`t=0`$. The curve
$$U(s)|\psi _0:=|\psi (s)$$
(3)
projects to a curve $`[\psi (s)]`$ on the projective Hilbert space $`PH`$.
If we further assume that $`U(s)`$ is such that at time $`t`$, $`[\psi (t)]=[\psi (0)]`$, i.e. we have a loop $`\gamma `$ on the projective space, then the phase that is transversed on the U(1) fiber is equal to
$$e^{_0^t𝑑s\psi (s)|\frac{d}{ds}iH(s)|\psi (s)}:=e^{iS[\psi (.)]},$$
(4)
where we wrote $`H(s)=U^1(s)\dot{U}(s)`$ and $`S`$ is the action out of which the Schrödinger equation is derived. The second term is a time dependent angle due to time evolution.
However, the first term is purely geometrical; it depends only on the transversed loop, and is equal to the holonomy of the Berry connection
$$e^{i_\gamma A}=\mathrm{exp}\left(_\gamma \psi |d\psi \right).$$
(5)
Note that the Berry phase does not change if we take different representatives $`|\psi `$ for the equivalence class $`[\psi ]`$.
The geometric phase may also be defined for open paths by exploiting the metric structure on $`PH`$ . It allows us to form a loop from any path on the projective Hilbert space, by joining its endpoints with a geodesic. The geometric phase of the loop thus constructed is defined to be equal to the geometric phase associated to the open path. Hence if $`\gamma =[\psi (.)]`$ is a path on $`PH`$, its associated geometric phase is proved to equal
$$e^{i\theta _g[\gamma ]}=\mathrm{exp}\left(_{t_i}^{t_f}𝑑t\psi (t)|\dot{\psi }(t)\right)\psi _i|\psi _f.$$
(6)
This expression is meaningful only if the endpoints are not orthogonal.
Hence, the Berry phase is strongly related with geometric and topological structures of the Hilbert space of quantum theory. These geometric structures are physically relevant because of Born’s probability interpretation: the single-time expectation values for observables do not change with phase transformations of the Hilbert space vector $`|ze^{i\varphi }|z`$.
## 3 Histories
A history is defined as a sequence of projection operators $`\alpha _{t_1},\mathrm{},\alpha _{t_n}`$, and it corresponds to a time-ordered sequence of propositions about the physical system. The indices $`t_1,\mathrm{},t_n`$ refer to the time a proposition is asserted and have no dynamical meaning. Dynamics are related to the Hamiltonian $`H`$, which defines the one-parameter group of unitary operators $`U(s)=e^{iHs}`$.
A natural way to represent the space of all histories is by defining a history Hilbert space $`𝒱:=_{t_i}_{t_i}`$, where $`_{t_i}`$ is a copy of the standard Hilbert space, indexed by the moment of time to which it corresponds. A history is then represented by a projection operator on $`𝒱`$. This construction has the merit of preserving the quantum logic structure and highlighting the non-trivial temporal structure of histories . Furthermore, one can also construct a Hilbert space $`𝒱`$ for continuous-time histories by a suitable definition of the notion of the tensor product.
Furthermore, to each history $`\alpha `$ we may associate the class operator $`C_\alpha `$ defined by
$$C_\alpha =U^{}(t_n)\alpha _{t_n}U(t_n)\mathrm{}U^{}(t_1)\alpha _{t_1}U(t_1).$$
(7)
It is important to note that time appears in two distinct places in the definition of the class operator $`C_\alpha `$: as the argument of the Heisenberg time evolution and as the parameter identifying the time at which a proposition is asserted. In what follows, we will show that this distinction is strongly related to the distinction between geometric and dynamical phase.
The decoherence functional is defined as a complex-valued function of pairs of histories: i.e. a map $`d:𝒱\times 𝒱𝐂`$. For two histories $`\alpha `$ and $`\alpha ^{}`$ it is given by
$$d(\alpha ,\alpha ^{})=Tr\left(C_\alpha \rho _0C_\alpha ^{}^{}\right)$$
(8)
The standard interpretation of this object is that when $`d(\alpha ,\alpha ^{})=0`$ for $`\alpha \alpha ^{}`$ in an exhaustive and exclusive set of histories <sup>2</sup><sup>2</sup>2 By exhaustive we mean that at each moment of time $`t_i`$ $`_{\alpha _{t_i}}\alpha _{t_i}=1`$ and by exclusive that $`\alpha t_i\beta _{t_i}=\delta _{\alpha \beta }`$. Note that by $`\alpha `$ we denote both the proposition and the corresponding projector., then one may assign a probability distribution to this set as $`p(\alpha )=d(\alpha ,\alpha )`$. The value of $`d(\alpha ,\beta )`$ is, therefore, a measure of the degree of interference between the histories $`\alpha `$ and $`\beta `$.
## 4 The geometric phase for histories
We now consider a time interval $`[t_0,t_f]`$ and a history with $`n+1`$ time steps $`\alpha _{t_0},\alpha _{t_1},\mathrm{}\alpha _{t_f}`$. We assume that the projectors are fine-grained, which means that they correspond to elements of the projective Hilbert space
$$\alpha _{t_i}=|\psi _{t_i}\psi _{t_i}|.$$
(9)
We first set the Hamiltonian equal to zero. The trace of the class operator $`C_\alpha `$ equals
$$TrC_\alpha =\psi _{t_0}|\psi _{t_n}\psi _{t_1}|\psi _{t_0}\psi _{t_2}|\psi _{t_1}\mathrm{}\psi _{t_n}|\psi _{t_{n1}}$$
(10)
and it is non-zero provided there is no value of $`i`$, for which the vector $`|\psi _{t_i}`$ is orthogonal to $`|\psi _{t_{i1}}`$.
Next, we assume that $`\mathrm{max}|t_jt_{j1}|=\delta t`$, and we choose the number of time steps $`n`$ very large, so that $`\delta tO(n^1)`$. Then $`|\varphi _{t_j}`$ approximates a path $`[\varphi (t)]`$ on $`PH`$. Hence,
$`\mathrm{log}TrC_\alpha `$ $`=`$ $`\mathrm{log}\psi _{t_0}|\psi _{t_n}+{\displaystyle \underset{i=1}{\overset{n}{}}}\mathrm{log}\psi _{t_i}|\psi _{t_{i1}}`$ (11)
$`=`$ $`\mathrm{log}\psi _{t_0}|\psi _{t_n}+{\displaystyle \underset{i=1}{\overset{n}{}}}\mathrm{log}\left(1\psi _{t_i}|\psi _{t_i}\psi _{t_{i1}}\right)`$
and the limit of large $`n`$ yields
$$\mathrm{log}TrC_\alpha =\mathrm{log}\psi _{t_0}|\psi _{t_n}\psi _{t_i}|\psi _{t_i}\psi _{t_{i1}}+O((\delta t)^2)$$
(12)
As $`\delta t0`$ the sum in the right-hand side converges to a Stieljes integral $`_{t_i}^{t_f}𝑑t\psi (t)|\dot{\psi }(t)`$. Hence for a continuous path we take
$$TrC_\alpha =e^{i\theta _g[\psi ()]}$$
(13)
Therefore, the map $`\alpha TrC_\alpha `$ assigns to each fine-grained “continuous-time” history $`\alpha `$ its corresponding Berry phase. In fact, the paths $`\psi ()`$ need not be continuous; it suffices that the Stieljes integral is defined.
Furthermore, one may use the above result to define the Berry phase, associated to a general coarse-grained history. Hence, if $`\alpha =(\widehat{\alpha }_{t_1},\mathrm{},\widehat{\alpha }_{t_n})`$ is a history, then we may write
$`\alpha Tr\left(\widehat{\alpha }_{t_n}\mathrm{}\widehat{\alpha }_{t_2}\widehat{\alpha }_{t_1}\right).`$ (14)
This defines a map from $`𝒱`$ to the complex numbers, that assigns to each history its corresponding geometric phase. In particular, if we decompose the projector $`\widehat{\alpha }_{t_i}`$ with respect to an orthonormal basis in the subspace, in which it projects
$$\widehat{\alpha }_{t_i}=\underset{r}{}|\psi _{t_i}^r\psi _{t_i}^r|,$$
(15)
we may then write the geometric phase for the coarse-grained histories as
$$\underset{r_1,\mathrm{}r_n}{}e^{i\theta _g[\psi _{r_1\mathrm{}r_n}()}]$$
(16)
In the continuum limit this can be written, suggestively, as a sum over all fine-grained paths $`\psi ()`$ compatible with the coarse-grained history $`\alpha `$
$$\underset{\psi ()\alpha }{}e^{i\theta _g[\psi ()]}$$
(17)
In view of this linearity, the map that assigns to each history the corresponding Berry phase can be described by a functional on $`𝒱`$. When $`𝒱`$ with a tensor product of single-time Hilbert spaces, this linear functional is naturally induced by the tensor product construction.
We must note here that, our definition of the geometric phase is structurally distinct from the standard one. The latter refers to the evolution of a state under a dynamical law. In histories formalism, the geometric phase is defined on observables, or, more precisely, on possible scenaria for the physical system. There is, therefore, no need to make any assumption about the dynamics: this definition of geometric phase makes sense even if the dynamics is non-unitary.
## 5 The structure of the decoherence functional
The standard form of the decoherence functional incorporates the histories $`\alpha `$ by means of the operator $`\widehat{C}_\alpha `$. This suggests an expression for the decoherence functional that can be written in terms of the geometric phase.
To this end, let us assume two “continuous-time” histories, which we shall denote as $`\alpha _{\varphi (.)}`$ and $`\alpha _{\psi (.)}`$. From Eq. (13), and for the decoherence functional written for vanishing Hamiltonian, we take
$`d(\alpha _{\psi (.)},\alpha _{\varphi (.)})`$ $`=`$ (18)
$`=\varphi (t_i)|\rho _0|\psi (t_i)\psi (t_f)|\varphi (t_f)\times \mathrm{exp}\left({\displaystyle _{t_i}^{t_f}}𝑑t\psi (t)|\dot{\psi }(t)+{\displaystyle _{t_i}^{t_f}}𝑑t\dot{\varphi }(t)|\varphi (t)\right)`$
The two histories form a loop on $`PH`$, provided that their endpoints coincide. For example, this is the case where the density matrix $`\rho _0`$ is pure, and hence equal to an one-dimensional projector that could be considered as part of the history. From Eq. (14) we conclude that, the value of the decoherence functional is the Berry phase, associated to this loop.
When the Hamiltonian is included, we find
$`d(\alpha _{\psi (.)},\alpha _{\varphi (.)})=\varphi (t_i)|\rho _0|\psi (t_i)\psi (t_f)|\rho _f|\varphi (t_f)e^{iS[\psi (.)]iS^{}[\varphi (.)]},`$ (19)
where the action operator $`S`$ is given by the expression
$`S[\varphi (.)]={\displaystyle _{t_i}^{t_f}}dt\varphi (t)|i{\displaystyle \frac{d}{dt}}H|\varphi (t)`$ (20)
Hence the phase change on the Hopf bundle enters the decoherence functional at the level of the most general fine-grained histories.
Let us now note the following:
First, the appearance of the action is contingent upon the dynamics given by a Hamiltonian. One may consider more general dynamics: they are incorporated in the decoherence functional through the map $`\widehat{\alpha }_t\widehat{\alpha }_t(t)`$ that assigns to each Schrödinger-picture projector $`\widehat{\alpha }_t`$, a corresponding Heisenberg-picture one $`\widehat{\alpha }_t(t)`$, at time $`t`$. In full generality, it suffices that the dynamics is generated by an one-parameter family of automorphisms of the algebra of operators on the Hilbert space $``$ (it does not even need to be an one-parameter group). Hence, even though the expression involving the action is suggestive and simple, it is not as fundamental and general as Eq. (18), which expresses the decoherence functional in terms of the Berry phase, prior to the introduction of the dynamics. One should keep in mind that one aim of the histories programme is to describe physical systems that have non-trivial temporal structure –as arising, for instance, in quantum gravity– and are, perhaps, not amenable to a Hamiltonian description. The equation (18) for the decoherence functional is of sufficient generality to persist even in such contexts.
Second, following our earlier reasoning, it is easy to show that the fine-grained expressions for the decoherence functional can be used to determine its values for general coarse-grained histories. In analogy to (17) they read
$$d(\alpha ,\beta )=\underset{\psi ()\alpha }{}\underset{\varphi ()\beta }{}\varphi (t_i)|\rho _0|\psi (t_i)\psi (t_f)|\rho _f|\varphi (t_f)e^{iS[\psi (.)]iS^{}[\varphi (.)]}$$
Finally, the knowledge of the geometric phase—for a set of histories and of the automorphism that implements the dynamics—is sufficient to fully reconstruct the decoherence functional – and hence all the probabilistic content of a theory. The contribution of the initial state can be obtained by convex combinations of a pure state at the initial moment of time. What is interesting, is that at this level there is no need for our system to be described by a Hilbert space. All that is needed is a space of paths—on any manifold, the $`U(1)`$ connection from which the functional giving the Berry phase will be constructed and the dynamical law in the form of an automorphism of the space of observables. This can be an important starting point for developing geometric procedures for quantisation of quantum mechanical histories.
## 6 Conclusions
From Eq. (13) we notice that the Berry phase arises solely from the ordering in time of the projection operators, as they appear in the decoherence functional. It eventually corresponds to the kinematical part of the action Eq. (20). The Hamiltonian part appears due to Heisenberg-type time evolution of the projectors. This distinction is a fundamental and impressive feature of history theories that was identified in . There exist two distinct ways, in which time appears in physical theories: as a distinction between past and future (partial ordering property of time) and as the parameter underlying the evolution laws (time as parameter of change).
One of us (N.S.) has shown that these notions of time are associated to the kinematical and dynamical part of the action functional respectively, and there exist distinct operators that generate time translation with respect to these two parameters. They are an irreducible part of any theory that is based on temporally extended objects, whether classical or quantum. This distinction is manifested in the two different ways the time parameter appears in the definition of the class operator $`C_\alpha `$. From Eqs. (1320) we see that this is identical to the distinction between the geometric and the dynamical phases of standard quantum theory. In a sense, this is the only non-trivial remnant in canonical quantum theory of the temporal structure of history theories. The reader is referred to for a fuller treatment of this issue and to for the merits of the quantisation scheme motivated by it.
The fact that the off-diagonal elements of the decoherence functional correspond to the difference in Berry phase between its histories, suggests that the current interpretation of probabilities in the consistent histories scheme is at least incomplete. The relative geometric phase between two histories is a measurable quantity, while the present interpretation gives physical meaning to the values of only the diagonal elements of the decoherence functional.
Of course, one might argue that the geometric phase is measured only by comparing statistical measurements in two different ensembles of systems. As such, it may be described as any other measurement in the scheme. However, the point we make is that, the off-diagonal elements of the decoherence functional have a clear geometric and operational meaning. Therefore, an interpretational scheme that ignores them might face a truncation of the physics it addresses. In addition, the Berry phase would constitute a quantity that cannot be explained in terms of the properties of an individual quantum system, even though it is measured in ensembles. This is extremely problematic for the aims of a realist interpretation of quantum theory.
Our results highlight the presence of the complex phases in time evolution at the purely kinematical level, as the main contributors in the non-additivity of the probability measure for histories.
This strongly suggests that the presence of complex numbers in quantum theory is intrinsically linked to its distinct “probabilistic” structure. To see this consider the following.
First, both classical and quantum probability theory at a single moment of time are described by an additive measure over a lattice of propositions. But when time-evolution takes place in quantum theory, there appear complex phases that render the probability measure non-additive (this is the essence of the interference of histories).
Second, the pure time evolution in standard quantum theory is of a Hamiltonian type on $`PH`$; the dynamical phases that are generated by the Hamiltonian, are structurally not different from any angle variables of classical mechanics. There is nothing inherently complex in them, as the Schrödinger equation can be written without any reference to an $`i`$. On the other hand the geometric phase appears due to the bundle structure of the quantum mechanical space of rays. The bundle structure arises in the first place, because single-time probabilities do not depend on phase. Hence, even in standard quantum theory there is an indirect relation between the Berry phase and the probability assignment. This is brought fully into focus in the histories formalism.
We explained in the introduction, that when we are restricted in a single moment of time, the structures of quantum theory are in one-to-one analogy with the ones of classical probability theory. From an operational perspective, quantum mechanics at a single moment of time may be formulated without making any reference to complex numbers; it can be stated solely in terms of real-valued observables, expectation values and probabilities. It is only, when we study physical systems in a temporal sense that complex numbers appear. However, their appearance cannot be attributed to the law of time evolution.
Dynamics, in general, appears as an automorphism of the space of observables: if the observables are defined as real-valued, they will remain real-valued when dynamically transformed. For example, Schrödinger’s equation does not need introduce the complex unit; it can be equally well written in a real Hilbert space .
Alternatively, one may substitute Schrödinger’s equation with a real, partial differential equation on phase space—using either the Wigner or the coherent state transforms. Hence, while complex numbers in quantum theory are unavoidable when we study properties of the system at different moment of time, they are not introduced by the dynamical law. Furthermore, it is the temporal ordering that introduces phases, in an irreducible way, into the decoherence functional, that it is encoded in the definition of the class operator $`C_\alpha `$.
In other words, the geometric phase is a genuinely complex-valued object; and it is only the fact that we measure such an object, that forces us to accept complex numbers as an irreducible part of quantum theory. Complex number are not a necessary consequence of any dynamical law.
Hence, we conclude that, the complex structure of quantum theory is intrinsically linked to both its probability structure and the way the notion of succession is encoded. After all quantum theory is a theory of amplitudes, and the results from the above analysis imply that all physically relevant amplitudes—contained in the decoherence functional—are constructed from the geometric phase. As such they are geometrical in origin.
This is an intriguing result. It is a structural characteristic of quantum probability that should persist in frameworks that attempt to generalise quantum theory in a way that the Hilbert space is not a necessary ingredient.
## Acknowledgments:
C.A. was supported from the NSF grant PHY90-00967 and N.S. by a gift from the Jesse Phillips Foundation.
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# Nonequilibrium theory of scanning tunneling spectroscopy via adsorbate resonances: nonmagnetic and Kondo impurities
## I Introduction
A considerable body of experience and wisdom within the area of solid state tunneling phenomenon was built up throughout tunneling’s “Golden Era of the Sixties”. It was during this period that many of the defining fundamental ideas, basic theoretical strategies and methodologies, and broad scope of new applications for tunnel structures were first realized. A general introduction to many of these achievements can be found in a number of comprehensive volumes and in the Nobel Lectures of Esaki, Giaver, and Josephson, who were awarded the 1973 Nobel Prize in Physics for “their \[independent\] discoveries regarding tunneling phenomena in solids” . It is against this background that the astounding achievements in contemporary tunneling studies utilizing the single atom spatial resolution of the scanning tunneling microscope (STM) are most meaningfully considered . One phenomenon of key interest here which was first considered in the “Golden Era” is that of impurity/adsorbate-assisted elastic tunneling. Two bodies of work are particularly relevant to the present study. The first is the recognition by Appelbaum and coworkers of the possible role of the Kondo effect in determining certain current-voltage characteristics (e.g. “zero-bias anomalies”) of metal-oxide-metal tunnel junctions containing localized paramagnetic impurity states near the metal-oxide interfaces . Second are the resonance tunneling studies involving valence electronic levels of single atoms adsorbed on metal surfaces, as probed in a field emission microscope configured for energy analysis (thus enabling electron spectroscopy) of the field emitted electrons . Many years later, useful parallels between the theory of single atom resonance tunneling developed in the “Golden Era” and the theory of the STM, in the single-atom-tip limit, were unambiguously established . Further discussion of these issues from the past will be offered throughout the text, when appropriate.
The basis for continuing interest and excitement in impurity/adsorbate- assisted tunneling is that the transparency of tunnel junctions can be dramatically enhanced by the presence of states localized within the barrier when they are in resonance with the tunneling electrons. Tunneling through such states is, for example, the origin of conductance fluctuations quantum dots exhibit in the Coulomb blockade . The tunneling probability in the presence of a “barrier” state is proportional to the spectral density produced by the hybridization of the localized state with the conduction electrons and in many situations the current is given by the Breit-Wigner formula
$$I\frac{\mathrm{\Gamma }^2}{(\omega ϵ_0)^2+(\mathrm{\Gamma }/2)^2}.$$
(1.1)
Here $`\mathrm{\Gamma }`$ is the width at half maximum of the resonance produced by the hybridization with the conduction electrons in the right and left lead and $`ϵ_0`$ is the energy of the local state. The value of $`\mathrm{\Gamma }`$ depends on the height and width of the barrier potential between the central region and the leads.
Recently, enhancements in the zero bias conductance in quantum dots due to the Kondo effect have been observed . It had been shown earlier that $`(dI/dV)`$, the zero-bias differential conductance, is proportional to the Fermi level density of states of the Kondo resonance on the quantum dot. Similarly, the Kondo resonance has been spectroscopically observed on single magnetic impurities adsorbed on metal surfaces using the STM . However in the case of the spectroscopic STM experiments, the resonance at the Fermi level appears to have an asymmetric shape and cannot be interpreted simply in terms of the local density of states of the impurity atom. Rather, the electron tunneling current – being a coherent quantum effect – is a result of interference between competing tunneling channels, as will soon be detailed. Unlike in the quantum dot where the tunneling can take place with appreciable magnitude only through the quantum dot region, the apparent tunneling current from the STM tip to the surface can either go through the resonance localized on the impurity or directly into the conduction states of the surface. The distinction between the conduction and local states will be discussed later. The notion that the tunneling conductance is proportional to the local density of states near the STM tip must then be modified.
In addition to its most common use for observing/determining atomic geometrical structure at surfaces, the STM is used as a sensitive probe of surface electronic structure. Various theoretical approaches to the STM conductance employ the tunneling Hamiltonian introduced by Bardeen and Golden Rule type expressions in which under certain limiting conditions of practical interest the STM conductance is indeed determined by the surface density of electronic states near the STM tip . Tersoff and Hamann developed a widely used model of the scanning tunneling microscope that includes the three dimensionality and spatial resolution of the tip.
The generic problem of a discrete state interacting with a continuum of states arises in many different areas of physics and chemistry . In condensed matter physics a frequently occurring realization is the electronic state of an impurity atom immersed within a host lattice. In the case of magnetic impurities, the interaction gives rise to nontrivial phenomena such as the Kondo effect .
Within the context of atomic physics, Fano discussed related effects, as they might appear on observable absorption lineshapes or resonant electron scattering cross sections which are due to the configuration interaction (CI) that couples a discrete two-electron excited atomic state with a continuum of ionization states . Whilst the “natural” lineshape of the resonance is Lorentzian, when studied by experiments in which an external probe interacts with the system, the resonance can appear to have an asymmetric lineshape. Such lineshapes are referred to as Fano resonances. Fano found that an asymmetry in absorption lineshapes is due to interference between the excitation or decay into CI-mixed discrete and continuum states which both couple to the external probe. If the coupling between the probe and continuum is expressed in terms of an energy independent matrix element $`t_c`$ and the interaction between the probe and the localized state (which has already been diluted by admixture into the continuum) by the matrix element $`\stackrel{~}{t}_a`$, then the lineshape detected has the form
$$I\frac{(q+2(\omega ϵ_0)/\mathrm{\Gamma })^2}{1+(2(\omega ϵ_0)/\mathrm{\Gamma })^2},$$
(1.2)
where $`q\stackrel{~}{t}_a/(2\pi Vt_c)`$ with $`V`$ being the hybridization (or CI) matrix element between the local state and continuum. The latter coupling results in the discrete state acquiring a width $`\mathrm{\Gamma }=2\pi \rho _sV^2`$ ($`\rho _s`$ is the density of continuum states).
In the present paper, we consider the problem of the discrete state embedded in a continuum using a probe such as the STM that has atomic scale spatial resolution. This work has been motivated by the recent STM experiments involving single Kondo impurities . The resonance observed in the conductance was interpreted by the authors in terms of the Fano interference. The fit of the resonance to the Fano formula – generalized to the case where intra-atomic Coulomb interaction on the impurity is taken into account – was based on the assumption that the tunneling into an Anderson impurity can be extended to include the tunneling into the continuum in a straightforward way. Upon further consideration, it appears that this generalization of the Fano result to the case of STM conductance is not as straightforward as has been presumed. In the present paper we obtain a more complicated expression than the elementary Fano formula, one which accounts for the correct asymptotic behavior for large tip-impurity separation. In particular, when the dependence of the probe’s distance from the local state is properly included, then observable consequences of the local state admixture with the conduction electrons show the correct asymptotic long range behavior in the large tip-to-impurity-separation-limit as they obviously should.
The difference between the Fano lineshape and the result obtained here is due to the different nature of the probe in the “multi-center” STM configuration (one “center” on the impurity/discrete state, the other on the STM tip/probe) compared with the “single-site” atomic physics processes. While Fano was concerned with light absorption or electron scattering where the system under study was always at the probe’s focus, the outcome of STM experiments must depend upon the variable spatial position of the probe (tip) with respect to the discrete state under investigation. Put another way, for the atomic physics applications considered by Fano, both the discrete state coupled to the continuum and also the initially excited decaying or autoionizing state (the “probe” state in our STM language) are atomic states spatially localized at the same site by the same atomic central potential. In contrast, the “S”(=scanning!) in STM assures that the tip, hence initial excited state, can be independently located with respect to the position of the “discrete state coupled to the continuum” and it is this extra degree of freedom that enriches the potential information content in STM lineshape analysis, but also requires a much more detailed theoretical treatment than merely fitting the atomic physics Fano lineshape, Eq. (1.2), to position-dependent STM spectra. This will be expanded upon in depth later. This realization demonstrates the importance and crucial need for considering the measurement process in quantum mechanical observations.
While the structuring of this comprehensive paper is based on a logical development of the subject matter, it may be useful to present a roadmap of key points and results to guide the casual as well as the dedicated reader. In section II, we introduce the model of the system and discuss our approximations of the tunneling matrix elements resulting in Eqs. (2.15), (2.16), and (2.20). In section III and the supporting appendices (A-C), we develop the general nonequilibrium theory of STM tunneling current and conductance in the presence of an adsorbate induced resonance. The more familiar equilibrium limit, asymptotically exact for large tip-surface separation but pragmatically useful even for moderate $`5`$Å separations, is treated as a special case in III B. The relationship between the equilibrium tunneling resonance lineshape observable in STM experiments and the asymmetric Fano lineshape is established in III C (Eqs. 3.24 and 3.25). The nonequilibrium contribution to the lineshape treated as a correction to the equilibrium limit is taken up in III D. The differential conductance is introduced in III E. The crucial role of the substrate electronic structure is explicitly considered in section III F where the substrate Green’s function is evaluated for a jellium surface and in III G where surface states and real electronic structure effects are discussed qualitatively.
In section IV, we illustrate the predictions and consequences of our theory on two models for the adsorbate: nonmagnetic and Kondo. Using a jellium substrate model, we first discuss the common features of the two adsorbate models in terms of the non-interacting Anderson model in IV A. Four different families of spectroscopic lineshape variations are taken up: dependence on (1) the relative strength of the tip-to-adsorbate vs. tip-to-substrate tunneling and on the substrate electronic structure; (2) tip-surface separation; (3) lateral tip position; (4) temperature due to Fermi level smearing. We explain why no variations in the resonance lineshape should be observed for lateral tip positions on length scales $`(12)`$Å – characteristic of the bulk $`k_F`$ – and only small lineshape variations should be expected for vertical tip variations (in experimentally accessible range). We leave the discussion of the tip and bias effects on lineshapes to section IV B 3.
Tunneling characteristics specific to Kondo systems are presented in IV B. We begin with conceptual issues in IV B 1. We then discuss the recent experiments on $`Co/Au(111)`$ and $`Ce/Ag(111)`$ and relate our work to related theoretical papers IV B 2). In particular, we show that the temperature dependence of the Kondo resonance is not easily extractable from the temperature dependence of the differential conductance. We reiterate that the stability of the experimental lineshape with the tip position is to be expected. Variations in the lineshapes would, however, occur at larger distance due to tunneling into the surface state. The effect of adsorbate-tip hybridization and bias induced nonequilibrium on the current (conductance) vs. bias measurements in Kondo systems are dealt with in section IV B 3. Finally, an enumeration of specific conclusions is offered in section V.
## II Model and approximations
Models of scanning tunneling microscopy are abundant in the literature of the last two decades and standard texts exist . We approach the problem as a nonequilibrium process and discuss the corrections to the Tersoff-Hamann formulation . Our intent is to develop such theory under general and self-consistent assumptions that accurately capture most of the qualitative aspects involved and do so in a way that make extension to more realistic calculations formally straightforward. We focus on the tunneling through adsorbate resonances. Throughout this paper, we adopt the convention that the energies are measured with respect to the respective Fermi levels of the substrate and tip unless specified otherwise and set $`\mathrm{}=1`$. When the tip is biased we explicitly shift the tip energies.
### A Model of the studied system
We consider a system which consists of a clean metallic surface with a single impurity atom adsorbed on it. The STM will be used to study the system by means of tunneling through a resonance produced by an electronic state of the impurity, such as the $`5f`$ orbital of $`Ce/Ag(111)`$ or $`3d`$ orbital in $`Co/Au(111)`$ . Unless otherwise noted, we place the origin of the coordinate system at a point on the surface of the metal directly below the adsorbate. This means that the position of the impurity is $`\stackrel{}{R}_0=(0,0,Z_0)`$. The system without the probe is described by the degenerate Anderson Hamiltonian
$`H_s(\stackrel{}{R}_0)`$ $`=`$ $`{\displaystyle \underset{a}{}}ϵ_0(\stackrel{}{R}_0)c_a^{}c_a+{\displaystyle \underset{a>a^{}}{}}U(\stackrel{}{R}_0)n_an_a^{}+`$ (2.1)
$`+`$ $`{\displaystyle \underset{ka}{}}ϵ_kc_{ka}^{}c_{ka}+{\displaystyle \underset{ka}{}}\{V_{ka}(\stackrel{}{R}_0)c_{ka}^{}c_a+\mathrm{H}.c.\}.`$ (2.2)
Here, $`ϵ_0`$ is the energy of the impurity state $`\psi _\sigma (\stackrel{}{r})`$, which we assume may be a multiplet of states described collectively by the quantum number $`a(m\sigma )`$. In the simplest case, $`a`$ correspond to the spin $`\sigma `$, but it may also include orbital degeneracy ($`m`$) in more complicated cases. In this paper, we discuss at most spin degenerate states with $`a=\sigma `$ and $`N=2`$ (degeneracy). We denote by $`c_a^{}`$ the creation operator for this state. The $`ϵ_k`$ is the conduction band state energy – independent of $`\sigma `$ in the absence of magnetic field – with $`c_{ka}^{}`$ being the creation operator for the corresponding Bloch state with symmetry (spin) $`a`$ common with the impurity state, and $`V_{ka}`$ is the matrix element for hybridization between the impurity and conduction states. The second term in (2.1) corresponds to the intra-atomic Coulomb interaction between electrons in the impurity state $`\psi _a`$.
If the renormalized energy $`ϵ_0`$ lies within the conduction band, the bound state broadens into a resonance which in the wide-band limit and with $`U=0`$, has a Lorentzian shape
$$\rho _0(\omega )=\frac{1}{2\pi }\frac{\mathrm{\Gamma }}{(\omega ϵ_0)^2+(\mathrm{\Gamma }/2)^2},$$
(2.3)
where
$$\mathrm{\Gamma }=2\pi \rho _sV^2,$$
(2.4)
with $`\rho _s`$ and $`V`$ the assumed-energy-independent density of conduction states and hybridization matrix element from (2.1). In most of our later considerations however, we will retain an energy dependence in the model density of states, $`\rho _s(\omega )`$.
### B Interaction of the system with the STM tip
When the STM tip is brought near the impurity, electrons can tunnel between the tip and the adsorbate state. This situation is expressed by adding an interaction term to the Hamiltonian
$$H_{at}(\stackrel{}{R}_t,\stackrel{}{R}_0)=\underset{pa}{}\{t_{ap}(\stackrel{}{R}_t,\stackrel{}{R}_0)c_a^{}c_p+H.c.\}.$$
(2.5)
The tip states are denoted by subscript $`p`$ while the conduction states of the metal by subscript $`k`$. The transfer matrix $`t_{ap}`$ depends on the position of both the adsorbate $`Z_0`$ and the tip $`\stackrel{}{R}_t(\stackrel{}{R}_{},Z_t)`$.
If the STM only coupled to the discrete state with transfer amplitude $`t_a`$ then the conductance would, in the wide-band limit ($`t_{ap}t_a`$, independent of $`p`$), be determined by
$$G|t_a|^2\frac{\mathrm{\Gamma }}{(\omega ϵ_0)^2+(\mathrm{\Gamma }/2)^2}$$
(2.6)
and the conductance would thus be directly related to the impurity density of states. This is reminiscent of the defining characteristics from field emission resonance tunneling spectroscopy in which tunneling from the substrate to vacuum is disproportionately smaller than that from “good” adsorbates . However, since in the STM geometry, tunneling directly between the tip and the metal surface can be comparable to (or in excess of) that between the tip and the impurity, the conductance exhibits a more complex behavior than that of a simple impurity local density of states. We take such processes into account through
$$H_{st}(\stackrel{}{R}_t)=\underset{pk}{}\{t_{kp}(\stackrel{}{R}_t)c_k^{}c_p+H.c.\},$$
(2.7)
where the tunneling matrix element $`t_{kp}`$ depends on the position of the tip $`\stackrel{}{R}_t`$. In the rest of this section, we will address the issue of the probe’s effect on the system itself. This is particularly important when the tip is brought very close to the adsorbate so that the tunneling $`t_{ap}`$ and $`t_{kp}`$ are comparable with $`V_{ka}`$. In this case, the width $`\mathrm{\Gamma }`$ of the adsorbate resonance is no longer given by (2.4) but rather by
$$\mathrm{\Gamma }=2\pi (\rho _sV^2+\rho _tt_a^2)$$
(2.8)
The effect of the tip on the conduction states can also be important. However, the most important consequence of the tip perturbation comes when finite bias is applied across the tunnel junction. The system will be out of equilibrium and the problem of the system-probe must be approached self-consistently with the tip included in the system it is probing. The typical operational mode of the STM during imaging and spectroscopic measurements is such that the tip distance from the adsorbate is several atomic units larger than the adsorbate-surface separation and therefore $`|t_{ap}|,|t_{kp}||V_{ka}|`$. In this limiting case, the probe has no effect on the adsorbate-metal complex other than as a source of hot tunneling electrons. This is a reasonable assumption as long as the tip separation is large enough to justify the approximation $`|t_{ap}|,|t_{kp}||V_{ka}|`$.
### C Model of the STM tip
An important property of the tip is its spatial resolution, as discussed in great detail by many . We will consider the tip to be well defined and terminated by a single atom through which the tunneling predominantly takes place. This is the s-wave model of Tersoff and Hamann , see Fig 1. The important features are the following: (a) the tip Hamiltonian is
$$H_t=\underset{p}{}ϵ_pc_p^{}c_p,$$
(2.9)
where $`c_p^{}`$ creates an electron in the state $`\psi _p(\stackrel{}{r})`$ with energy $`ϵ_p`$ measured from the Fermi level of the tip $`ϵ_{Ft}`$; (b) when the tip is positioned near the surface, tunneling into and out of a state $`\psi _p`$ can take place; (c) the states are filled up to the chemical potential $`ϵ_{Ft}`$ controlled by the bias; (d) the tip states are characterized by a density of states which we denote by $`\rho _t(\omega )=_p\delta (\omega ϵ_p)`$; and (e) the asymptotic form of the tip eigenstate $`\psi _p`$ in the vacuum region extending towards the metal surface is characterized by the atomic orbital of the apex atom. The wavefunction $`\psi _p`$ can be found based on simple physical arguments without solving the complete problem. If $`\varphi _t`$ is the work function of the tip and $`\kappa _t\sqrt{2m_t^{}(\varphi _tϵ_p)}`$, then following Tersoff and Hamann,
$$\psi _p(\stackrel{}{r})Re^{\kappa _tR}\frac{\mathrm{exp}(\kappa _tr)}{r}$$
(2.10)
where R is the radius of curvature of the tip about its center which is located at the origin of this “tip-defining” coordinate system. While (2.10) represents an “s-wave tip”, more generally $`\psi _p`$ would carry whatever symmetry was possessed by the relevant atomic orbital centered at the tip apex . The wavefunction tail, controlled by $`\kappa `$, depends on the bias and tip-surface separation. Both factors modify the height of the vacuum barrier, hence effective work function $`\varphi `$ determining $`\kappa `$. These modification can be essentially included by renormalizing the wavefunction tails and the densities of states by position and energy dependent factors via the tunneling matrix elements. We discuss the tunneling matrix elements next.
### D Approximations for the tunneling matrix elements
An important role in our formulation is played by the tunneling (hybridization) matrix elements $`V_{ka}`$, $`t_{ap}`$, and $`t_{kp}`$ since they include the dependence on electronic structure and the tip and adsorbate position. The desired quantitative accuracy of the model for the tunneling process is to a large degree determined by the approximations made in the evaluation of these matrix elements. We do this first for a general adsorbate-metal system and then for a jellium model. It is rather straightforward to include band structure effects using a realistic electronic structure calculation of the substrate Green’s function.
We begin with the discussion of the matrix elements $`V_{ka}`$ and $`t_{kp}`$ that contain the metal wavefunctions $`\psi _\stackrel{}{k}(\stackrel{}{r})`$. They have the form
$$M_{kl}(\stackrel{}{R}_l)=d^3r\psi _\stackrel{}{k}^{}(\stackrel{}{r})v_{sl}(\stackrel{}{r};\stackrel{}{R})\psi _l^{}(\stackrel{}{r}),$$
(2.11)
where $`v_{sl}`$ is the potential representing the mutual interaction of the two systems. The wavefunction $`\psi _\stackrel{}{k}(\stackrel{}{r})`$ is a Bloch state of the unperturbed metal and $`\psi _l^{}(\stackrel{}{r})`$ is the wavefunction (in the coordinate system of the metal) of the adsorbate state $`a`$ in the case of $`V_{ka}`$ or the tip wavefunction in the case of $`t_{kp}`$. Either way, $`\psi _l^{}(\stackrel{}{r})`$ can be written in terms of the wavefunction with the origin at the adsorbate (tip) as $`\psi _l^{}(\stackrel{}{r})=\psi _l(\stackrel{}{r}\stackrel{}{R}_l)`$, where $`\stackrel{}{R}_l`$ is the position of the adsorbate ($`\stackrel{}{R}_0`$) or the tip apex atom ($`\stackrel{}{R}_t`$) measured from a reference point on the surface. The energies $`ϵ_k`$ of the metal electrons are written in terms of the perpendicular and parallel components as $`ϵ_kϵ_{k_z}+ϵ_k_{}`$. We follow the convention that $`ϵ_k_{}`$ is measured from the bottom of the 2-D band and $`ϵ_{k_z}`$ is measured with respect to $`ϵ_{Fs}`$. For example, we write for the jellium model $`ϵ_{kz}=k_z^2/2m_s^{}D`$ and $`ϵ_k_{}=k_{}^2/2m_s^{}`$, where $`(D)`$ is the energy of the bottom of the band with respect to the Fermi level. The Bloch states can generally be written in the relevant region outside the metal as
$$\psi _{n\stackrel{}{k}}(\stackrel{}{r})=e^{\kappa _{ns}z}u_{n\stackrel{}{k}_{}}(\stackrel{}{\rho },z)e^{i\stackrel{}{k}_{}\stackrel{}{\rho }}.$$
(2.12)
where $`n`$ is the band index, $`u_{n\stackrel{}{k}_{}}(\stackrel{}{\rho },z)`$ is a function weakly dependent on $`z`$ outside the surface and periodic in $`\stackrel{}{\rho }`$, the electron coordinate in the plane of the surface. This form is equally valid for the metal band gap surface states that seem ubiquitous to STM studies on (111) noble metal surfaces since both the z-propagating Bloch states and the localized surface states are eigenstates of the same Hamiltonian with different eigenvalues at a given $`k`$. At the surface, evanescent states into the bulk that appear upon analytic continuation of the band structure into the domain of complex k-vectors cannot be rejected on the basis of physical considerations as they were for the perfectly periodic interior of the solid.
For jellium, $`u_{n\stackrel{}{k}_{}}(\stackrel{}{r})`$ is constant and the the metal states are then simply plane waves along the surface with exponentially decaying amplitude into the vacuum. Here $`\kappa _{ns}=\sqrt{2m_s^{}(\varphi _sϵ_{nk_z})}=\sqrt{2m_s^{}(\varphi _sϵ_{nk}+ϵ_{nk_{}})}`$ with $`m^{}`$ the metal electron effective mass (number) and $`\varphi _s`$ the height of the tunneling barrier for a Fermi level electron. For bias voltages much smaller than the work function, $`\varphi _s`$ is equal to the metal work function. We omit the band index $`n`$ in the rest of the paper unless we explicitly discuss the electronic structure effects. We also define $`\lambda _\omega ^1=\sqrt{2m_s^{}(\varphi _s\omega )}`$, where the energy factor $`(\varphi _s\omega )`$ represents the effective tunneling potential barrier for an electron with energy $`\omega `$. For small bias voltages, $`\omega 0`$, and we can replace $`\lambda _\omega `$, which depends weakly on energy in this range, by its Fermi level value $`(\lambda )`$.
We shift the integration variable in the integral (2.11) to an origin centered on the adatom (tip), which gives
$`M_{kl}(\stackrel{}{R}_l)`$ $`=`$ $`\psi _k^{}(\stackrel{}{R}_l){\displaystyle _{z>Z_l}}d^3re^{\kappa _sz}e^{i\stackrel{}{k}_{}\stackrel{}{\rho }}\times `$ (2.13)
$`\times `$ $`{\displaystyle \frac{u_\stackrel{}{k}^{}(\stackrel{}{r}+\stackrel{}{R}_l)}{u_\stackrel{}{k}^{}(\stackrel{}{R}_l)}}v_{sa}^{}(\stackrel{}{r},R_l)\psi _l(\stackrel{}{r}).`$ (2.14)
Matrix elements of this type have been the focus of intense study in the context of charge transfer processes at surfaces. The main contribution to the integral in (2.13) comes from the region just outside the surface where the $`z`$-dependence of the metal states is that of a decaying exponential. In the case of both the adatom and the tip, the integrand is assumed to be well localized to the $`z>Z_0`$ region by virtue of the spatial properties of $`\psi _l`$. Consequently, the integral will be reasonably constant for all $`k_{}1/r_l`$ where $`r_l`$, the radial length scale for the atom or tip, sets the range of $`k_{}`$. Since, as will soon be demonstrated, other factors in the full problem will provide a much more severe $`k_{}`$ cutoff, it is sufficient to represent the integral by a constant. We then write the matrix element approximately $`M_{kl}(\stackrel{}{R}_l)M_0\psi _\stackrel{}{k}^{}(\stackrel{}{R}_l)`$, where $`M_0`$ is the overlap integral defined in (2.13) that contains the dependence of the matrix element on the symmetry of the atomic orbitals near the tip and on selection rules.
Since the distance of the adsorbate from the surface is small, it will be sufficient for the purpose of $`V_{ka}`$ to write the Bloch state in the form $`\psi _\stackrel{}{k}(Z_0)=e^{Z_0/\lambda }\psi _\stackrel{}{k}(0)`$ with the decay constant independent of the k-vector. With the choice of our coordinate system $`\stackrel{}{R}_0=(0,Z_0)`$, we write $`V_{ka}`$ in the form of a separable product of a k-independent function of adsorbate position multiplied by a function of $`\stackrel{}{k}`$
$$V_{ka}(\stackrel{}{R}_0)V_a(Z_0)\psi _\stackrel{}{k}^{}(0)$$
(2.15)
where $`V_a(Z_0)=V_0e^{Z_0/\lambda }`$. We cannot make this last approximation in the tip-to-surface tunneling matrix element $`t_{kp}`$, since the interplay between $`k`$ and $`Z_t`$ dependence has an important role in the tunneling process. However, from the conceptual point of view, we find it convenient to isolate a $`k`$-independent dependence on the tip position in $`t_c(Z_t)=t_0e^{Z_t/\lambda }`$ and write the matrix element in the form
$$t_{kp}(\stackrel{}{R}_t)=t_c(Z_t)e^{Z_t/\lambda }\psi _\stackrel{}{k}^{}(\stackrel{}{R}_t).$$
(2.16)
We note that the s-wave tip-to-surface tunneling matrix element $`t_{kp}`$ has been shown by Tersoff and Hamann under quite general assumptions to assume the form (2.16) with the tip wavefunction given by Eq. (2.10), which indicates that our simple qualitative arguments seem to be supported by more detailed analysis.
The tip-to-impurity matrix element
$`t_{ap}(\stackrel{}{R}_t,\stackrel{}{R}_0)=`$ (2.17)
$`={\displaystyle d^3r\psi _\stackrel{}{p}^{}(\stackrel{}{r}\stackrel{}{R}_t)v_{at}(\stackrel{}{r};\stackrel{}{R}_t,\stackrel{}{R}_0)\psi _a(\stackrel{}{r}\stackrel{}{R}_0)}.`$ (2.18)
depends on the position of both the adsorbate ($`Z_0`$) and the tip ($`\stackrel{}{R}_t`$). Since the tunneling from the tip takes place predominantly through the apex atom, the wavefunction $`\psi _p`$ in the last expression, a generalization of (2.10), can be written (with the tip at origin) as
$$\psi _\stackrel{}{p}(\stackrel{}{r})=e^{\kappa _tR}\frac{e^{\kappa _tr}}{r}\psi _s(\widehat{r})$$
(2.19)
where $`\psi _s`$ is the angular part of the orbital localized on the apex atom whose center is located at $`\stackrel{}{R}_t`$. The tunneling matrix element $`t_{ap}`$ will thus reflect the symmetry and spatial dependence of the states localized on the adsorbate and in the tip apex. Within the s-wave tip model, $`\psi _s`$ is just an innocuous constant. Since $`\psi _p`$ and $`\psi _a`$ appearing in Eq. (2.17) are both atomic-like functions in the relevant region of overlap centered respectively on the tip and on the adatom, in a broad sense $`t_{ap}`$ is similar to a common two-center hybridization/hopping integral defining the binding in a diatomic molecule . Typically the magnitudes of these integrals are exponentially decreasing functions of their separation (possibly multiplied by a mildly oscillatory function accounting for the nodal structure of the atomic functions). Based on this analogy from quantum chemistry, $`t_{ap}`$ given by Eq. (2.17) should take the form
$$t_{ap}(\stackrel{}{R}_t,Z_0)t_ae^{|\stackrel{}{R}_tZ_0\widehat{i}_z|/\alpha }t_a(\stackrel{}{R}_t,Z_0)$$
(2.20)
Here, $`\alpha ^1\left(\kappa _t+\sqrt{2m(\varphi _sϵ_0)}\right)`$ is an effective decay constant evaluated for states at the Fermi level of the tip, $`|\stackrel{}{R}_tZ_0\widehat{i}_z|=\sqrt{R_{}^2+(Z_tZ_0)^2}`$ is the tip-to-atom separation, as depicted in Fig. 1, where $`\stackrel{}{R}_{}`$ is the parallel component of $`\stackrel{}{R}_t`$. The decay constant $`\kappa _t`$ depends on the energy of the tip state $`(ϵ_p)`$, but this dependence is very weak for small biases considered here and we neglect it. In this case, the matrix element is well approximated by the $`p`$ independent form $`t_a(\stackrel{}{R}_t,Z_0)`$. The matrix element $`t_a`$ may be taken real.
## III The nonequilibrium theory of the tunneling current and differential conductance
The tunneling between the adsorbate-metal complex and a biased tip is a nonequilibrium process. Although we frequently make the assumption in this paper that the tip-system interaction is weak enough so that local equilibrium is maintained to a good approximation, the assumption is less valid when the tip is near the surface. For this reason, we develop our theory of the tunneling process within the Keldysh-Kadanoff framework for the nonequilibrium Green’s functions and discuss the nonequilibrium corrections.
### A General expression for the tunneling current in terms of nonequilibrium Green’s functions
We define the tunneling current as the flow of electrons through a closed surface around the tip. It is expressed in terms of the continuity equation as
$$I=e\frac{dn_t(t)}{dt}$$
(3.1)
where $`n_t=_pc_p^{}c_p`$ is the number operator for the tip electrons, and the brackets signify the ensemble average, which in the local equilibrium case is the thermal average over the tip states. The time derivative is found from the Schrödinger equation of the total Hamiltonian of the tip-substrate-adsorbate system $`H_{tot}=H_s+H_t+H_{at}+H_{st}`$. Since the number operator commutes with $`H_s`$ and $`H_t`$, the only contribution comes from the interaction terms $`H_{at}`$ and $`H_{st}`$ and the current is
$$I=\frac{2e}{\mathrm{}}\mathrm{I}m\{\underset{kp}{}t_{kp}c_k^{}(t)c_p(t)+\underset{ap}{}t_{ap}c_a^{}(t)c_p(t)\},$$
(3.2)
where we omitted the arguments in $`t_{kp}`$ and $`t_{ap}`$ for convenience and used the relation $`t_{ap}=t_{pa}^{}`$ and $`t_{kp}=t_{pk}^{}`$. We write the arguments explicitly only when we wish to emphasize their dependence. We define the time loop Green’s functions $`G_{pa}(t,t^{})=iT_Cc_p(t)c_a^{}(t^{})`$ and $`G_{pk}(t,t^{})=iT_Cc_p(t)c_k^{}(t^{})`$, where $`T_C`$ orders the times along a contour $`C`$ in the complex time plane. The contour can be taken to be the Kadanoff-Baym contour , the Keldysh contour , or a more general choice. The discussion of nonequilibrium Green’s function is available in standard books and review articles and we refer the reader to these references for further details. Our notation follows closely that of the more detailed discussion in reference . In the present paper, we are mainly interested in steady state tunneling current (time independent) and therefore we work with the Fourier transformed quantities in frequency rather than time space. The current can be written as
$$I=\frac{2e}{h}\mathrm{I}m_{\mathrm{}}^{\mathrm{}}𝑑\omega \{\underset{kp}{}t_{kp}G_{pk}^<(\omega )+\underset{ap}{}t_{ap}G_{pa}^<(\omega )\},$$
(3.3)
where $`G_{pk}^<(\omega )`$ and $`G_{pa}^<(\omega )`$ are the Fourier transforms of $`G_{pk}^<(t,t^{})=c_k^{}(t^{})c_p(t)`$ and $`G_{pa}^<(t,t^{})=c_a^{}(t^{})c_p(t)`$ – the analytic pieces on the real time axis of the Green’s functions introduced above. Equivalently, the current may be calculated from $`\frac{dn_s}{dt}+\frac{dn_a}{dt}`$. It is easy to see that this approach also leads to the equation (3.2).
The problem of finding the current thus reduces to finding the “lesser” Green’s functions $`G_{pa}^<`$ and $`G_{pk}^<`$. This is done using the equation of motion method for the time ordered Green’s functions in Appendix B and the rules for analytic continuation described in Appendix C. In order to see the interference between the two scattering channels giving rise to the Fano lineshape, it is also useful to express the current using (B2) and (B6) in Eq. (3.3) as
$`I`$ $`={\displaystyle \frac{2e}{h}}\mathrm{I}m{\displaystyle \underset{p}{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega `$ (3.6)
$`\{G_p^{(0)}[{\displaystyle \underset{a}{}}|t_{pa}|^2G_a+{\displaystyle \underset{ka}{}}t_{ap}t_{pk}G_{ka}+`$
$`+{\displaystyle \underset{ka}{}}t_{kp}t_{pa}G_{ak}+{\displaystyle \underset{kk^{}}{}}t_{pk^{}}t_{kp}G_{k^{}k}]\}^<.`$
The first term corresponds to the tunneling into the adsorbate state hybridized with the metal electrons. The fourth term is the contribution to the current from the direct tunneling into the conduction states perturbed by the presence of the discrete adsorbate state. The second and third terms give the interference between the two channels.
We substitute the solutions (B4) and (B9) for $`G_{pk}`$ and $`G_{pa}`$ from Appendix B into Eq. (3.3) and write the tunneling current in the form
$`I=`$ $`{\displaystyle \frac{2e}{h}}\mathrm{I}m{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega \times `$ (3.7)
$`\times `$ $`{\displaystyle \underset{pp^{}}{}}\left\{G_{pp^{}}^R(\omega )T_{p^{}p}^<(\omega )+G_{pp^{}}^<(\omega )T_{p^{}p}^A(\omega )\right\},`$ (3.8)
where the retarded function $`T_{pp^{}}^R(\stackrel{}{R}_t,\stackrel{}{R}_0,\omega )`$ plays the role of the T-matrix for scattering of the tip electrons from the adsorbate-metal complex. It is defined by
$`T_{pp^{}}`$ $`=`$ $`{\displaystyle \underset{k}{}}t_{pk}G_k^0t_{kp^{}}+`$ (3.9)
$`+`$ $`{\displaystyle \underset{a}{}}\stackrel{~}{t}_{pa}G_a(t_{ap^{}}+{\displaystyle \underset{k}{}}\stackrel{~}{V}_{ak}G_k^0t_{kp^{}}),`$ (3.10)
with
$$\stackrel{~}{V}_{ka}=V_{ka}+\underset{p}{}t_{kp}G_p^0t_{pa}$$
(3.11)
and
$$\stackrel{~}{t}_{pa}=t_{pa}+\underset{k}{}t_{pk}G_k^0V_{ka}$$
(3.12)
the hybridization and tunneling matrix elements for the adsorbate modified by the tip-substrate interaction. The matrix $`T_{pp^{}}`$ incorporates the properties of the tip as well as the adsorbate into the expression for current. We discuss its physical meaning more in the next section.
### B Equilibrium limit of the tunneling current at large tip-surface separation
We define the equilibrium tunneling current as the large tip-surface separation limit of (3.7) when the tip, adsorbate, and substrate are all in local equilibrium. This is equivalent to keeping only the lowest order terms in $`t_{kp}`$ and $`t_{pa}`$. In our formalism, this is achieved by replacing $`\stackrel{~}{G}_{pp^{}}G_p^0\delta _{pp^{}}`$ in Eq. (3.7), $`\stackrel{~}{V}_{ka}V_{ka}`$ in $`T_{pp^{}}`$, and by using the fluctuation-dissipation relation $`G_i^<(\omega )=f_i(\omega )\rho _i(\omega )`$. The subscript $`i`$ stands for tip (t), adsorbate (a), and metal (s), respectively. The adsorbate is in equilibrium with the metal, i.e. $`f_a(\omega )=f_s(\omega )`$. The matrix $`T_{pp^{}}`$ is expressed entirely in terms of the Green’s functions of the system and the tunneling matrix elements $`t_{pk}`$, $`t_{pa}`$. Since these matrix elements reflect the symmetry of the apex atom wavefunction, but are only weakly dependent on $`p`$ on the energy scale of the resonance width, the matrix $`T_{pp^{}}`$ will also have this property. We therefore make an additional assumption that $`_pG_p^0T_{pp}(_pG_p^0)T_t`$, where
$`T_t={\displaystyle \underset{k}{}}t_{pk}G_k^0t_{kp}+{\displaystyle \underset{a}{}}\stackrel{~}{t}_{pa}G_a\stackrel{~}{t}_{ap}`$ (3.13)
is only a function of the atomic tip orbital independent of $`p`$. We define a tip-specific quantity observable by the STM, which is related to the local density of states
$$\stackrel{~}{\rho }_{sat}(\stackrel{}{R}_t,\stackrel{}{R}_0;\omega )=\frac{1}{\pi }\mathrm{I}mT_t^R(\stackrel{}{R}_t,\stackrel{}{R}_0;\omega ),$$
(3.14)
and write the equilibrium current $`I_{eq}(\stackrel{}{R}_t,\stackrel{}{R}_0,V)`$ as
$$I_{eq}=\frac{2e}{h}_{\mathrm{}}^{\mathrm{}}𝑑\omega \left[f_t(\omega ^{})f_s(\omega )\right]\rho _t(\omega ^{})\stackrel{~}{\rho }_{sat}(\omega ),$$
(3.15)
where $`\rho _t=_p\delta (\omega ϵ_p)`$ is the density of tip states and $`\omega ^{}=\omega eV`$ with $`V`$ being the bias voltage. This equations is easily related to traditional formulations given in terms of an integral product of an electron “supply function” multiplied by a tunneling or transmission probability when it is realized that $`\stackrel{~}{\rho }_{sat}`$, as defined here, already contains within it factors ($`|t|^2`$) representing the role of the tunneling probability.
This expression has a form similar to the standard tunneling theories which express the current as a product of the local densities of states of the two systems evaluated at a common point and a difference in the corresponding Fermi functions. Kawasaka et al. who studied the STM current through a Kondo resonance used as a starting point of their considerations the Tersoff and Hamann expression
$$I_{eq}𝑑\omega \left[f_t(\omega ^{})f_s(\omega )\right]\rho _t(\omega ^{})\rho _{sa}(\stackrel{}{R}_t,\omega )$$
(3.16)
according to which the current at zero temperature is related to the LDOS of the adsorbate plus metal electrons $`\rho _{sa}(\stackrel{}{R}_t,\omega )`$ at the position of the tip, where the local density of states is
$$\rho _{sa}(\stackrel{}{R}_t,\omega )=\frac{1}{\pi }\mathrm{I}m\stackrel{}{R}_t|G^R(\omega )|\stackrel{}{R}_t.$$
(3.17)
The LDOS is expressed in terms of unperturbed metal and adsorbate states by inserting $`_k|kk|+|aa|(1)`$ on both sides of $`G`$ in (3.17). This is strictly valid only for orthogonal orbitals, $`a|k=0`$. The four resulting terms, proportional to $`G_a`$, $`G_{ka}`$, $`G_{ak}`$, $`G_{kk^{}}`$, reflect the fact that the LDOS includes both the adsorbate and metal states perturbed by their mutual interaction. Inserting this expansion into Eq. (3.17) gives for $`\rho _{sa}(\stackrel{}{R}_t,\omega )`$
$`\rho _{sa}`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{I}m\{{\displaystyle \underset{a}{}}|\psi _a^{}|^2G_a^R+{\displaystyle \underset{ka}{}}\psi _\stackrel{}{k}G_{ka}^R\psi _a^{}+`$ (3.18)
$`+`$ $`{\displaystyle \underset{ka}{}}\psi _a^{}(\stackrel{}{R}_t)G_{ak}^R\psi _\stackrel{}{k}^{}+{\displaystyle \underset{kk^{}}{}}\psi _\stackrel{}{k}^{}G_{k^{}k}^R\psi _\stackrel{}{k}^{}\},`$ (3.19)
where the wavefunctions are evaluated at $`\stackrel{}{R}_t`$. The four terms correspond to the terms in the square brackets in (3.6). We can also rewrite $`\rho _{sa}`$ using the expressions for $`G_{ak}`$ and $`G_{kk^{}}`$ in Appendix B. We define $`\stackrel{~}{\psi }_a=\psi _a^{}+_k\psi _kG_k^0V_{ka}`$ and write
$$\rho _{sa}=\frac{1}{\pi }\mathrm{I}m\{\underset{k}{}|\psi _k|^2G_k^R+\underset{a}{}|\stackrel{~}{\psi }_a|^2G_{ka}^R\}.$$
(3.20)
Comparison of (3.20, 3.16) with ( 3.13, 3.15) shows the difference between our equilibrium limit and the transfer Hamiltonian method . The tunneling current and differential conductance in the equilibrium limit provides information about $`\stackrel{~}{\rho }_{sat}`$ – a local density of states modified by the tunneling matrix elements – rather than the LDOS. In the case when the tunneling takes place into distinct orbitals with different symmetry, $`\stackrel{~}{\rho }_{sat}`$ can be rather different from $`\rho _{sa}`$ and the statement that the STM is a measure of local density of states must be understood in this context.
### C The equilibrium tunneling current and the Fano lineshape
In this section, we evaluate $`\stackrel{~}{\rho }_{sat}`$ in Eq. (3.15) and write the equilibrium current using the approximations (2.15), (2.16), and (2.20) for the tunneling matrix elements. We note that the approximations do not require any specification of the substrate electronic structure. The final form of the current allows the discussion of the tunneling resonances in terms of the well established Fano lineshapes.
We introduce new quantities in terms of which the current is expressed. First, we define the “bulk” density-of-states (DOS) for the substrate and the tip as $`\rho _s(\omega )=_k\delta (\omega ϵ_k)`$ and $`\rho _t(\omega )=_p\delta (\omega ϵ_p)`$. The impurity width without the STM tip is defined as
$$\mathrm{\Gamma }_{as}(\stackrel{}{R}_0,\omega )=2\pi \rho _s(\omega )V_a^2(\stackrel{}{R}_0).$$
(3.21)
The adsorbate perturbation on the local density of conduction states at some lateral position between the tip and the adsorbate is discussed in terms of the unperturbed substrate Green’s function
$$G_0^+(\stackrel{}{r},\stackrel{}{r}^{};\omega )=\underset{k}{}\frac{\psi _k(\stackrel{}{r})\psi _k^{}(\stackrel{}{r}^{})}{\omega ϵ_k+i\eta }.$$
(3.22)
We define two dimensionless quantities related to the real and imaginary parts of the Green’s function
$$\mathrm{\Lambda }(\stackrel{}{R},\omega )=e^{Z/\lambda }\frac{\mathrm{R}eG_0^+(\stackrel{}{R},0;\omega )}{\pi \rho _s(\omega )}$$
(3.23)
and
$$\gamma (\stackrel{}{R},\omega )=e^{Z/\lambda }\frac{\mathrm{I}mG_0^+(\stackrel{}{R},0;\omega )}{\pi \rho _s(\omega )}.$$
(3.24)
These two functions carry the information about both the spatial extent of the metal electron perturbation at arbitrary $`\stackrel{}{R}`$ in the surface region due to a localized perturbation at $`\stackrel{}{R}_{}=0`$ and also the spatial resolution of the tip, as we will see later. We have included the exponential factor $`e^{Z/\lambda }`$ in the definition (3.23), (3.24), and (3.25) because we explicitly take the $`k`$-independent part of the exponential dependence on position to be part of the tunneling matrix elements $`V_a(Z_0)`$, $`t_a(\stackrel{}{R}_t,Z_0)`$, and $`t_c(Z_t)`$. We postpone further discussion of $`G_0^+(\stackrel{}{r},\stackrel{}{r}^{},\omega )`$ to the subsection III F.
Finally, we define a dimensionless quantity as the normalized density of the substrate states at a position $`\stackrel{}{R}`$ above the metal surface
$$\nu (\stackrel{}{R},\omega )=e^{2Z/\lambda }\frac{\rho _s(\stackrel{}{R},\omega )}{\rho _s(\omega )}=e^{2Z/\lambda }\frac{\mathrm{I}mG_0^+(\stackrel{}{R},\stackrel{}{R};\omega )}{\pi \rho _s(\omega )}.$$
(3.25)
The tunneling current $`I_0`$ into a clean metal is given by the first term in Eq. (3.13). The current $`I_0`$ for small bias $`(V/\varphi _s1)`$ can be written with the above definitions in a familiar form
$`I_0(\stackrel{}{R}_t,V)={\displaystyle \frac{2e}{h}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega \times `$ (3.26)
$`\times \rho _t(\omega ^{})\left[f_t(\omega ^{})f_s(\omega )\right]\rho _s(\omega )\pi t_c^2(\stackrel{}{R}_t)\nu (\stackrel{}{R}_t;\omega ).`$ (3.27)
Here, $`f_s(\omega )`$ and $`f_t(\omega )`$ are the substrate and STM tip Fermi functions, respectively, and $`\omega ^{}=\omega eV`$. The tip and substrate are assumed to have common chemical potential $`ϵ_{Fs}=ϵ_{Ft}=0`$ at zero bias $`eV=0`$ and we adopt the convention of measuring the energies in the substrate-adsorbate complex and in the tip from their respective Fermi levels at finite bias. The bias $`V`$ is measured with respect to $`ϵ_{Fs}`$ and is defined as positive when the chemical potential of the tip $`ϵ_{Ft}`$ is raised. The functions $`\rho _s`$ and $`\rho _t`$ are the substrate and tip densities of states, respectively.
It follows from Eq. (3.26) that the tunneling current $`I_0`$ is independent of temperature if $`\rho _t`$, $`\rho _s`$, and $`\nu (\stackrel{}{R}_t,\omega )`$ are independent of energy in the relevant energy range. If $`\rho _s`$ shows structure on the scale of the temperature $`T`$ while $`\rho _t`$ is constant, the current will depend on the temperature of the tip only, and vice versa. These statements are not limited to the case of clean metal surfaces, but also hold for the substrate with an impurity. The same is true for $`𝒢`$, the differential conductance.
The equilibrium current in the presence of the adsorbate is written by expressing $`\stackrel{~}{\rho }_{sat}`$ with the notation and approximations that lead to Eq. (3.26). We define a modified matrix element $`\stackrel{~}{t}_a(\stackrel{}{R}_t,\stackrel{}{R}_0;\omega )`$ for tunneling from tip to the adsorbate state as
$$\stackrel{~}{t}_a=t_a+\pi t_c\rho _s\mathrm{\Lambda }V_a.$$
(3.28)
The second term represents a coherent process of tip-to-surface tunneling, through-surface-propagation, and surface-to-adsorbate hopping. This is completely isomorphic with Fano’s coupling of an excited state (here the tip state) with the originally discrete state “modified by admixture of states of the continuum”. The reader is enthusiastically directed to the original Fano paper for further enlightenment on this point.
We introduce the Fano parameter $`q(\stackrel{}{R}_t,\stackrel{}{R}_0;\omega )`$ as
$$q=\frac{\stackrel{~}{t}_a}{\pi t_cV_a\rho _s}.$$
(3.29)
We will see later that this definition of $`q`$ makes the expression for differential conductance formally equivalent with the Fano formula in certain limits. It is rather straightforward now to evaluate $`\stackrel{~}{\rho }_{sat}`$ using (3.13) and (3.14) in (3.15). After rearranging the terms, we write the current $`I_{eq}(\stackrel{}{R}_t,\stackrel{}{R}_0,\omega )`$ in the presence of the adsorbate resonance as
$`I_{eq}(\stackrel{}{R}_t,\stackrel{}{R}_0,V)={\displaystyle \frac{2e}{h}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega \rho _t(\omega ^{})\times `$ (3.30)
$`\times \left[f_t(\omega ^{})f_s(\omega )\right]\rho _s(\omega )\pi t_c^2(\stackrel{}{R}_t)Y(\stackrel{}{R}_t,\stackrel{}{R}_0,\omega ),`$ (3.31)
with
$`Y=\nu +{\displaystyle \underset{a}{}}{\displaystyle \frac{\mathrm{\Gamma }_{as}}{2}}\left\{(\gamma ^2q^2)\mathrm{I}mG_a^R+2q\gamma \mathrm{R}eG_a^R\right\}.`$ (3.32)
In our approximation, the localized nature of the tip and the adsorbate enters through the position dependence of $`t_a(\stackrel{}{R}_t,Z_0)`$ and the substrate Green’s function $`G_0^+(\stackrel{}{R}_t,0;\omega )`$. The matrix element $`t_a`$ gives an exponentially decreasing amplitude with increasing tip-adsorbate distance and the substrate Green’s function gives decreasing amplitude due to the phase difference between electrons entering (or leaving) the surface at the adsorbate site and leaving (or entering) at $`(\stackrel{}{R}_{},z=0)`$ and also due to the exponential decay of the tip wavefunction with increasing $`k_{}`$). We note that, in the wide band limit for the substrate and with the tip near the surface above the adsorbate, $`\stackrel{~}{t}_at_a`$, since in this limit $`\mathrm{R}eG^+`$ and thus $`\mathrm{\Lambda }`$ vanish.
### D Nonequilibrium effects at stronger tip-surface coupling
We now generalize Eq. (3.30) for the equilibrium tunneling current – obtained in the lowest order in $`t_{ap}`$ and $`t_{kp}`$ – by including nonequilibrium effects. The general problem of tunneling for arbitrary relative strength between the tunneling amplitudes $`t_{ap}`$ and $`t_{kp}`$ and the hybridization matrix $`V_{ak}`$ and for finite bias is formulated in Eq. (3.7), but the expression is quite complicated to evaluate in practice. In a typical STM experiment, the tunneling matrix elements $`t_{ap}`$ and $`t_{kp}`$ are much smaller than than $`V_{ak}`$. We can expect the nonequilibrium effects to be important when, at small separations, the magnitude of the two tunneling matrix elements is not a negligible fraction of $`|V_{ak}|`$. However, we can always safely assume that $`|t_{ak}|`$, $`|t_{ap}|`$ are smaller than $`|V_{ak}|`$ in the STM experiments under all realistic conditions.
Therefore we make additional simplifications which are justified by these relations. First of all, we replace $`\stackrel{~}{V}_{ak}`$ by $`V_{ak}`$ inside Eq. (3.9) and (A5). We neglect the modifications to the tip and substrate wavefunctions, i.e. replace $`\stackrel{~}{G}_{kk^{}}`$ by $`\delta _{kk^{}}G_k^0`$ and $`\stackrel{~}{G}_{pp^{}}`$ by $`\delta _{pp^{}}G_p^0`$. We also neglect any deviations from thermal electronic distribution in the substrate and tip, i.e. we assume the validity of the fluctuation-dissipation theorem for the tip and substrate Green’s functions. On the other hand, when the tip-adsorbate coupling is not negligible with respect to the adsorbate-metal hybridization, the current into the resonance can be large enough to produce significant nonequilibrium electronic population on the adsorbate since the time scales for electron dissipation from the resonance into the metal and tip, respectively, are comparable. In this case, the fluctuation-dissipation theorem $`G_a^<=f_s\rho _a`$ is no longer valid for the adsorbate Green’s function and we must use the full nonequilibrium $`G_a^<(\omega )`$ instead of $`f_s(\omega )\rho _a(\omega )`$ in Eq. 3.30.
Under these assumptions, we find it convenient to write the total current with the nonequilibrium effects as $`I_{tot}=I_{eq}+\delta I_{non}`$, where $`I_{eq}`$ is formally given by (3.30) and $`\delta I_{non}`$ is
$`\delta I_{non}`$ $`=`$ $`{\displaystyle \frac{2e}{h}}{\displaystyle \underset{a}{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega \pi ^2\rho _t\rho _s^2t_c^2V_0^2\times `$ (3.33)
$`\times `$ $`(f_s\mathrm{I}mG_a^R+\pi G_a^<)(q^2+\gamma ^2),`$ (3.34)
where all adsorbate and substrate densities and Green’s functions are evaluated at energy $`\omega `$ and $`\rho _t`$ at $`\omega ^{}=\omega eV`$. We omitted the spatial arguments for simplicity. The bias dependence enters through the self-consistent solution of the adsorbate spectral density $`\rho _a(\omega )=\frac{1}{\pi }\mathrm{I}mG_a^R(\omega )`$ and the “lesser” Green’s function $`G_a^<(\omega )`$. In the case of noninteracting system, $`(U=0)`$, the spectral density does not depend on the bias and the only nonequilibrium (finite bias) effect is given by the difference between the equilibrium $`G_{a,eq}^<(\omega )=f_s(\omega )\rho _a(\omega )`$ and the nonequilibrium density of occupied states $`G_a^<(\omega )`$, as featured in $`\delta I_{non}`$.
On the other hand, the spectral density $`\rho _a(\omega )`$ of Kondo systems depends on the bias. This means that $`I_{eq}`$ also contains nonequilibrium effects and is different from the equilibrium current despite the subscript “eq” and its identical form. The effect of bias on the spectral function depends on the tip hybridization with the discrete impurity level and is similar to that of temperature for $`eVT_K`$ where it broadens the Kondo resonance. At larger biases the broadening increases further and a second peak may develop at the Fermi level of the tip, depending on the strength of the adsorbate-to-tip hybridization $`\mathrm{\Gamma }_{at}=2\pi \rho _tt_a^2`$ compared to $`\mathrm{\Gamma }_{as}=2\pi \rho _sV_a^2`$ for the relevant impurity orbital. In Fig. 2, we show for different bias voltages the spectral function and electron occupation of the resonance for a model Kondo system with $`\mathrm{\Gamma }_{at}`$ equal to $`10\%`$ of $`\mathrm{\Gamma }_{as}`$ and under an additional assumption that $`|t_{kp}||t_{ap}|`$. The model will be discussed in more detail in section IV B.
### E Differential conductance in the limit of large tip-surface separation
The differential conductance is obtained directly from (3.30) by differentiating it with respect to the bias, i.e. $`𝒢=dI/dV`$. We do this here under the assumption that the bias voltage is varied across a sufficiently narrow range so that the density of tip states may be taken constant. Under these assumptions the differential conductance $`𝒢_{eq}`$ is
$`𝒢_{eq}(\stackrel{}{R}_t,\stackrel{}{R}_0,V)`$ $`=`$ $`{\displaystyle \frac{2e^2}{h}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega \rho _t(\omega ^{})\left({\displaystyle \frac{f_t(\omega ^{})}{\omega }}\right)`$ (3.35)
$`\times `$ $`\rho _s(\omega )\pi t_c^2(\stackrel{}{R}_t)Y(\stackrel{}{R}_t,\stackrel{}{R}_0;\omega )`$ (3.36)
and for the clean metal
$`𝒢_0(\stackrel{}{R}_t,V)`$ $`=`$ $`{\displaystyle \frac{2e^2}{h}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega \rho _t(\omega ^{})({\displaystyle \frac{f_t(\omega ^{})}{\omega }})\times `$ (3.37)
$`\times `$ $`\rho _s(\omega )\pi t_c^2(\stackrel{}{R}_t)\nu (\stackrel{}{R}_t,\omega ).`$ (3.38)
These expressions neglect any changes to the tunneling barrier from the finite bias voltage. When these approximation are not justified, the conductance must be obtained by differentiating the expression for current (3.30) and (3.26). This is always the case for $`I_{tot}`$ of the previous section when nonequilibrium effects are important.
It is known from the Anderson “compensation theorem” that in the wide band limit ($`\rho _s(\omega )`$ is constant and unbounded), the presence of an impurity does not affect the conduction electron density of states at all. However, the density of conduction states is affected locally even in this limit, as can be seen by setting $`\psi _a=0`$ in Eq. (3.18). The perturbation of the conduction electrons is probed by the STM directly if $`t_a=0`$, i.e. when tunneling into the local state is absent due to either large tip-adsorbate separation or because of symmetry. Under these conditions if the STM conductance shows the signature of the local resonance, it is a result of the perturbation of the LDOS of conduction electrons.
An important feature of the final result (3.30) and (3.35) is that the role of the impurity resonance on the tunneling conductance is contained in the Green’s function $`G_a`$ of the local adsorbate state. It is then possible to separate the problem into two steps. First, the solutions for the adsorbate and substrate Green’s functions are found for a given system; then the tunneling conductance is calculated using the solution in the expression (3.30) or (3.35). Since the equations (3.30) and (3.35) have been obtained under very general assumptions, they can be used as a starting point in the study of a variety of tunneling problems with appropriate approximations for the Green’s functions and the tunneling matrix elements. We demonstrate this in the next section, where we first study noninteracting and then Kondo systems. If the approximations to the tunneling matrix elements employed here are too crude, the more general expression (3.7) or (3.15) must be used.
We find that the lineshape $`(𝒢`$ vs. $`V)`$ of the adsorbate resonance depends sensitively on (1) the relative strength of the tunneling matrix elements $`t_{ap}`$ and $`t_{kp}`$, (2) the perturbation of the conduction electron states, (3) the lineshape of the local resonance, i.e. the spectral function $`\frac{1}{\pi }\mathrm{I}mG_a`$ and $`\mathrm{R}eG_a`$, and on (4) temperature in the case when $`T\mathrm{\Gamma }`$, the width of the resonance. Since the observed perturbation of metal states and the tunneling matrix elements depend on the tip position, so will the lineshape.
### F The substrate Green’s function $`G_0^+`$ and perturbation of the conduction electrons: jellium surface
There are two ways in which the adsorbate state affects the tunneling conductance: (A) direct tunneling into the discrete state; (B) perturbation of the conduction electron states by the discrete state which consequently contributes to the tip-to-continuum tunneling current. Both contributions drop-off with increasing tip-adsorbate separation. The direct tunneling into the resonance is controlled by $`t_{ap}(\stackrel{}{R}_t,Z_0)`$ which is a function of the overlap between the tip and adsorbate wavefunction and thus decays exponentially with the distance. The perturbation of the continuum also vanishes at large distances from the adsorbate. However, its spatial extent shows a more complicated behavior and depends on the details of the electronic structure of the substrate in resonance with the broadened discrete state. It is anticipated that this contribution will show a significantly longer range influence than the decaying exponential, much in the spirit of Friedel oscillations.
The position dependence of the perturbation enters through the Green’s function $`G_0^+(\stackrel{}{R}_t,0;\omega )`$. We note that the imaginary part $`\gamma `$ appears explicitly in the expression for conductance, (3.35), while the real part $`\mathrm{\Lambda }`$ enters the definition of $`\stackrel{~}{t}_a`$. The STM-observable effects of the spatially dependent perturbation of the conduction electrons caused by the local impurity state are thus also controlled by $`G_0^+`$ (3.22). In our formulation, a non-trivial $`G_0^+`$ is intimately related to the dependence of the tunneling element $`t_{kp}(\stackrel{}{R}_t)`$ on the lateral tip position $`\stackrel{}{R}_{}`$. With this in mind we focus on the specific problem of tunneling as a function of the tip-adsorbate separation.
We consider a simple approximation for $`G_0^+`$ based on the assumption that in the relevant surface region the surface corrugations are smoothed out (jellium model) and both the Bloch and/or surface state $`\psi _k`$ (2.12) is given by
$$\psi _\stackrel{}{k}(\stackrel{}{r})e^{\kappa _sz}e^{i\stackrel{}{k}_{}\stackrel{}{\rho }}.$$
(3.39)
The states with the smallest $`\kappa _s`$ have the longest tail into the vacuum region and thus will be the most important ones in the tunneling process. These are the states with the smallest $`ϵ_k_{}`$. It is then reasonable to represent $`\kappa _s`$ in terms of the Taylor expansion around the minimum of $`ϵ_k_{}`$ with $`ϵ_k`$ equal to the bias. In most cases, it is reasonable to replace $`ϵ_k`$ by its Fermi level value. We expand $`ϵ_k_{}`$ around its minimum as $`ϵ_k_{}k_{}^2/2m^{}`$, and write $`\kappa _s=\lambda ^1+\lambda k_{}^2/2`$ plus higher order terms which we neglect. For states far from $`ϵ_{F_s}`$ and at small tip-surface separation, the expansion should be made around a different value of $`\lambda `$. For the purpose of this paper, it is sufficient to consider the Fermi level value $`\lambda `$. We then write
$$\psi _\stackrel{}{k}(\stackrel{}{r})e^{z/\lambda }e^{\lambda zk_{}^2/2}e^{i\stackrel{}{k}_{}\stackrel{}{\rho }}.$$
(3.40)
As we will show later, the second exponential $`e^{\lambda Z_tk_{}^2/2}`$ is a measure of the tunneling current carrying $`k_{}`$, the property that gives the STM tip its spatial resolution, and the third exponential, $`e^{i\stackrel{}{k}_{}\stackrel{}{R}_{}}`$, controls the dependence of the tunneling current on the lateral tip position.
With this approximation for Bloch states in the surface region the substrate Green’s function (3.22) is
$$G_0^+(\stackrel{}{R}_t,0;\omega )=e^{Z_t/\lambda }\underset{k}{}\frac{e^{\lambda Z_tk_{}^2/2}e^{i\stackrel{}{k}_{}\stackrel{}{R}_{}}|\psi _k(0)|^2}{\omega ϵ_k+i\eta }.$$
(3.41)
For the bulk band state propagation, it is easy to show using Eq. (3.24) that
$$\gamma (\stackrel{}{R}_t,\omega )=_0^1𝑑xJ_0(k_\omega R_{}\sqrt{1x^2})e^{\lambda Z_tk_\omega ^2(1x^2)/2}$$
(3.42)
and
$$\mathrm{\Lambda }(\stackrel{}{R}_t,\omega )=\frac{1}{\pi \rho _s(\omega )}𝒫_0^{2D}𝑑ϵ\rho _s(ϵ)\frac{\gamma (\stackrel{}{R}_t,ϵ)}{\omega ϵ},$$
(3.43)
where $`J_0`$ is the zeroth order Bessel function and $`k_\omega `$ is the wavevector of the substrate state of energy $`\omega `$. The normalized density of (STM-accessible) conduction states $`\nu (Z_t)`$ a distance $`Z_t`$ from the surface is
$$\nu (Z_t,\omega )=_0^1𝑑xe^{\lambda Z_tk_\omega ^2(1x^2)}.$$
(3.44)
In calculating $`\mathrm{\Lambda }`$, $`\gamma `$, and $`\nu `$ we assumed jellium-like dispersion relation $`\omega =k_\omega ^2/2m^{}`$ and use parabolic density of states $`\rho _s(\omega )=1\omega ^2/D^2`$. The incompatibility of the density of states with the dispersion relation is not important for the purpose of demonstrating the important band structure effects at this level of simplification.
Although the expressions (3.42)-(3.44) are valid for a very simple model of the surface, we believe they contain the most important features of more realistic bulk electronic structures. We now discuss these features beginning with $`\gamma `$, Eq. (3.42). At large $`Z_t`$, the dominant contribution to the integral in $`\gamma `$ comes from small values of the argument $`y`$ in $`J_0(y)`$. In this case, we use the mean value theorem to write (3.42) as $`\gamma (\stackrel{}{R}_t,\omega )=J_0(\overline{k}R_{})e^{\lambda Z_t\overline{k}^2/2}`$ where $`\overline{k}=\alpha k_\omega `$ with $`\alpha (0,1)`$. Clearly, $`\alpha 0`$ as $`Z_t\mathrm{}`$ and $`\gamma (\stackrel{}{R}_t,\omega )`$ is independent of the lateral tip position. Since at the same time $`\mathrm{\Lambda }0`$ and $`t_a0`$, the STM has no spatial resolution in this limit. As the tip moves closer to the surface the spatial resolution increases. In the limit $`Z_t=0`$, the integral in (3.42) can be evaluated and $`\gamma (\stackrel{}{R}_t,\omega )=j_0(k_\omega R_{})`$ where $`j_0`$ is the spherical Bessel function of zeroth order.
The resonance lineshape depends on the ratio of the direct tip-adsorbate tunneling amplitude ($`t_a`$) to the amplitude for tunneling into the perturbed metal states ($`t_c`$) and on the interference of conduction electrons scattering from the impurity. The latter contribution is represented here by the substrate Green’s function $`G^+`$. Therefore, the presence of an impurity on the surface can be sensed spectroscopically even if the direct tunneling into the resonance is negligible as is the case, for instance, of $`Ce/Ag(111)`$ . As our estimates for $`t_a(R_{})`$ and $`\gamma (R_t)`$ indicate, the direct tunneling matrix element $`t_a(R_{})`$ falls off much more rapidly with $`R_{}`$ than does $`\gamma `$. This is due to the limited spatial extent of the tightly bound impurity orbital. Therefore, the relative importance of tunneling into the perturbed continuum is likely to increase with the lateral tip-adsorbate separation.
Using the simple model for $`G^+`$, we show the typical length scales in Fig. 3(A). We plot $`t_a`$ (bold solid line) parameterized as $`t_ae^{R_{}/\alpha }`$ and normalized to one for $`R_{}=0`$ together with $`\gamma `$ evaluated at three different positions $`Z_t`$ above the surface and with $`k_\omega =1.2`$ Å<sup>-1</sup>. The exponential fall-off for conduction states at the Fermi level with work functions in the range $`45`$ eV would be characterized by $`\lambda =0.9`$ Å. The tightly bound discrete state will have larger decay constant and we parameterize it by $`\alpha =0.75`$ Å. Clearly, $`\gamma `$ decays much slower than $`t_a`$ at small tip-surface separations, but the difference in fall-off becomes smaller with increasing $`Z_t`$. We also see by comparison with the Bessel function $`J_0(k_FR_{})`$ (light dotted) that the spatial frequency decreases with increasing $`Z_t`$ and the oscillations eventually disappear entirely. This is due to the fact that smaller $`k`$-vectors have larger weight in the tunneling at greater $`Z_t`$ (see integral (3.42)). The real part $`\mathrm{\Lambda }`$, shown in panel (C) for the same k-vector, has a similar behavior. However, being a Hilbert transform of the imaginary part, the nodes in $`\mathrm{\Lambda }`$ appear at the positions of local extrema of $`\gamma `$ and vice versa. As we will see later, this property would lead to significant variations in the lineshape with $`R_{}`$ if it survived in the real electronic structure. Our results suggest that this is possible only at small $`Z_t`$. We discuss the band structure effects in the following section.
The panels (B) and (D) show the same as (A) and (C) but for smaller wavevector $`k_\omega =0.6`$ Å<sup>-1</sup>. Comparison between the right and left thus demonstrates the strong dependence of the substrate Green’s function on the wavevector itself, not just the product $`k_\omega R_{}`$. The two most significant features are that with decreasing $`k_\omega `$: (1) the frequency and damping of the oscillations with $`R_{}`$ decrease and (2) the dependence on $`Z_t`$ weakens. In comparing the two different energies, we assumed that the damping constant $`\lambda `$ (i.e. the tunneling barrier) is identical in the two cases. This would be the case in metals with identical work functions for states at the Fermi energy, in one of which the bottom of the band were closer to the Fermi level (smaller $`k_\omega `$). We note that $`k_\omega =1.2`$ Å<sup>-1</sup> corresponds to energy $`\omega =5`$ eV in the middle of the parabolic band with our parameterization. Therefore, the value of $`\mathrm{\Lambda }(\stackrel{}{R}_t=0)=0`$ at this energy but is negative for smaller energies, e.g. for $`k_\omega =0.6`$ Å<sup>-1</sup>, since in this case there are more high energy continuum states repelling the discrete state downward than low energy states pushing it up. We also see that the value of $`\mathrm{\Lambda }`$ at $`R_{}=0`$ can change sign with $`Z_t`$ depending on the energy $`\omega `$. We note that, since $`\mathrm{\Lambda }`$ enters the expression for $`q`$, the Fano parameter could also be negative and the asymmetry of the resonance lineshape could be reversed.
### G Electronic structure effects and the surface states on (111) noble metals
In the previous section we introduced a simple model of $`G^+`$ based on the unperturbed jellium surface. In general, more realistic behavior of $`G^+`$ can be obtained from electronic structure calculations. Here we discuss qualitatively the electronic structure effects with special attention to the (111) surfaces of noble metals frequently used in STM studies.
It is well known that (111) surfaces of noble metals contain Shockley surface states inside the projected two dimensional band gap that forms on these surfaces. Both the surface state and bulk wavefunction are given by the same general expression (2.12) outside of the metal surface. However their overall degree of localization at the surface is determined by the position of $`ϵ_{ss}`$, the surface state eigenvalue, with respect to the band gap edges. All other things being equal, the most localized surface state occurs when $`ϵ_{ss}`$ is at midgap. As $`ϵ_{ss}`$ moves towards either band edge, the extension of the evanescent oscillatory tail of the surface state wave function into the bulk increases, ultimately becoming identical to a periodic Bloch function when $`ϵ_{ss}`$ hits the band edge. From elementary normalization considerations, surface state extension into the bulk and amplitude at the surface, as reflected in the scale factor (or normalization constant) for the surface state tails (3.39) extending into vacuum, are intimately related; greater population within the bulk means lesser in the surface region. This surface state delocalization into the bulk allows for the local density of bulk states at the surface to greatly exceed that of the surface states, in which case the relative importance of the surface state in the tunneling current will be small near the surface. However, its importance increases with increasing distance from the surface because the bulk states with shorter wavefunction tails are eliminated from the tunneling. The surface state accounts for about $`50\%`$ of the total signal in typical STM tunnel junctions in Au(111) and is known to be responsible for the interference effects observed on these surfaces near edges, impurities, and in quantum corrals. It will also play a disproportionately important role in the resonance tunneling at large lateral tip-adsorbate distance because its contribution to $`G_0^+`$ does not decay as quickly as that for the bulk states.
We see the different behavior of the bulk and the surface states in the STM when we consider the propagator $`G_0^+`$ for the Shockley state. This is again given by Eq. (3.41). However, the k-sum now only extends over the 2-D wavevector $`k_{}`$. Assuming parabolic dispersion for the surface state, $`\gamma `$ is given by
$$\gamma (\stackrel{}{R}_t,\omega )=J_0(k_\omega R_{})e^{\lambda Z_tk_\omega ^2/2}$$
(3.45)
and
$$\mathrm{\Lambda }(\stackrel{}{R}_t,\omega )=\frac{1}{\pi \rho _s(\omega )}_0^{2D}𝑑ϵ\rho _s(ϵ)\frac{\gamma (\stackrel{}{R}_t,ϵ)}{\omega ϵ},$$
(3.46)
where, as before, $`k_\omega `$ is the 2-D wavevector of the substrate state corresponding to energy $`\omega `$. The contribution of the Shockley state to the normalized density of conduction states $`\nu (Z_t)`$ is given by
$$\nu (Z_t,\omega )=e^{\lambda Z_tk_\omega ^2}.$$
(3.47)
The propagator $`G_0^+`$ for the surface state is essentially equal to the Bessel function $`J_0(k_\omega R_{})`$ weighted by the exponential $`e^{\lambda Z_tk_\omega ^2/2}`$. Therefore the oscillations are not damped with increasing $`Z_t`$ and only their overall amplitude is diminished. Since the surface state on the noble metal surfaces (111) have a short $`k_F0.150.2`$ Å<sup>-1</sup>, its propagator will have a much longer spatial extent than that of the bulk states. The corresponding oscillations thus have a spatial period of about 10 times that of the Bessel function $`J_0`$ in Fig. 3(A) in agreement with the experimental observation of Friedel oscillations. The contribution of the surface state to the total current can carry information about an impurity on the surface over a long distance. The perturbation of the surface state by the impurity should persist over several tens of angströms.
It is also known that the spectral weight of the surface state decreases near surface imperfections. We expect the same to be true near the adsorbate. While we have explicitly taken into account the interaction of the conduction states with the discrete state “$`a`$” through the adsorbate Green’s function $`G_a`$, all other adsorbate-metal interactions, such as potential scattering of the conduction electrons from the adsorbate and hybridization of the outer shell adsorbate electronic states with the conduction electrons, are neglected in our model. In principle, these “residual” adsorbate-metal interactions can be included by modifying $`G^+`$ and $`t_{kp}`$. Although a realistic calculation of the system electronic structure is necessary to see the effect of the adsorbate on the behavior of $`G^+`$ around the adsorbate, we believe that it will not produce oscillatory behavior in $`G^+`$. In a typical metal, several bands with anisotropic dispersion relations $`ϵ_\stackrel{}{k}`$ contribute to $`G^+`$ giving rise to more complicated behavior with no single frequency. This will further reduce any oscillatory behavior seen in Fig. 3. At larger distance from the adsorbate, the band structure of the clean surface will be reestablished and, as a result, tunneling into the surface state.
Since the importance of the direct tunneling into the tightly bound impurity orbital “$`a`$” relative to the tunneling into the metal should be weak and decreases with increasing $`R_{}`$, it is useful to study the asymptotic behavior of the conductance $`𝒢`$ in the limit $`t_a=0`$. This is equivalent to replacing the Fano parameter $`q(\stackrel{}{R}_t,\stackrel{}{R}_0;\omega )`$ by $`\mathrm{\Lambda }(\stackrel{}{R}_t;\omega )`$ inside $`Y(\stackrel{}{R}_t,\stackrel{}{R}_0;\omega )`$, (3.32), in the expression for conductance. It then follows that, if the oscillations in $`G^+`$ persist, the lineshape should change with $`R_{}`$ and antiresonances should form at positions where $`\mathrm{\Lambda }^2>\gamma ^2`$. Using $`G^+=\pi \rho _se^{Z/\lambda }(\mathrm{\Lambda }i\gamma )`$ and
$`\mathrm{I}m\left\{G^+(\stackrel{}{R}_t,\stackrel{}{R}_0)G_a^RG^+(\stackrel{}{R}_0,\stackrel{}{R}_t)\right\}=`$ (3.48)
$`\pi ^2\rho _s^2e^{2Z_t/\lambda }\left\{(\gamma ^2\mathrm{\Lambda }^2)\mathrm{I}mG_a^R+2\mathrm{\Lambda }\gamma \mathrm{R}eG_a^R\right\}`$ (3.49)
we can write $`\mathrm{\Delta }𝒢_{eq}G_{eq}G_0`$ at zero temperature by replacing (-$`f(\omega )/\omega )`$ by $`\delta (\omega V)`$ and using $`t_c(Z)=t_0e^{Z/\lambda }`$ as
$`\mathrm{\Delta }𝒢_{eq}(V)`$ $`=`$ $`{\displaystyle \frac{2e^2}{h}}t_0^2\rho _t(0)V_a^2\times `$ (3.50)
$`\times `$ $`\mathrm{I}m\left\{G^+(\stackrel{}{R}_t,\stackrel{}{R}_0;V)G_a^R(V)G^+(\stackrel{}{R}_0,\stackrel{}{R}_t;V)\right\}.`$ (3.51)
We see that the resonance in the conductance is a result of an interference between different conduction states scattering resonantly from the impurity. Its long range behavior on the (111) noble metal surfaces is controlled by the surface states. Whether the resonance can be observed at the large distances ($`20`$ Å) depends on the spectral weight of the surface state and on its hybridization $`(V_a^2)`$ with the impurity orbital “$`a`$” (usually $`d`$ or $`f`$). Interesting spatial effects may be realized in system with suitable boundary conditions. We believe that Eigler’s quantum mirage of the Kondo resonance inside the elliptical corral falls into this category. Based on the results of the previous section, we do not expect “Friedel” oscillations at smaller distances and with period of a few angströms characteristic of the bulk k-vector, although we cannot completely rule these out for small $`Z_t`$.
## IV Discussion and examples
The equations (3.30) and (3.35) were derived under rather general assumptions. They are suitable as a starting point for numerical investigations given the necessary input from electronic structure calculations. In the rest of the paper, we discuss the implications of our theory for several specific cases of interest. In all of these cases we use our simple model for $`G^+`$ based on the jellium surface and the DOS given by $`\rho _s(\omega )=\rho _t(\omega )=1\omega ^2/D^2`$ with $`D=5`$ eV the band half width. In order to eliminate the exponential fall-off in the tunneling conductance with the tip-surface separation and the background distortions, we plot the normalized change in conductance due to the additional impurity defined as
$$\mathrm{\Delta }𝒢_{eq}(V)h(𝒢_{eq}(V)𝒢_0(V))/(2\pi e^2\rho _t(0)t_c^2),$$
(4.1)
where $`𝒢_{eq}`$ is given by (3.35) and $`𝒢_0`$ by (3.37). This is equivalent to replacing $`Y`$ by $`\mathrm{\Delta }Y=Y\nu `$ in the expression for $`𝒢`$.
Although we were motivated by the experimental observation of the Kondo resonance and this work is mostly applied to the tunneling through the Kondo resonance, we discuss many of the tunneling properties on the simple noninteracting model. We do this primarily because most of the STM observable characteristics of the Fano resonance are common to the single particle and Kondo resonances, despite the difference in processes that give rise to the two resonances. We wish to point out these general features on a model that is conceptually far simpler and more familiar to the surface science community than the Kondo model, and emphasize that the resonances can also be observed in systems with nonmagnetic impurities with a tightly bound orbital near the Fermi level. Finally, the connection with Fano result and the consequences of the spatial resolution of the STM become more transparent when the same noninteracting Anderson Hamiltonian is used.
### A Noninteracting adsorbate
We begin our discussion with an adsorbate-metal system described by the non-interacting Anderson model ($`U=0`$). The impurity resonance is characterized by its energy $`ϵ_0`$ and the width $`\mathrm{\Gamma }_{as}`$. The retarded Green’s function $`G_a^R(\omega )`$ for the adsorbate state is
$$G_a^R=(\omega ϵ_0\mathrm{R}e\mathrm{\Sigma }_a+i\mathrm{\Gamma }_{as}/2)^1$$
(4.2)
where $`\mathrm{\Gamma }_{as}(Z_0,\omega )`$ is defined in Eq. (3.21) and $`\mathrm{R}e\mathrm{\Sigma }_a(Z_0,\omega )=𝒫_k|V_{ak}(Z_0)|^2(\omega ϵ_k)^1`$ is the real part of the self-energy for the noninteracting Anderson model (not to be confused with the real part of the substrate Green’s function $`\mathrm{\Lambda }`$ given in Eq. (3.23)). Following Fano we now define the dimensionless energy parameter $`ϵ(\stackrel{}{R}_0,\omega )`$ by
$$ϵ=\frac{2(\omega ϵ_0\mathrm{R}e\mathrm{\Sigma }_a)}{\mathrm{\Gamma }_{as}}.$$
(4.3)
We neglect all nonequilibrium effects since they are likely to be insignificant for the noninteracting system under most experimentally realizable conditions. The differential conductance, in lowest order in $`t_c`$ and $`t_a`$, is given by Eq. (3.35) where $`Y`$ for the noninteracting system takes the form
$$YY_0=\nu +\frac{q^2\gamma ^2+2ϵ\gamma q}{1+ϵ^2}.$$
(4.4)
All terms are evaluated at energy $`\omega `$ and at the appropriate tip position. We note that $`Y_0(0,0,\omega )Y_{00}`$ characterizing the unphysical case of the STM tip in contact with surface at the position of the adsorbate (embedded in the surface) has the analytic form obtained by Fano
$$Y_{00}=\frac{(q+ϵ)^2}{1+ϵ^2},$$
(4.5)
although the inherent energy dependence of $`q`$ (through $`\mathrm{\Lambda }(\omega )`$) could distort the pure Fano character of the lineshape, even for this “almost atomic physics” STM example. At all other tip positions, the shape of the resonance will be described by the more general expression (4.4), which is equivalent to a sum of a “Lorentzian” plus a “Fano” profile,
$$Y_0=\frac{\nu ^{}}{1+ϵ^2}+\nu \frac{(q^{}+ϵ)^2}{1+ϵ^2}$$
(4.6)
with $`\nu ^{}\nu +q^2\gamma ^2q^2\gamma ^2/\nu `$ and $`q^{}q\gamma /\nu `$. Decomposition of $`Y_0`$ in this manner may be useful when analyzing experimental lineshapes. Note that the additional intrinsic energy dependences of $`\nu `$ and $`\gamma `$ could further distort the standard lineshape.
#### 1 Lineshape dependence on electronic structure and on the relative strength of $`t_a`$ and $`t_c`$
We first show (Fig. 4) the dependence of $`\mathrm{\Delta }𝒢_{eq}`$ on the ratio $`(t_a/t_c)`$ for $`\stackrel{}{R}_t=0`$ and a resonance at the Fermi level. The solid line corresponds to a resonance at the center of a parabolic band (symmetric around its center), i.e. $`ϵ_0=5`$ eV from the bottom of the band and the dashed line corresponds to a resonance at $`ϵ_0=2`$ eV from the bottom of the band. The two energies correspond to the band energy $`ϵ_k`$ with wavevector $`k=1.2`$ Å<sup>-1</sup> and $`k=0.6`$ Å<sup>-1</sup> in our jellium model (Fig. 3). The resonance width is $`\mathrm{\Gamma }_{as}=0.2`$ eV in both cases. At zero temperature, from Eq. (3.35) and (4.4) we write
$$𝒢_{eq}(V)=\frac{2e^2}{h}\rho _s(V)\pi \rho _t(0)t_c^2Y_0(V)$$
(4.7)
In order to make connection with the Fano result, we plot the conductance for small tip-metal separation with the tip above the adsorbate ($`\stackrel{}{R}_t=0`$). The lineshape $`Y_0`$ is then given by the Fano formula (4.5). The plotted quantity $`\mathrm{\Delta }𝒢_{eq}(V)`$ in Fig. 4 is then given by
$$\mathrm{\Delta }𝒢_{eq}(V)=\rho _s(V)\left(\frac{(q+ϵ)^2}{1+ϵ^2}1\right).$$
(4.8)
The Fano parameter $`q`$ depends not only on the ratio $`(t_a/t_c)`$ but also on energy and electronic structure. We see this most clearly in the first panel where $`t_a=0`$. The resonance placed at the center of the band produces a symmetric dip in $`\mathrm{\Delta }𝒢`$ characteristic of $`q=0`$, whereas the resonance at $`ϵ_0=2`$ eV has an asymmetric lineshape due to the negative contribution from $`\mathrm{\Lambda }`$ to $`q`$ (see Fig. 3). Its lineshape actually becomes symmetric at finite value of $`t_a`$. The value of $`t_a`$ inside each panel is given in units of $`t_c`$. The inset in the upper panel shows the model density of conduction states $`\rho _s`$ and the lower panel inset shows the spectral function $`\rho _a=\frac{1}{\pi }\mathrm{I}mG_a^R`$ (solid) and $`\mathrm{R}eG_a^R/\pi `$ (dotted) for the level at the center of the band.
As the strength of the direct tunneling $`t_a`$ increases with respect to $`t_c`$, the resonance develops its characteristic asymmetric shape and, eventually, at large $`t_a/t_c1`$ it acquires the shape nearly indistinguishable from that of the impurity spectral function $`\rho _a(\omega )`$. With increasing tip-adsorbate separation, the signal from the resonance must disappear as both the tunneling element $`t_a`$ and $`G^+(\stackrel{}{R}_t,\omega )`$ tend to zero. The differential conductance is then determined by the density of states of the clean surface. This property is not present in the Fano expression. We now discuss this behavior.
#### 2 Lineshape dependence on the tip-surface separation
Using the same model system as in the previous section with the resonance at the center of the band ($`ϵ_0=5`$ eV from the bottom of the band), we demonstrate the dependence on $`Z_t`$ (with $`\stackrel{}{R}_{}=0`$) in Fig. 5. We make the following model for the tunneling matrix element $`t_a(\stackrel{}{R}_t,\stackrel{}{R}_0)`$ and $`t_c(\stackrel{}{R}_t)`$. The exponential fall-off of the metal and adsorbate wavefunctions is controlled by different decay constants. The adsorbate state $`\psi _a`$ is tightly bound especially for narrow resonances of interest here. The conduction electron wavefunctions, on the other hand, typically belong to the outer $`s`$ or $`p`$ orbitals and have longer tails into the vacuum. As a consequence, the ratio $`t_a/t_c`$, and thus also the Fano parameter $`q`$, changes with $`Z_t`$. In order to incorporate this property, we use the matrix elements (2.20) and $`t_c=t_0e^{Z_t/\lambda }`$, where $`\alpha =0.75`$ Å, $`\lambda 0.9`$ Å, $`t_0=25`$ meV, and $`t_a=0.1t_c`$ at $`Z_t=2`$ Å. Under these conditions, the $`q`$ parameter tends to zero with increasing $`Z_t`$.
The panel (a) shows the normalized $`\mathrm{\Delta }𝒢_{eq}`$ for this model. The lineshape undergoes only moderate changes with $`Z_t`$ within the experimentally relevant range. We expect this to be a general property. In order to understand the behavior, we discuss the lineshape dependence on tip-surface separation conceptually in terms of two contributions: (1) different decay constants for the discrete $`\psi _a`$ and metal $`\psi _k`$ states at the Fermi level, and (2) different decay constant for metal states at $`ϵ_{Fs}`$ with different $`k_{}`$. We separate the observable consequences of these two effects in panels (b) and (c). The first contribution produces changes in $`q`$ due to the changing relative strength between $`t_a`$ and $`t_c`$. We demonstrate this in Fig. 5(b) where only this contribution is taken into account by setting $`Z_t=0`$ inside the substrate Green’s function $`G^+`$, or equivalently by setting $`\gamma =\nu =1`$ and $`\mathrm{\Lambda }=0`$ in Eq. (4.4). This limit does not correspond to a real situation and does not lead to the correct $`Z_t\mathrm{}`$ limit. It is shown here merely as an example of the contribution (1) to the $`Z_t`$ dependence of the tunneling conductance. With our parameterization, this case is identical with $`q`$ changing from $`q0.8`$ at $`Z_t=3`$ Å to $`q0.2`$ at $`Z_t=9`$ Å. As $`Z_t`$ increases further, $`q0`$ and the resonance becomes symmetric. However, we see that the normalized conductance $`\mathrm{\Delta }𝒢_{eq}`$ does not vanish in the limit $`Z_t\mathrm{}`$.
Fig 5(c) takes the second contribution (2) into account while leaving out the first one. We chose $`t_a/t_c=0=`$ constant, which would be the case if there were no direct tunneling into the discrete state. In this extreme limit, any changes in lineshape are a consequence of the varying weight that different $`k_{}`$ metal states play in the tunneling at different $`Z_t`$. This occurs because the $`k_{}=0`$ metal wavefunctions given by Eq. (3.40) have the greatest extension into the vacuum and as a result the spatial resolution of the tip decreases. Therefore the signature of the resonance in $`𝒢`$ decreases even after normalization of the current for different $`Z_t`$ as the ratio $`\mathrm{\Delta }Y/\nu 0`$ with $`Z_t\mathrm{}`$. Fig 5(a) shows the combined effect of the contribution and represents realistic conditions. It accounts correctly for the changing lineshape, as well as its disappearance. We again emphasize that realistic band structure is desirable for making quantitative statements.
Finally, we comment on the experimental issues. It is clear that the lineshape dependence on $`Z_t`$ will be observable only if it can be studied over a reasonably large range of $`Z_t`$, this being limited by the experimental resolution and detection capabilities. The most favorable case is one in which the direct tunneling $`t_a`$ into the resonance is strong at small $`Z_t`$, i.e. $`q`$ is large, and the difference in relevant decay constants for the adsorbate and metal electrons, $`a`$ and $`\lambda `$ respectively, is large. This is not the case in the experiments where $`q`$ is small. Therefore, we do not expect significant changes in the lineshape with $`Z_t`$ in these systems. Since our model is based on realistic parameterization, we expect the behavior shown in Fig. 5 to serve as a guide for order of magnitude estimates for the spectral dependence on $`Z_t`$.
The direct effect of the STM tip on the system and thus also on the lineshapes is not taken into account here. This issue is discussed in section IV B 3.
#### 3 Lineshape dependence on the lateral tip position
As we already discussed in section III F, the resonance lineshape depends on the relation between the spatial dependence of the direct tunneling and the propagation of the adsorbate-induced perturbation through the metal. Fig. 3 shows that the direct tunneling into the resonance $`(t_a)`$ is expected to fall off faster than the perturbation. Therefore, at large $`R_{}`$, the lineshape will be given by the $`t_a=0`$ limit conductance. We show the dependence of $`\mathrm{\Delta }𝒢_{eq}`$ on $`R_{}`$ in Fig. 6. We do this again for the model described in the previous section (Fig. 5) with $`t_a=0`$ and a resonance at the center of the band, i.e. $`ϵ_0=5`$ eV from the bottom of the band which corresponds to $`k_\omega =1.2`$ Å<sup>-1</sup> in Fig. 3. The solid line corresponds to $`Z_t=5`$ Å and the dashed line to $`Z_t=0`$ Å.
The unphysical case of $`Z_t=0`$ (dashed) is shown to emphasize the possible consequences of the oscillations in $`G_0^+`$ displayed in Fig. 3. We chose the lateral tip positions in the figure to coincide with the nodes and zeros of $`\mathrm{\Lambda }`$ and $`\gamma `$ to show the dramatic changes in the lineshape with $`R_{}`$ due to the oscillations in $`\mathrm{\Lambda }`$ and $`\gamma `$. Since the spatial decay of the oscillations is small at $`Z_t=0`$, the sequence of resonances and antiresonances appear in the range $`R_{}(0,10)`$ Å. The possibility for such antiresonances is discussed implicitly in the work of Kawasaka and explicitly by Schiller and Hershfield . However, this behavior is not observed in the experiments by Madhavan et al. and Li et al. due to the smoothing of the electronic structure with increasing distance from the surface that we discussed in III F.
In fact we would not expect the dramatic variations in line shape with $`R_{}`$ reported by Schiller and Hershfield to be observed. The reason is apparent from the behavior of $`G^+`$ as a function of $`Z_t`$ (Fig. 3). As the tip distance from the surface increases the oscillations are destroyed by the increasing weight of the lower frequency (small $`k_{}`$) components at larger $`Z_t`$ interfering destructively with those given by $`k_\omega `$. For this value of $`k_\omega `$, the oscillations are effectively damped when $`Z_t5`$ Å and the shape of the resonance does not change significantly as shown by the bold line in Fig. 6. We expect that band structure effects will suppress the oscillations even further.
We also find that the spatial extent of the resonance in the spectrum should decrease as the STM is retracted, as long as the signal is due to the bulk states. At $`Z_t=0`$ Å, the resonance is still visible at $`R_{}10`$ Å but only to about $`R_{}4`$ Å at $`Z_t=5`$ Å. This is a somewhat shorter distance than that found experimentally for $`Co/Au(111)`$ and $`Ce/Ag(111)`$ . Although the Fermi wavevector $`k_F1.2`$ Å<sup>-1</sup> used in Fig. 6 is close to the free electron value of $`k_F`$ for the noble metals, the disagreement is not surprising since we made no real attempt at realistic electronic structure description. Smaller values of $`k_F`$ would increase the spatial extent as would smaller values of $`Z_t`$ and $`\lambda `$.
Interestingly, the $`Z_t=0`$ \[dotted\] lineshape progression shown in Fig. 6 is qualitatively similar to the family of lineshapes that would be expected from surface state propagation, but with $`R_{}`$, the lateral tip-adatom separation rescaled upward by nearly an order of magnitude. This claim is based on the qualitative similarity between the bulk $`G_0^+`$ at $`Z_t=0`$ and the surface $`G_0^+`$. The bulk $`\gamma (R_{},ϵ_F)=j_0(k_FR_{})`$ \[the dashed curve in Fig. 3\] at $`Z_t=0`$ the analogous surface state $`\gamma (R_{},ϵ_F)=J_0(k_FR_{})`$ \[dotted curve in Fig. 3\] both exhibit long range oscillations unlike the bulk state at $`Z_t5`$Å. However since $`k_F(0.10.2)`$Å<sup>-1</sup> for the surface state band, $`J_0(k_FR_{})`$ shown in Fig. 3(A) for $`k_F=1.2`$Å<sup>-1</sup> should be plotted with this smaller $`k_F`$ when referring to actual noble metal surface state bands, in which case the observable $`R_{}`$-dependent lineshape evolution in Fig. 6 would still be representative, but with $`R_{}`$ rescaled by the factor $`1.2/0.15=8`$. From this it is easy to appreciate that the dramatic lineshape variations will occur mainly at very large lateral separations. Clearly, realistic electronic structure calculations are necessary to answer the more quantitative questions.
At large values of $`R_{}`$ and for broader resonances, an additional mechanism for distortion of the lineshape is possible if the relative change in the length of $`k_\omega `$ in the energy range given by the resonance width $`(\mathrm{\Gamma })`$ near $`ϵ_0`$ is large. In principle, this gives rise to the possibility of the oscillatory behavior of $`J_0`$ with respect to energy exhibiting itself in the spectrum. However, the conditions for this effect in $`𝒢(\omega )`$ would require the relative change $`\mathrm{\Delta }k_\omega /k_ϵ\sqrt{1+\mathrm{\Gamma }/2ϵ_0}\sqrt{1\mathrm{\Gamma }/2ϵ_0}`$ to be $`>2\pi /k_ϵR_{}`$. In this case, the argument of the Bessel function in (3.42) will vary over several periods starting at a small value at the bottom of the resonance. It is however clear that the width of the resonance and the distance from the impurity would have to be much larger than in the recent experiments with Kondo impurities. Moreover, these oscillations could only be observed through the surface state due to the damping of oscillations in the bulk.
#### 4 Temperature dependence of differential conductance
The temperature dependence of the tunneling conductance through a narrow resonance is an important issue especially in considerations of tunneling through the Kondo resonance which itself is temperature dependent. We revisit this issue in the section IV B 2 on Kondo effect. Here, we demonstrate the effect of the broadening in the Fermi function on the spectrum when the impurity spectral function is independent of temperature.
We assume the density of tip states and the substrate conduction electrons to be constant on the scale of the width $`\mathrm{\Gamma }`$ near the tip Fermi level and around the resonance. The temperature dependence of the conductance then follows from equation (3.35) with (4.4). We consider only the case of the tip directly above the adsorbate for simplicity and write
$$𝒢(V)=\frac{2e^2}{h}\pi \rho _t\rho _s_{\mathrm{}}^{\mathrm{}}𝑑\omega \left(\frac{f_t(\omega ^{})}{\omega }\right)t_c^2Y_0(\stackrel{}{R}_t,\omega ).$$
(4.9)
It follows from this expression that the differential conductance only depends on the temperature of the STM tip. This is a consequence of the assumption that the density of tip states is constant in the relevant energy range and the tunneling barrier is not modified by the bias. The differential conductance is then a function of the spectral density of the substrate states and is independent of their occupation. Generally, if the DOS in the STM tip varies significantly on energies $`\mathrm{\Gamma }_{as}`$, the substrate temperature would also enter. If the density of states in both the substrate and the tip were constant, no temperature dependence would be observed. It follows from Eq. (4.9) that the temperature $`T`$ must be $`\mathrm{\Gamma }`$, in order to have a significant effect on the conductance.
We show the temperature dependent $`𝒢_{eq}`$ in Fig. 7. The spectral density $`\rho _a(\omega )`$ with the narrow resonance at the Fermi level (the same system parameters as in Fig. 4) is independent of $`T`$. The STM tip is directly above the resonance with $`Z_t=5`$ Å. In the range $`T0.1\mathrm{\Gamma }`$, no temperature dependence is noticeable. However, when $`T(0.20.5)\mathrm{\Gamma }_{as}`$ the differential conductance shows a rather strong dependence on temperature, and at $`T\mathrm{\Gamma }`$ the sensitivity of the STM to the resonance disappears. The temperature dependence in Fig. 7 comes entirely from the broadening of the Fermi function of the tip.
### B Tunneling into Kondo resonance
In the previous section IV A, we discussed the STM conductance in tunneling through a noninteracting impurity \[$`U=0`$ in Hamiltonian (2.1)\], frequently referred to as the resonant level model (RLM). We now turn to the case of magnetic impurities and tunneling through a Kondo resonance. We begin with the case of a weak tip-metal coupling. However, for the Kondo systems this assumption is more restrictive than for the RLM model, and for this reason, we later take advantage of our nonequilibrium approach to account for the direct effect of the tip on the impurity spectral density, while still neglecting the tip’s effect on the metal states.
#### 1 Conceptual and theoretical approach
Since our earlier derivation of the current and conductance is valid for arbitrary interaction \[$`U0`$ in Eq. (2.1)\], the final results (3.30) and (3.35) also hold in the Kondo and mixed-valent regimes of the Anderson model (i.e. $`U\mathrm{\Delta }`$). The properties of the adsorbate enter through the Green’s function $`G_a`$. The problem is thus reduced to finding the one electron Green’s function $`G_a`$.
However, we first consider the tunneling for a spin $`1/2`$ ($`a\sigma `$) impurity in the Kondo limit, $`(ϵ_{Fs}ϵ_0)\mathrm{\Gamma }`$ and $`(ϵ_0ϵ_{Fs}+U)\mathrm{\Gamma }`$. The Kondo resonance has a very small weight and is due to spin fluctuations. The possible tunneling channels in this case are shown in Fig. 8 as processes (1) and (3). The system can be described by the Kondo Hamiltonian in this limit
$`H_s(Z_0)`$ $`=`$ $`{\displaystyle \underset{k\sigma }{}}ϵ_kc_{k\sigma }^{}c_{k\sigma }+{\displaystyle \underset{p\sigma }{}}ϵ_pc_{p\sigma }^{}c_{p\sigma }+`$ (4.10)
$`+`$ $`{\displaystyle \underset{kp\sigma }{}}\{t_{kp}(\stackrel{}{R}_t)c_{k\sigma }^{}c_{p\sigma }+\mathrm{H}.c.\}+`$ (4.11)
$`+`$ $`J_s{\displaystyle \underset{kk^{}\sigma \sigma ^{}}{}}(c_{k\sigma }^{}\stackrel{}{s}_{\sigma \sigma ^{}}c_{k^{}\sigma ^{}})\stackrel{}{S}+`$ (4.12)
$`+`$ $`J_t{\displaystyle \underset{pp^{}\sigma \sigma ^{}}{}}(c_{p\sigma }^{}\stackrel{}{s}_{\sigma \sigma ^{}}c_{p^{}\sigma ^{}})\stackrel{}{S}+`$ (4.13)
$`+`$ $`J_{st}{\displaystyle \underset{kp\sigma \sigma ^{}}{}}\{(c_{k\sigma }^{}\stackrel{}{s}_{\sigma \sigma ^{}}c_{p\sigma ^{}})\stackrel{}{S}+H.c.\}`$ (4.14)
where the first three terms were also present in the total Hamiltonian introduced in section II and describe the unperturbed metal and tip states and the coupling between the two. The remaining terms give rise to spin fluctuations in the presence of the magnetic impurity. The terms with couplings $`J_s`$ and $`J_t`$ correspond to the exchange interaction of the local spin with the substrate and tip electrons, respectively. The last term ($`J_{st}`$) corresponds to the effective tip-substrate exchange interaction in which charge is transported between the tip and the surface. This Hamiltonian can be obtained from $`H_{tot}`$ of section II using the Schrieffer-Wolf transformation which relates $`J_s`$, $`J_t`$ and $`J_{st}`$ to $`V_a`$ and $`t_a`$. For the symmetric Anderson model, $`J_s=4V_0^2/U`$, $`J_{st}=4t_aV_a/U`$, and $`J_t=4t_a^2/U`$. Using the continuity equation (3.1), the current is
$`I`$ $`=`$ $`{\displaystyle \frac{2e}{\mathrm{}}}\mathrm{I}m\{{\displaystyle \underset{kp\sigma }{}}t_{kp}c_{k\sigma }^{}c_{p\sigma }+J_{st}{\displaystyle \underset{kp\sigma \sigma ^{}}{}}c_{k\sigma }^{}\stackrel{}{s}_{\sigma \sigma ^{}}c_{p\sigma ^{}}\stackrel{}{S}+`$ (4.15)
$`+`$ $`J_t{\displaystyle \underset{pp^{}\sigma \sigma ^{}}{}}c_{p\sigma }^{}\stackrel{}{s}_{\sigma \sigma ^{}}c_{p^{}\sigma ^{}}\stackrel{}{S}\}.`$ (4.16)
The first term is identical with the first term in Eq. (3.3). In the lowest order of the tip-system couplings ($`t_{kp},J_{st},J_t`$), the third term does not contribute. The first term corresponds to the direct tip-substrate tunneling channel – process (1) in Fig. 8 – which includes the scattering of conduction electrons from the local moment. The second term corresponds to the direct tunneling into the magnetic impurity – process (3) in Fig. 8. We note that the spin flip scattering that gives rise to the Kondo effect is a higher order process. In the lowest order, the channel (1) and the spin-flip component of (3) do not give rise to interference because the final states have different spin states. The lowest spin-flip process that does interfere with (1) is of second order in $`J`$ and proportional to $`J_sJ_{st}`$.
In the limit of large tip-metal separation, equivalent to the condition $`(J_sJ_{st}J_t)`$, the third term in Eq. (4.15) as well as higher order contributions from $`J_{st}`$ are neglected and all other exchange processes are included in principle. This is equivalent to assuming that the state of the metal-adsorbate system is determined only by $`J_s`$ and is unaffected by the presence of the tip. Theoretically, the problem then reduces to finding the spectral properties of the system without the tip and using them in the expansion for tunneling via the two terms in Eq. (4.15).
As the system parameters move away from the Kondo limit – that is either $`ϵ_0`$ shifts towards $`ϵ_{Fs}`$ or $`U`$ becomes smaller – valence fluctuations appear. The Kondo resonance is then due to both the spin and charge fluctuations. The separate energy scale due to the spin fluctuations eventually disappears in the mixed-valent regime and the Kondo peak merges with the broad resonance centered at $`ϵ_0`$. In the intermediate regime, where both charge and spin fluctuations coexist on the impurity, another tunneling channel exists. This channel is denoted by (2) in Fig. 8. It also includes the contribution from higher order non-flip processes similar to (3). We study the system in this regime with the Hamiltonian defined in section II.
We adopt the slave-boson technique of Coleman and find the adsorbate Green’s function using the non-crossing approximation (NCA) . Following the theory of section III, the final expression for current in Eq. (3.7) is valid, as well as all the consequent steps and approximation in III. We insert the solution for the Green’s function $`G_a`$ of the $`(U=\mathrm{})`$ interacting system in Eq. (3.35). This is equivalent to including the three tunneling channels in Fig. 8 to lowest order in the tip-system couplings. We now turn to the discussion of the results based on this approach.
#### 2 Results for large tip-substrate separation
In order to model $`Co/Au(111)`$ studied both experimentally and theoretically , we choose a parameterization that gives the Kondo temperature $`T_K70`$ K appropriate for the system. Our simplified model has degeneracy $`N=2`$ with no orbital degeneracy, band width $`2D=10`$ eV, and the adsorbate level at $`ϵ_0=0.75`$ eV with the width $`2\mathrm{\Gamma }=1`$ eV (the width of a multiplet with an occupied level is $`N\mathrm{\Gamma }`$ rather than $`\mathrm{\Gamma }`$!). We show the corresponding spectral function and the real part of $`G_a`$ in the inset of Fig. 9.
Fig. 9 shows the spectral properties of the system using $`\mathrm{\Delta }𝒢_{eq}`$ over the whole energy range of the conduction band and for different values of $`q`$ at the Fermi level. The spectrum contains information about the broad resonance at $`ϵ_0=0.75`$ eV below the Fermi level, as well as the prominent feature due to the Kondo resonance at zero bias. We show the large bias voltage results only for completeness since we do not expect the STM experiments to be able to provide spectroscopic information about the system over the whole energy range shown.
The resonance lineshapes both in $`Co/Au(111)`$ and $`Ce/Ag(111)`$ correspond to small values of $`q`$. Madhavan et al. fitted the observed resonances to Fano lineshapes with $`q0.7`$. Our best fit would give approximately the same value of $`q`$. In the case of $`Ce/Ag(111)`$, the observed feature is an almost symmetric antiresonance corresponding to $`q0`$. Due to the contribution from the substrate electronic structure to $`q`$, its value cannot be directly used to make quantitative statements about the relative strength of the tunneling into the discrete state $`d`$ $`(f)`$ with respect to that into the continuum. However, in agreement with Li et al. and Lang we conclude that the STM probes mostly the $`sp`$ wave functions and the tunneling into the $`f`$-orbital is rather weak at the tip-adsorbate distances used in the $`Ce/Ag(111)`$ experiment. The resonance is mostly the result of interference between conduction electrons scattering from the impurity. The larger value of $`q`$ in $`Co/Au(111)`$ indicates stronger contribution from the coupling of the STM to the $`d`$ orbital. This is expected because the $`3d`$ orbital is not as tightly bound.
The recent work of Kawasaka et al. deals with the spatial and spectroscopic profiles of the Kondo resonance. They begin with the Tersoff-Hamann expression for current (3.16) and use the local density of states given by (3.18). They insert the self energy correction in the Green’s function $`G_a`$ due to the intra-adsorbate Coulomb correlations using perturbation theory $`(T>T_K)`$ and Yamada’s expansion in $`U`$ $`(T<T_K)`$ to study the temperature dependence in the whole temperature range. They neglect the additional temperature effects due to the Fermi surface broadening, replace $`(f/\omega )`$ by the delta function, and evaluate the conductance at the tip bias.
One of the main conclusions of their work is that the calculated temperature dependence of the resonance in the differential conductance is indicative of the temperature dependence of the Kondo resonance itself. They show results at the experimentally relevant low temperatures for $`Co/Au(111)`$ and $`Ce/Ag(111)`$ in the range of temperatures $`(T0.1T_K)`$. They find a rather weak temperature dependence, due entirely to the temperature dependence of the spectral function $`\rho _a`$. It is qualitatively the same and comparable in magnitude with that found in Fig. 7 for a temperature independent resonance of the noninteracting system for temperatures $`T\mathrm{\Gamma }`$. In our case, the temperature dependence in $`\mathrm{\Delta }𝒢_{eq}`$ is the consequence of the Fermi surface broadening in the STM tip. Therefore a careful deconvolution is necessary even at these low temperatures to extract information about the temperature dependence of the Kondo resonance. The other possibility is to eliminate variations in the Fermi surface broadening of the tip.
We show the temperature dependence for a Kondo system in Fig. 10. Since the validity of our approximation is limited to temperatures of order $`T_K`$ and higher, we show our results only in this temperature range. Panel (c) shows the temperature dependence one would observe with the tip at $`T=0`$ K and with varying substrate temperature, i.e. when only the temperature dependence of the spectral function is taken into account. Panel (b) assumes the substrate is at a constant temperature $`T=T_K`$, which determines the shape of the Kondo resonance, while the tip temperature is varied. We see that the two contributions produce a very similar broadening of the Fano resonance. Only a close look can uncover the difference. Panel (a) shows the combined effect when the tip and substrate are kept at a common temperature. Obviously, it would be difficult to determine the contribution from the broadening of the Kondo resonance.
In addition to the temperature effects just discussed, Kawasaka et al. also predicted the existence of weak, long range oscillations in the current as a function of the lateral tip position. The particular long wavelength, long range character of these predicted oscillations are a consequence of their assumption that the tip-to-metal tunneling takes place into the surface states of the (111) noble metal surfaces. The observed resonances are at variance with these expectations. On the other hand, the limited spatial extent of the resonance observed at lateral tip positions up to $`10`$ Å is consistent with the rapid spatial decay determined by the bulk $`G_0^+`$. No significant changes in the resonance lineshape are expected on this length scale (see discussion in section III G and Fig. 6). The surface state would, however, be responsible for lineshapes variations on larger length scales of order $`20`$ Å. The fact that no resonance is observed at such a distance from the impurity indicates that the surface state contribution is indeed weak. Since Friedel oscillations have been observed over long tip-impurity separations, we believe that a weak tunneling resonance most likely persists in the conductance over comparable distances, but more sensitive experiments are necessary. In this case, changes in the lineshape with $`R_{}`$ are expected. However, unlike Schiller and Hershfield , we do not expect variations in the lineshape due to the dominant contribution from the bulk states on the length scale of $`5`$ Å as discussed in section IV A 3.
#### 3 Nonequilibrium and hybridization effects at small tip-substrate separation
In typical STM experiments, the tip-substrate separation can be varied from the point of contact where the tunneling resistance $`R`$ is a few $`100k\mathrm{\Omega }`$ to distances where $`R1G\mathrm{\Omega }`$. Experimental constraints limit the STM usefulness to the near Fermi level spectroscopy – especially at small $`Z_t`$ – because of exponentially increasing tunneling currents with bias. However, it is likely to be possible to investigate the Kondo resonance – which only requires biases of the order of $`10`$ meV – with very small tip-adsorbate separations. It is therefore useful to analyze the physical consequences of the small tip-metal separation on the resonance in tunneling conductance.
In this case, nonequilibrium effects, as well as the tip-adsorbate interaction, become important in the spectroscopy of Kondo systems. First of all, as $`\mathrm{\Gamma }_{at}`$ increases and becomes a significant fraction of $`\mathrm{\Gamma }_{as}`$ at small distances, the tip-adsorbate hybridization will contribute to the width $`\mathrm{\Gamma }`$ of the resonance and to the renormalization of the level $`ϵ_0`$. As a result, the Kondo temperature, which depends sensitively on $`\mathrm{\Gamma }`$ and $`ϵ_0`$, will change. This could be particularly important for systems with very low bulk $`T_K`$, such as $`Fe/Au`$ with $`T_K1`$ K. The Kondo temperature for an impurity adsorbed on the surface of the metal is even lower than its bulk $`T_K`$ because the lower coordination number for the adsorbate makes the width $`\mathrm{\Gamma }`$ narrower. If $`T_KT`$, the Kondo resonance will not be observed. In certain systems and in the right temperature regime, it may be possible for the Kondo resonance to reappear at smaller tip-adsorbate distance as a result of the increased hybridization. This could also be achieved by incorporating the adsorbate into the top surface layer. The recent study of transition-metal impurities at the surface of gold did not find any sign of the Kondo effect in $`V`$, $`Cr`$, $`Mn`$, or $`Fe`$. We believe that, in the case of iron, this is due to the low $`T_K`$ and may be an example of a candidate system for the conditions discussed here. On the other hand, the smaller impurity-metal hybridization at the surface can lead to magnetic behavior for systems which are nonmagnetic in the bulk, such as $`Ni/Cu`$. There is a possibility for observing the transition between magnetic and nonmagnetic behavior on a single system induced either by embedding or by the proximity of the STM tip.
We show an example of the changing $`T_K`$ with hybridization in Fig. 11, where the spectral function $`\rho _a(\omega )`$ is plotted at zero bias as a function of the partial width $`\mathrm{\Gamma }_{at}`$, i.e. tip-metal separation for a model system. We choose $`D=5`$ eV, $`\mathrm{\Gamma }_{as}=0.25`$ eV, $`ϵ_a=1`$ eV, and $`T=30`$ K. The Kondo temperature for this model in the limit $`t_a=0`$, is $`T_K30`$ mK, much smaller than the temperature $`T`$. Therefore the Kondo resonance in the spectral function is very weak. When the tip is brought closer to the adsorbate, the Kondo resonance acquires more spectral weight as the Kondo temperature increases to $`T_K100`$ mK at $`\mathrm{\Gamma }_{at}=0.01\mathrm{\Gamma }_{as}`$, $`T_K0.2`$ K at $`\mathrm{\Gamma }_{at}=0.04\mathrm{\Gamma }_{as}`$, and $`T_K1.5`$ K at $`\mathrm{\Gamma }_{at}=0.25\mathrm{\Gamma }_{as}`$. Based on the justifications in Appendix A, we neglected the effect of the direct metal-tip interaction on the spectral function and treat the effect of the tip as another hybridization channel for the impurity state.
The effect of varying hybridization – due to either the presence of the STM tip or due to embedding or changing the environment of the adsorbate – on the tunneling resonance depends on the relation between $`T_K`$ and $`T`$. For instance, when $`T_KT`$, increased hybridization would produce stronger (sharper) tunneling resonance of the same width since the spectral weight in the Kondo resonance increases while its width remains almost constant until $`T_KT`$. Also the experimental resolution is limited by temperature in this regime. When, on the other hand $`T_KT`$, additional hybridization would not only increase the spectral weight in the Kondo resonance but also its width. The two cases should thus be distinguishable experimentally from each other and from the possible lineshape variations with $`Z_t`$ as a result of changing $`q`$.
The second important effect of the strong tip-adsorbate interaction is the breakdown of equilibrium relations at finite bias such as the fluctuation-dissipation theorem $`G^<(\omega )=f(\omega )\rho _a(\omega )`$ which consequently cannot be used in deriving the expression for current (3.30). This is true in general because the electron occupation of the tip, metal, and adsorbate electrons will no longer be thermal, i.e. will not be given by $`f_t(\omega )`$ and $`f_s(\omega )`$ but rather will be characterized by a nonequilibrium distribution produced by the injected tunnel electrons. The differential conductance is no longer proportional to the local density of states and cannot be obtained using Eq. (3.35). In Kondo systems, the hot electrons not only modify the electronic distribution on the impurity, but also modify the Kondo resonance itself.
This is shown Fig. 2 where the spectral function and density of occupied states is plotted at selected bias voltages in the limit of $`|t_{pk}||t_{pa}|`$. The impurity has a resonance at $`ϵ_0=1`$ eV below the Fermi level and total width $`\mathrm{\Gamma }=0.5`$ eV produced by the hybridization with both the tip and the substrate with the partial widths $`\mathrm{\Gamma }_{at}=0.1\mathrm{\Gamma }_{as}`$. Temperature is of the order of $`T_K`$ in this example. We see that the Kondo resonance broadens even more with increasing bias. This is due to the increase in the rate of incoherent scattering by $`eV_a/T`$ – an effect similar to temperature. At the same time, the electron occupation develops a non-thermal profile due to the large tip-adsorbate current. This is particularly visible for negative biases where the density of states is larger. Fig. 2(a) shows the equilibrium spectral function (dotted) and the electron population on the resonance (solid bold). The equilibrium spectral density is shown (dotted) in all panels. In addition, the spectral density (solid) and occupation (bold solid) are shown for the biases indicated in the figure by the labeled arrow. If the coupling to the tip were comparable with the metal-adsorbate hybridization, a double peak structure would develop. This has been predicted by Wingreen and Meir in the context of the nonequilibrium Kondo effect in quantum dots, also discussed by Plihal et al. . We see the onset of the double peak structure in panels (b)-(d) where a small cusp develops at the chemical potential of the tip. In summary, the bias has a significant effect on the spectral density even when $`\mathrm{\Gamma }_{at}0.1\mathrm{\Gamma }_{as}`$.
We show the tunneling current (right) and the corresponding differential conductance (left) in Fig. 12 for this model of Kondo impurity and for the STM geometry defined by $`\mathrm{\Gamma }_{at}=0.1\mathrm{\Gamma }_{as}`$ and $`|t_{pk}||t_{pa}|`$. The panels correspond to $`q=0.6`$, $`q=1.2`$, and $`q=2.4`$, respectively. The current on the right is calculated using $`I_{tot}=I_{eq}+\delta I_{non}`$ of section III D with $`\delta I_{non}`$ given by Eq. (3.33). The differential conductance $`𝒢_{tot}`$ on the left is obtained by differentiating the results displayed on the right. It cannot be calculated from the expressions in the text since the dependence of $`\rho _a(\omega )`$ and $`G_a^<(\omega )`$ on the bias voltage modifies the contributions to the current in a wide energy range, and $`𝒢_{tot}`$ is not related in simple terms to the properties at the Fermi level of the tip. We compare $`I_{tot}`$ and $`𝒢_{tot}`$ (circles) with $`I_{eq}`$ and $`𝒢_{eq}`$ (solid). The equilibrium quantities were obtained in the lowest order in $`t_{ap}`$ and $`t_{kp}`$ and with the equilibrium $`G_a`$ of Fig. 2(a).
We see that the broadening and disappearance of the Kondo resonance with increasing bias at strong tip-adsorbate coupling is weakened in the nonequilibrium calculation of $`𝒢_{tot}`$, because the contribution $`\delta I_{non}`$ compensates partially for the spectral function effect. The most consistent effect on the lineshape for various values of $`q`$ is the suppression of the resonance maximum and as a consequence a more symmetric appearance. This behavior is qualitatively different from both the hybridization effect and that of the changing Fano parameter $`q`$ – due to different decay constant of the impurity and metal states. Although the dependence of the tunneling resonance on the tip-substrate separation $`Z_t`$ will contain all three contributions, the hybridization and nonequilibrium contribution should only be important at extremely small tip-adsorbate separations. The variations in $`q`$ should not be important as it depends on the difference of the wavefunction tails and the two remaining contributions should leave distinguishable signatures in the tunneling resonance. It remains to be seen if the nonequilibrium condition play an important role in the tunneling between STM tip and Kondo impurity.
Finally, we note that the limit $`|t_{pk}||t_{pa}|`$ discussed here in connection with the nonequilibrium effects is not appropriate for the recent STM experiments, where $`t_{pa}`$ is likely much weaker than $`t_{pk}`$, even though the importance of $`t_{pa}`$ will increase relative to $`t_{pk}`$ with decreasing $`Z_t`$. We will address the more general case in a future work.
## V Conclusions
We used the Keldysh-Kadanoff method to study the spectroscopic features of adsorbate resonances in the STM tunneling experiments. The central results of our theory are the general expression for current Eq. (3.7) to all orders in the tunneling matrix elements and its equilibrium limit (3.15). Both are valid for arbitrary intra-adsorbate electron correlations and thus apply to both noninteracting ($`U=0`$), as well as magnetic (large $`U`$) systems. The discussion of the Fano resonances is based on additional approximations for the tunneling and hybridization matrix elements that lead to expressions (3.30) for the tunneling current and Eq. (3.35) for differential conductance in the lowest order in the tip-to-system tunneling matrix elements $`t_a`$ and $`t_c`$, i.e. at large tip-surface separation, and the nonequilibrium correction to the current in Eq. (3.33).
In the equilibrium limit, our theory of the tunneling current and conductance differs from the standard theories of STM, in that the dependence on LDOS is replaced by a tip specific quantity related to the LDOS (Eq. 3.13). The current is expressed entirely in terms of the adsorbate Green’s function $`G_a^R`$, the tip density of states, the tunneling matrix elements and the substrate Green’s function. We used the formulation to study the resonance lineshape as a function of temperature, tip-substrate separation, and lateral tip position. We summarize our findings as follows.
(1) The role of impurity state resonances in tunneling can be discussed in terms of two limiting cases. When direct tunneling across the barrier is weak, the resonance within the barrier provides an additional tunneling channel and can significantly enhance the tunneling current. This is the case of quantum dots in Coulomb blockade regime. If on the other hand, the tunneling into the continuum is strong, the presence of an “impurity” state could suppress the tunneling current due to the additional scattering of the conduction electrons in the metal from the impurity, i.e. increased resistance. The tunneling into the Kondo resonance in the recent STM experiments seems to be closer to the latter limit.
(2) The information about electron correlations and the Kondo resonance enters the tunneling problem through the impurity Green’s function $`G_a`$ while the position dependence of the conductance is controlled by the electronic structure of the metal.
(3) The spatial decay of the observed Fano resonance in the recent experiments is consistent with the conclusion that tunneling into the bulk conduction and hybridized $`sp`$ impurity states gives rise to most of the signal. The absence of any observable resonance at distances larger than $`10`$ Å suggests that the contribution from the surface state on Au(111) and Ag(111) to the resonant tunneling is not important in these experiments. However, the surface states are important in special cases, as indicated by the recent corral experiments in which the contribution of the surface states is enhanced by scattering from the walls of the corral.
(4) At large $`Z_t`$, tunneling into conduction states with $`k_{}`$ having the smallest parallel component corresponding to energy $`\omega =ϵ_k_{}+ϵ_k_{}`$ is strongly favored. This leads to the disappearance of the current oscillation vs. the lateral tip position due to tunneling into the bulk states which should otherwise be observed with period of about $`12`$ Å (corresponding to the bulk $`k_F`$) for typical experimental conditions. Therefore no oscillations in the lineshape should be observed on this length scale for typical tip-surface separation. The occurrence of an antiresonance with tip position at certain neighboring sites predicted by Schiller and Hershfield has its origin in these oscillations. It is a result of a simplified model for the surface electronic structure and we believe is unphysical. The small current oscillations predicted by Kawasaka et al. assume that the surface states are all-important in the spatial dependence of the resonance which seems to contradict the experimental results. We believe the surface state should be important at larger distances since on the (111) noble metal surfaces $`k_F0.150.2`$ Å<sup>-1</sup> and the corresponding period of oscillations is about $`20`$ Å (as observed experimentally as Friedel oscillations). We expect changes in the resonance lineshapes with this spatial period if the contribution from the surface state if the signature of the resonance is detectable at such distances.
(5) From the lineshapes observed in $`Co/Au(111)`$ and $`Ce/Ag(111)`$, we conclude that the direct tunneling into the discrete ($`d`$ or $`f`$) state is quite weak – stronger in $`Co/Au(111)`$. This confirms that the STM is mostly a probe of the delocalized $`sp`$ states and couples only weakly to the tightly bound $`d`$ or $`f`$ orbitals at typical tip-surface separations. Therefore the dominant process giving rise to the resonance lineshape is the tip-to-metal tunneling and interference between conduction electrons scattering from the local moment.
(6) The temperature dependence in differential conductance does not reflect only the temperature dependence of the Kondo resonance, but includes also the effect of Fermi surface broadening (mostly of the tip). The two contributions are of the same order of magnitude and qualitatively indistinguishable. Therefore, the temperature dependence in the differential conductance cannot be used directly to make conclusions about the temperature dependence of the resonance without controlling the tip Fermi surface broadening or without deconvolution.
(7) At small tip-surface separations, nonequilibrium effects as well as the additional tip-adsorbate hybridization may play an important role – especially in Kondo systems. The main effect of the finite bias voltage in this case is to broaden the Kondo resonance and produce nonequilibrium electron population on the adsorbate. The observed Fano resonance in differential conductance also broadens and its maximum is suppressed. The effect of the tunneling current on the Kondo resonance should thus leave a characteristic dependence of the lineshape on $`Z_t`$.
## VI Acknowledgments
We thank R. Celotta, E. Hudson, M. Stiles, and J. Stroscio for fruitful discussions and for helping us understand the experimental issues more clearly.
## Appendix A Adsorbate Green’s function
An important quantity in the theory of tunneling current through adsorbate resonances is the adsorbate Green’s function $`G_a`$. Using the equation of motion method, we find the expression for $`G_a`$ defined as the Fourier transform of
$$G_a(t,t^{})=iT_Cc_a(t)c_a^{}(t^{}).$$
(A1)
We do this for the case of arbitrarily strong coupling between the tip and the adsorbate with the intent to describe the nonequilibrium effects at finite bias. However, in this paper we consider the effect of the direct tip-metal interaction on $`G_a`$ to be weak and neglect it. Extension to the full description will be considered in future work. We believe the approximations adopted here capture the most important nonequilibrium effects.
We discuss both the noninteracting $`(U=0)`$ and interacting $`(U=\mathrm{})`$ model. Since the solution in both limits for the adsorbate-metal interaction is well known, we limit our discussion to the issues specific to the addition of the biased tip and refer reader to standard texts for the details. The $`(U=\mathrm{})`$ model is solved using the slave boson technique and NCA. In this approach a new pseudofermion is introduced by the transformation $`c_ac_ab^{}`$ in the Hamiltonian (2.1), where $`b^{}`$ is the creation operator for the slave boson. This eliminates the interaction term $`U`$ from the Hamiltonian as discussed by Coleman.
The time ordering operator $`T_C`$ orders the time according to their position on contour in the complex time plane . It is important to note that the equations must be first solved in the complex time domain and then analytically continued to the real axis as was pointed out in the previous appendix. The analytic continuation is performed before the Fourier transform, so we must be careful about how we deal with the Fourier transformed equations. Relevant details are in Appendix C. Here we discuss the equations of motion satisfied by the Fourier transforms of the time ordered Green’s functions and the analytic continuation is performed at the end according to the rules in Appendix C. All Green’s functions and self-energies in the following expressions are function of frequency $`\omega `$ and, therefore, we omit their argument to simplify the notation.
The Green’s function for the impurity state $`G_a`$ can be written in a standard way
$$G_a=(\omega ϵ_0\mathrm{\Sigma }_a)^1,$$
(A2)
using the self energy $`\mathrm{\Sigma }_a(\stackrel{}{R}_t,Z_0;\omega )`$. The solution for $`G_a`$ is thus reduced to finding $`\mathrm{\Sigma }_a`$. We first treat a closed shell or nonmagnetic open shell ($`V_{ka}U`$) adsorbate for which electron correlations can be neglected. We begin by considering the tip-substrate system without the adsorbate. We define $`\stackrel{~}{G}_{kk^{}}`$ and $`\stackrel{~}{G}_{pp^{}}`$ in analogy with $`G_a`$ (A1) as the Green’s functions of the metal and tip states, respectively, in the absence of the adsorbate. These are not identical with the Green’s functions $`G_{kk^{}}`$ and $`G_{pp^{}}`$ for the full system introduced in Appendix B. The bare metal-tip system is described by the Hamiltonian of section II with $`ϵ_0=U=t_{ap}=V_{ak}0`$. Using the equations of motion, we can write
$$(\omega ϵ_k)\stackrel{~}{G}_{kk^{}}=\delta _{kk^{}}+\underset{k^{\prime \prime }}{}\mathrm{\Sigma }_{kk^{\prime \prime }}\stackrel{~}{G}_{k^{\prime \prime }k^{}}$$
(A3)
and
$$(\omega ϵ_p)\stackrel{~}{G}_{pp^{}}=\delta _{pp^{}}+\underset{p^{\prime \prime }}{}\mathrm{\Sigma }_{pp^{\prime \prime }}\stackrel{~}{G}_{p^{\prime \prime }p^{}}$$
(A4)
where the self-energies are $`\mathrm{\Sigma }_{kk^{}}=_pt_{kp}G_p^0t_{pk^{}}`$ and $`\mathrm{\Sigma }_{pp^{}}=_kt_{pk}G_k^0t_{kp^{}}`$ and $`G_k^0=(\omega ϵ_k+i\eta _k)^1`$ and $`G_p^0=(\omega ϵ_p+i\eta _p)^1`$ are the Green’s functions for the clean metal and tip, respectively, without their mutual interaction. The coupled equations are shown diagrammatically in Fig. 13. The solutions for $`\stackrel{~}{G}_{kk^{}}`$ and $`\stackrel{~}{G}_{pp^{}}`$ can be formally written as the inverse of $`D_{kk^{}}=\delta _{kk^{}}(\omega ϵ_k)\mathrm{\Sigma }_{kk^{}}`$ and $`D_{pp^{}}=\delta _{pp^{}}(\omega ϵ_p)\mathrm{\Sigma }_{pp^{}}`$. We also define the adsorbate-metal and adsorbate-tip hybridization matrices modified by the tip-substrate interaction as $`\stackrel{~}{V}_{ka}=V_{ka}+_pt_{kp}G_p^0t_{pa}`$ and $`\stackrel{~}{t}_{pa}=t_{pa}+_kt_{pk}G_k^0V_{ka}`$. With these definitions and with $`\stackrel{~}{G}_{kk^{}}`$ and $`\stackrel{~}{G}_{pp^{}}`$ obtained through (A3) and (A4), the solution for the noninteracting $`\mathrm{\Sigma }_a`$ – shown diagrammatically in Fig. 13 – is formally given by
$$\mathrm{\Sigma }_a=\underset{kk^{}}{}V_{ak}\stackrel{~}{G}_{kk^{}}\stackrel{~}{V}_{k^{}a}+\underset{pp^{}}{}t_{ap}\stackrel{~}{G}_{pp^{}}\stackrel{~}{t}_{p^{}a}.$$
(A5)
The evaluation of the self energy $`\mathrm{\Sigma }_a`$ is rather complicated in the general case of strong tip-to-substrate coupling. We proceed with formulation of the general nonequilibrium theory for the tunneling current using this self energy (section III A) and then we discuss two limiting cases: (a) the equilibrium limit $`|t_{kp}|,|t_{ap}||V_{ka}|`$ (section III C) in which case the second term in (A5) is neglected and $`\stackrel{~}{V}_{ka}`$, $`\stackrel{~}{G}_{kk^{}}`$ replaced by $`V_{ka}`$, $`G_k^0\delta _{kk^{}}`$; and (b) the nonequilibrium case under the assumption $`|t_{kp}||t_{ap}||V_{ka}|`$, in which case we keep both terms in (A5) and replace $`\stackrel{~}{V}_{ka}`$, $`\stackrel{~}{t}_{ap}`$, $`\stackrel{~}{G}_{kk^{}}`$, and $`\stackrel{~}{G}_{pp^{}}`$ by $`V_{ka}`$, $`t_{ap}`$, $`G_k^0\delta _{kk^{}}`$, and $`G_p^0\delta _{pp^{}}`$. Section IV B 3 deals with tunneling through a Kondo impurity in this limit. The case (b) includes the effect of the increased hybridization of the discrete state due to the tip presence and the onset of nonequilibrium population on the adsorbate at finite bias.
In order to study these corrections in the limit (b), $`(|t_{pk}||t_{pa}||V_{ka}|)`$, we replace the Green’s functions $`\stackrel{~}{G}_{kk^{}}`$ and $`\stackrel{~}{G}_{pp^{}}`$ by the noninteracting ones, i.e. $`\stackrel{~}{G}_{kk^{}}=\delta _{kk^{}}G_k^0`$ and $`\stackrel{~}{G}_{pp^{}}=\delta _{pp^{}}G_p^0`$ and the modified $`\stackrel{~}{V}_{ka}`$ and $`\stackrel{~}{t}_{pa}`$ by $`V_{ka}`$ and $`t_{pa}`$. The self energy $`\mathrm{\Sigma }_a`$ then simplifies to
$$\mathrm{\Sigma }_a^0=\underset{k}{}|V_{ak}|^2G_k^0+\underset{p}{}|t_{ap}|^2G_p^0.$$
(A6)
The largest source of error in writing the approximate self energy is the neglect of the possibly significant interference effects at larger $`t_{kp}`$ as a result of the phase difference between $`\stackrel{~}{t}_{pa}`$ and $`t_{pa}`$. It is always reasonable to replace $`\stackrel{~}{V}_{ka}`$ by $`V_{ka}`$, as long as the adsorbate is on the surface rather than on the STM tip. These general case will be the topic of a future study. If the tip distance from the adsorbate is much larger than the adsorbate-metal separation, so that $`t_{ap}V_{ak}`$, the self-energy is well described by the first term only. In such a case, the STM does not strongly modify the studied system. It is then reasonable to characterize the system without the presence of the STM tip and then consider the tunneling.
Finally, we discuss the Green’s function $`G_a`$ in the limit (b) for a Kondo impurity which is likely to show stronger dependence on the bias and tip interaction. We find the nonequilibrium Green’s function $`G_a`$ under the same assumption that lead to $`\mathrm{\Sigma }_a^0`$ for the noninteracting Anderson Hamiltonian. We solve the interacting system in the limit of $`U=\mathrm{}`$ using the NCA approximation, shown diagrammatically in Fig. 14. The self energy is not a simple sum of the two contributions from the metal and tip as it was in the noninteracting system, because the occupation of the resonance is limited to one electron and the hybridization is now correlated – formally through the slave boson Green’s function $`B(\omega )`$. The two coupled equations in Fig. 14 are solved selfconsistently. The expressions for the NCA self energy is obtained in a standard with the help of the diagrams in Fig. 14
## Appendix B Equations of motion for $`G_{pk}`$ and $`G_{pa}`$
In this appendix, we find the solution for $`(_{kp}t_{kp}G_{pk}+_{ap}t_{ap}G_{pa})`$ entering the expression for tunneling current (3.3) for the general case of arbitrary tip-system coupling. Ultimately, the interesting regime in connection with typical STM experiments is one in which $`t_{kp},t_{ap}<V_{ak}`$. However, we want to be able, in principle, to study the system when the coupling of the STM tip to the system and the tunneling current are strong. This creates nonequilibrium occupation on the adsorbate resonance and modifies the spectroscopic properties of the system. We therefore proceed by deriving the most general expression valid for arbitrary coupling strength $`t_{kp}`$ and $`t_{ap}`$ and discuss an approximation (b) that allows us to take into account the most important nonequilibrium effects as described in the previous appendix. For this purpose we introduce Green’s function $`G_{pa}(\omega )`$, $`G_{pk}(\omega )`$, $`G_{kk^{}}(\omega )`$, $`G_{pp^{}}(\omega )`$, and $`G_{ka}(\omega )`$ as the Fourier transform of
$$G_{ij}(t,t^{})=iT_Cc_i(t)c_j^{}(t^{}).$$
(B1)
We now turn to the equations of motion for the Green’s functions relevant for the tunneling current. The following expressions are valid for arbitrary interaction $`U0`$ and the nature of the intra-adsorbate interactions are contained fully in the solution for $`G_a`$ discussed in the previous appendix. The first term in the current (3.3) contains the tip-adsorbate propagator which satisfies
$$(\omega ϵ_p)G_{pa}=t_{pa}G_a+\underset{k}{}t_{pk}G_{ka}.$$
(B2)
It is expressed in term of $`G_a`$ already solved within a given approximation in the previous appendix through (A2) and in terms of the metal-adsorbate Green’s function
$$(\omega ϵ_k)G_{ka}=V_{ka}G_a+\underset{p}{}t_{kp}G_{pa}.$$
(B3)
The last two equations are coupled and need to be solved self-consistently. We do this by substituting Eq. (B3) for $`G_{ka}`$ in (B2), and vice versa. The solutions are then expressed in terms of $`G_a`$, $`\stackrel{~}{G}_{pp^{}}`$, and $`\stackrel{~}{t}_{pa}`$ discussed in the previous appendix as
$$G_{pa}=\underset{p^{}}{}\stackrel{~}{G}_{pp^{}}\stackrel{~}{t}_{p^{}a}G_a.$$
(B4)
We will also need the solution for $`G_{ak}`$. The tip-induced correction to $`V_{ka}`$ contributes to the phase of $`V_{ka}`$, as well as its magnitude, and could thus affect the lineshape significantly in the strong coupling limit. But it should be particularly weak when $`t_{kp},t_{pa}V_{ka}`$ and it will be safe to ignore it. We write
$$G_{ak}=G_a\underset{k^{}}{}\stackrel{~}{V}_{ak^{}}\stackrel{~}{G}_{k^{}k}.$$
(B5)
The second term in (B3) is negligible when the tip-adsorbate separation is much larger than the adsorbate-metal separation. Neglecting this term is equivalent to replacing $`\stackrel{~}{G}_{pp^{}}G_p^0\delta _{pp^{}}`$ in (B4) and $`\stackrel{~}{G}_{kk^{}}G_k^0\delta _{kk^{}}`$, $`\stackrel{~}{V}_{ka}V_{ka}`$ in (B5). The tip-adsorbate Green’s function $`G_{pa}`$ is then expressed entirely in terms of $`G_a`$ and the unperturbed conduction electron Green’s functions.
The second term in (3.3) contains the tip-metal propagator $`G_{pk}`$, which satisfies
$$(\omega ϵ_p)G_{pk}=\underset{a}{}t_{pa}G_{ak}+\underset{k^{}}{}t_{pk^{}}G_{k^{}k}.$$
(B6)
It is expressed in terms of $`G_{ak}`$ (B5) discussed in the previous paragraph and in terms of $`G_{k^{}k}`$, the Green’s function for the substrate conduction electrons
$$(\omega ϵ_k^{})G_{k^{}k}=\delta _{kk^{}}+\underset{a}{}V_{k^{}a}G_{ak}+\underset{p}{}t_{k^{}p}G_{pk}.$$
(B7)
We see that $`G_{kk^{}}`$ couples to $`G_{pk}`$, (B6), and also to $`G_{ak}`$, (B5), already solved in terms of $`G_a`$ and $`\stackrel{~}{G}_{kk^{}}`$. The last two equations can be solved self-consistently to give
$$G_{k^{}k}=\stackrel{~}{G}_{k^{}k}+\underset{k_1k_2}{}\stackrel{~}{G}_{k^{}k_1}\stackrel{~}{V}_{ak_1}G_a\stackrel{~}{V}_{k_2a}\stackrel{~}{G}_{k_2k}.$$
(B8)
and
$$G_{pk}=\underset{p^{}}{}\stackrel{~}{G}_{pp^{}}(t_{p^{}k}+\underset{ak^{}}{}\stackrel{~}{t}_{p^{}a}G_a\stackrel{~}{V}_{ak_1})G_k^0.$$
(B9)
For the purpose of analytic continuation, it is important to keep track of the order in which the Green’s functions appear in the product in the above equations. The “lesser” Green’s functions are then obtained according to rules stated in Appendix C. The equilibrium limit of the theory is achieved by neglecting the last term in (B7) along with the equivalent approximations for $`G_a`$ and $`G_{ak}`$ discussed above. This removes the self-consistency requirement and neglects the effect of the tip on the substrate conduction electrons, but not on the tunneling current. The solution for $`G_{kk^{}}`$ is then identical to that of the system without the tip.
## Appendix C Rules for dealing with nonequilibrium Green’s functions in frequency space
In this appendix we review the process of analytic continuation of the complex time contour expression to integrals on the real time axis. We follow Langreth’s generalization of Kadanoff-Baym’s method described in detail by Haug and Jauho . Four Green’s functions appear in the nonequilibrium theory “lesser” $`G^<`$, “greater” $`G^>`$, retarded $`G^R`$, and advanced $`G^A`$ . We frequently need to find the retarded (advanced) and “lesser” Green’s function corresponding to Green’s function $`A`$ time ordered in the complex time plane expressed as a product of $`N`$ time ordered functions $`B\mathrm{}Z`$ , i.e.
$$A(t,t^{})=_Cd\tau _1\mathrm{}d\tau _2B(t,\tau _1)..Z(\tau _2,t^{})$$
(C1)
where all functions are assumed fermion-like. The desired expressions analytically continued onto the real time axis are
$$A^{R(A)}(t,t^{})=_{\mathrm{}}^{\mathrm{}}𝑑\tau _1\mathrm{}𝑑\tau _2B^{R(A)}(t,\tau _1)\mathrm{}Z^{R(A)}(\tau _2,t^{})$$
(C2)
which consists of only one term and
$`A^{\genfrac{}{}{0pt}{}{>}{<}}(t,t^{})={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau _1\mathrm{}𝑑\tau _3\mathrm{}𝑑\tau _2`$ (C3)
$`[\mathrm{}+B^R(t,\tau _1)\mathrm{}C^{\genfrac{}{}{0pt}{}{>}{<}}(t,\tau _3)\mathrm{}Z^A(\tau _2,t^{})+\mathrm{}]`$ (C4)
where each of the $`N`$ terms in the integral has exactly one function of the type $`f^{\genfrac{}{}{0pt}{}{>}{<}}`$, all functions to the left (right) of it are retarded (advanced), and each of the terms has the $`f^{\genfrac{}{}{0pt}{}{>}{<}}`$ in a different position.
We are dealing here with the case of time independent perturbations (steady state current). The double time propagators then only depend on the time difference ($`tt^{}`$). The equations (C2), (C3) can be Fourier transformed and we can write the time ordered Green’s function $`A(\omega )`$ as a simple product
$$A(\omega )=B(\omega )\mathrm{}Z(\omega ).$$
(C5)
The rules for writing the expression for the “lesser” and retarded function in the frequency space are then directly carried out by Fourier transforming the equations (C2) and (C3). Leaving out the frequency arguments, we write
$$A^{R(A)}=B^{R(A)}\mathrm{}Z^{R(A)}$$
(C6)
and
$$A^{\genfrac{}{}{0pt}{}{>}{<}}=\mathrm{}+B^R\mathrm{}C^{\genfrac{}{}{0pt}{}{>}{<}}\mathrm{}Z^A+\mathrm{}$$
(C7)
We note that it is important to keep track of the order in which the function $`B\mathrm{}Z`$ appear in the time integral when the Fourier product is formed.
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# Presupernova collapse models with improved weak-interaction rates
\[
## Abstract
Improved values for stellar weak interaction rates have been recently calculated based upon a large shell model diagonalization. Using these new rates (for both beta decay and electron capture), we have examined the presupernova evolution of massive stars in the range $`15`$$`40\mathrm{M}_{}`$. Comparing our new models with a standard set of presupernova models by Woosley and Weaver, we find significantly larger values for the electron-to-baryon ratio $`Y_e`$ at the onset of collapse and iron core masses reduced by approximately $`0.1\mathrm{M}_{}`$. The inclusion of beta-decay accounts for roughly half of the revisions, while the other half is a consequence of the improved nuclear physics. These changes will have important consequences for nucleosynthesis and the supernova explosion mechanism.
\]
The late stages of massive stellar evolution are strongly influenced by weak interactions which act to determine the core entropy and electron-to-baryon ratio, $`Y_e`$, of the presupernova star. Electron capture reduces the number of electrons available for pressure support, while beta decay acts in the opposite direction. Both processes generate neutrinos which, for densities $`\rho 10^{11}`$g/cm<sup>3</sup>, escape the star carrying away energy and entropy from the core.
Electron capture and beta decay during the final evolution of a massive star are dominated by Fermi and Gamow-Teller (GT) transitions. While the treatment of Fermi transitions (important only in beta decays) is straightforward, a correct description of the GT transitions is a difficult problem in nuclear structure. In their pioneering work on the subject Fuller, Fowler and Newman (FFN) estimated electron-capture rates assuming a single GT resonance. The properties of this resonance were derived on the basis of the independent particle model, supplemented by Fermi contributions and experimental data for low-lying transitions, whenever available. These authors also noted the importance of the ‘backresonances’ for beta decay. These are excited states in the decaying nucleus which are connected by strong GT transitions to low-lying states in the daughter nucleus and, by thermal population and with increased phase space, can significantly contribute to the stellar beta decay rates.
Recent experimental data shows that the GT distributions in nuclei are quenched, compared to the independent particle model value, and strongly fragmented over many states in the daughter nucleus. Both effects are caused by residual interaction among the valence nucleons and an accurate description of these correlations is essential for a reliable evaluation of the stellar weak interaction rates due to the strong phase space energy dependence, particularly of the stellar electron capture rates. The shell model is the only known tool to reliably describe GT distributions in nuclei . Its application to iron mass nuclei in the middle of the $`pf`$-shell as required in the presupernova collapse, however, has long been inhibited due to the extremely large model space dimensions involved. After significant progress in shell-model programming and hardware development the situation has changed very recently and in Ref. it has been demonstrated that state-of-the-art diagonalization studies, typically involving a few 10 million configurations, are indeed able to reproduce all relevant ingredients (GT<sub>±</sub> strength distributions for changing protons (neutrons) into neutrons (protons), level spectra and half-lives) and hence have the predictive power to reliably calculate stellar weak interaction rates. This program has recently been finished and stellar weak-interaction rates for nuclei with $`A=45`$–65 have been calculated based on the shell-model results, supplemented by experimental data, wherever available. The shell-model rates have been discussed and validated in . It has been found that for $`pf`$-shell nuclei the shell-model electron-capture rates are smaller than the FFN rates by, on average, an order of magnitude, for the reasons explained in . The situation is different for the beta decay as the shell model and FFN rates are of the same magnitude for the most relevant nuclei to be identified below.
To study the influence of the shell model rates on presupernova models we have repeated the calculations of Woosley and Weaver (henceforth WW) keeping all the stellar physics, except for weak rates, as close to the original studies as possible. The present calculations have incorporated the new shell-model weak interaction rates (including electron capture, positron emission, and beta decay) for nuclei with mass numbers $`A=45`$–65, supplemented by rates from Oda et al. for lighter nuclei. In practice, the weak rates for these lighter nuclei were not very important for determining the presupernova structure, but only dominate prior to silicon burning. We also note that, for $`sd`$-shell nuclei with $`A=1739`$, the FFN rates agree rather well with the shell model rates previously determined by Oda et al..
The earlier calculations of WW, to which we compare, used the FFN rates for electron capture and an older set of beta-decay rates taken from and . Shortly after the models of WW were calculated, it was recognized that these older beta-decay rates were inadequate and that use of the larger values from FFN would appreciably alter the results . These are the first models of the WW variety to incorporate realistic beta-decay rates and electron-capture rates in a complete stellar model, even though beta decays were included in some other models of massive stellar evolution . In a separate paper , we will present the comparison of models that use the full FFN rate set and the new rate set, and in a future paper we will examine the evolution of stars of other mass and metallicity.
Fig. 1 shows the late evolution, following core oxygen burning, of the central temperature and density in $`15\mathrm{M}_{}`$ and $`25\mathrm{M}_{}`$ stars as well as the central value of the electron-to-baryon ratio, $`Y_e`$. Time is measured here backwards from the time of iron core collapse, which is arbitrarily t = 0. Silicon burning ignites with a mild “flash” and the core becomes convective and expands. For the $`25\mathrm{M}_{}`$ star the temperature and density trajectories in the new and old models are similar, but the calculations with shell model rates have significantly larger values of $`Y_e`$. This difference persists throughout the iron core, not just at its center. Larger values of $`Y_e`$ also result in the $`15\mathrm{M}_{}`$ star, but there even the density and temperature structures of the presupernova star are appreciably altered.
To understand the origin of these differences, we will now explore the role of the weak-interaction rates in greater detail. Weak processes become particularly important in reducing $`Y_e`$ below 0.50 after oxygen depletion ($`10^7`$ s and $`10^6`$ s before core collapse for the $`15\mathrm{M}_{}`$ and $`25\mathrm{M}_{}`$ stars, respectively) and $`Y_e`$ begins a decline that becomes precipitous during silicon burning. Initially electron capture occurs much more rapidly than beta decay. Since the shell model rates are generally smaller than the FFN electron capture rates, the initial reduction of $`Y_e`$ is smaller in the new models. The reduction is less pronounced during silicon shell burning both because the evolution time scale becomes quite short and because nuclei near the valley of beta-stability, with smaller weak-interaction rates, have already been produced (from a composition that initially had $`N=Z`$). During this period the core matter is composed of nuclei with $`A<65`$ which are carried in the calculation.
An important feature of the new models is that beta decay becomes temporarily competitive with electron capture after silicon depletion in the core and during silicon shell burning. That this would occur was foreseen by on the basis of one-zone models. Here we see it occurring in a complete stellar model. Moreover, the new electron capture rates are smaller than FFN and thus offer less resistance to beta decay. Dynamic weak equilibrium, in the sense described by , thus occurs at larger values of $`Y_e`$. Interestingly, by the time the iron core is actually collapsing, weak equilibrium no longer exists. The increase in density closes the phase space for beta decay and electron capture again predominates by a large factor. It is this special characteristic of “presupernova models” that lead some researchers in the past to miss the importance of beta decay during a transient stage an hour or so prior to collapse.
While dynamic weak equilibrium is achieved in the $`15\mathrm{M}_{}`$ model, with the new rates, it is not in the $`25\mathrm{M}_{}`$, though beta-decay still offers a non-negligible resistance to electron capture even there. This is in part due to the shorter time scale of silicon shell burning in the more massive star and also the larger value of $`Y_e`$ in the cores of stars with higher entropies (Fig. 1).
The presence of an important beta-decay contribution has two effects. Obviously it counteracts the reduction of $`Y_e`$ in the core, but equally important, beta decays are an additional neutrino source and thus add to the cooling of the core and a reduction in entropy . This cooling can be quite efficient, as often the average neutrino energy in the involved beta decays is larger than for the competing electron captures. As a consequence the new models have significant lower core temperatures than the WW models after silicon burning, which is particularly pronounced for the $`15\mathrm{M}_{}`$ star. During the contraction stage electron capture is again more important than beta decays, associated with the increased electron Fermi energy. Although the shell model rates are individually smaller than the FFN electron capture rates, the effective electron capture rate is larger in the new models as the evolution now proceeds along a trajectory with larger $`Y_e`$ values involving nuclei with smaller Q-values thus making electron capture energetically easier.
In summary, the shell model weak interaction rates result in significantly larger $`Y_e`$ values during the presupernova evolution of a 15 and 25 $`\mathrm{M}_{}`$ star than calculated by WW. Part - about half - of the change is due to including beta-decay and the other half is due to slower rates for electron capture. Fig. 2 shows that this is a general finding for all stars in the mass range $`11`$$`40\mathrm{M}_{}`$. The central values of $`Y_e`$ at the onset of core collapse are increased by $`0.01`$$`0.015`$. This is a significant effect. We note that the new models also result in lower core entropies for stars with $`M20\mathrm{M}_{}`$. The larger core $`Y_e`$ values and the lower entropies suggest that these stars will have a larger homologous core than currently assumed. For $`M20\mathrm{M}_{}`$, the new models actually have a slightly larger entropy. Thus one might expect that increased electron capture rate on (the more abundant) free protons will partly counteract the increase in $`Y_e`$ values during the subsequent core collapse phase. In general, core collapse calculations with detailed neutrino transport are required before definite conclusions about the explosion mechanism can be drawn. We will provide our models to those wishing to attempt such calculations.
Another important property for the core collapse is the size of the iron core, which we define as the mass interior to the point where the composition becomes at least $`50\%`$ of iron group elements $`(A48`$). As is shown in Fig. 2, the iron core masses are generally smaller in the new models. The effect is larger for stars more massive than $`20\mathrm{M}_{}`$, while for the most common supernovae ($`M20\mathrm{M}_{}`$) the reduction is about $`0.05\mathrm{M}_{}`$. This reduction appears to be counterintuitive at first glance with respect to the slower electron capture rates in the new models. It is, however, related to changes in the entropy profile during silicon shell burning which reduces the growth of the iron core just prior to collapse . Clearly, though we have, for reasons of space concentrated upon central values, it is the entire distribution of entropy and $`Y_e`$ in the stellar interior that determines its final evolution.
Confidence in the shell model rates stems from the recent measurements of GT distributions in iron mass nuclei which are all well reproduced by the shell model calculations . However, the energy resolutions of these pioneering (n,p) charge-exchange studies performed at TRIUMF has been rather limited ($`1`$$`1.5`$MeV) and they have been performed for stable nuclei only. These limitations are likely to be overcome in the near future as measurements with charge-exchange reactions like $`(`$d$`,^2`$He), $`(`$t$`,^3`$He) promise data with an order of magnitude improved resolution. Furthermore, after radioactive ion beam facilities will become operational it will be possible by inverse techniques also to determine the GT distributions for unstable nuclei. Of course, laboratory experiments cannot measure directly the relevant stellar rates as these involve, for example, electron capture or beta decays from excited states. Nevertheless, high-resolution charge-exchange or beta-decay experiments are important for two reasons: First, they are stringent constraints for the nuclear models and their predictive powers. Second, such experiments can determine the energy positions of the daughter states for the GT transitions which can then be used directly in the determination of the rates.
To guide such experiments we have attempted to identify the most important nuclei for electron-capture and beta decay during the final stages of stellar evolution. The relevant quantity is the product of abundance of the nuclear species in the core composition and electron capture rate (or beta decay rate). Table I lists the most important nuclei at selected points during the final evolution of our $`15\mathrm{M}_{}`$ and $`25\mathrm{M}_{}`$ stars. Since the $`Y_e`$ values are significantly larger in the new models, the important flows now involve nuclei much closer to stability so that several of the most important electron capturing nuclei (e.g., <sup>54,56,58</sup>Fe, <sup>55</sup>Mn, <sup>53</sup>Cr) are stable. Beta decay, on the other hand, mainly involves unstable manganese and cobalt isotopes. However, the lifetimes are long enough to allow for an experimental determination of the relevant GT strength distribution once radioactive ion beam facilities become operational. We finally mention that the backresonances contribute noticeably to the stellar beta decay rates for these isotopes making also measurement of the GT<sub>+</sub> strength on the daughter nuclei (e.g. iron and nickel isotopes) very useful.
We would like to thank G. M. Fuller for helpful discussions. This work has been partly supported by the Danish Research Council, the Carlsberg Foundation, by NATO grant CRPG973035, by the National Science Foundation (NSF-AST-9731569), by the US Department of Energy (DOE/LLNL B347885), and by the Alexander von Humboldt-Stiftung (FLF-1065004).
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# On the Observation of Phase Transitions in Collisions of Elementary Matter
## 1 Motivation
The primary goal for the investigation of heavy-ion collisions is to test the equation of state (EoS) of hot and dense matter far off the ground state, especially with view on possible phase transitions, e.g. to the Quark-Gluon-Plasma (QGP) . Indeed, collective flow phenomena are sensitive indicators for thermodynamically abnormal matter . In the case of a first-order phase transition to a QGP, an isentropic expansion proceeds through a stage of phase coexistence which should lead to signatures in the observables. The first ideas to investigate this phenomenon occured in the mid-seventies . In this paper we investigate the excitation function of directed flow, as well as the compactness in heavy-ion collisions.
## 2 Model
To investigate quantitatively the experimental observables, we perform 1-fluid and 3-fluid (3+1)-dimensional relativistic hydrodynamic calculations. That is, we solve numerically the continuity equations for the energy-momentum tensor, $`_\mu T^{\mu \nu }=0`$, and the net baryon current, $`_\mu N_B^\mu =0`$. Detailed discussions of (3+1)-d numerical solutions for hydrodynamical compression and expansion can be found e.g. in . We shall employ two different equations of state for $`P(e,\rho )`$:
* A relativistic mean field (RMF) hadron fluid corresponding to baryons and antibaryons interacting via exchange of massive scalar and vector bosons, plus free thermal pions; the parameters of the Lagrangian are fitted to the ground state of infinite nuclear matter, in particular the nuclear saturation density, the energy per particle, and the incompressibility.
* The same EoS as in i) for the low density phase, but supplemented by a Bag Model EoS with a bag constant $`B^{1/4}=235`$ MeV for the quark-gluon (QGP) phase. The phase coexistence region corresponding to this first-order transition is constructed employing the Gibbs’ condition of phase equilibrium, $`P_{HG}(T,\mu _B)=P_{QG}(T,\mu _B)`$, where $`T`$ and $`\mu _B`$ denote the temperature and the baryon-chemical potential, respectively. For example, for $`\rho =0`$ we find $`T_C170`$ MeV, while at $`T=0`$ phase coexistence sets in at $`\rho 4.6\rho _0`$. A more detailed discussion of these EoS can be found in .
Further, we employ the three-fluid model with a dynamical local unification procedure . The three-fluid model treats the nucleons of the projectile and target nuclei as two different fluids, since they populate different rapidity regions in the beginning of the reaction. The same holds for the newly produced particles around midrapidity, which are therefore collected in the third fluid. Thus, the three-fluid model accounts for the non-equilibrium situation during the compression stage of heavy-ion collisions. The coupling between the projectile and target fluids is calculated assuming free binary $`NN`$-collisions .
The unification of fluids $`i`$ and $`j`$ consists of adding their energy-momentum tensors and net-baryon currents in the respective cells,
$$T_i^{\mu \nu }(x)+T_j^{\mu \nu }(x)=T_{\mathrm{unified}}^{\mu \nu }(x),N_i^\mu (x)+N_j^\mu (x)=N_{\mathrm{unified}}^\mu (x)$$
(1)
and common values for $`e,P,\rho `$ and $`u^\mu `$ are obtained from $`T_{\mathrm{unified}}^{\mu \nu }=(e+P)u^\mu u^\nu Pg^{\mu \nu }`$, $`N_{\mathrm{unified}}^\mu =\rho u^\mu `$, and the given EoS $`P=P(e,\rho )`$. The local criterion for unification is $`(P_i+P_j)/P>90\%`$. Here, $`P_{i,j}`$ denotes the pressure in $`T_{i,j}^{\mu \nu }`$, and $`P`$ the pressure in $`T_{\mathrm{unified}}^{\mu \nu }`$.
## 3 Directed flow and softest point of the EoS
In order to measure the EoS, i.e. the pressure $`P(e,\rho )`$ as a function of energy density $`e`$ and baryon density $`\rho `$ in the local rest frame of a fluid element, the transverse momentum in the reaction plane, $`p_x`$, is investigated. This quantity is proportional to the pressure created in the hot and dense collision zone :
$$p_xPdA_{}dt.$$
(2)
The pressure $`P`$ is exerted on a transverse area element $`A_{}`$. Directed flow has therefore been proposed as a measure for the pressure and a possible “softening” of the EoS .
Fig. 2 shows the excitation function of directed flow $`p_x^{\mathrm{dir}}/N`$ calculated in the three-fluid model in comparison to that obtained in a one-fluid calculation . The one fluid calculations show that for increasing bombarding energy, the flow, $`p_x`$, first increases (for a collision without phase transition), as the compression and thus the pressure grow. At large $`E_{\mathrm{Lab}}^{\mathrm{kin}}`$ the time span of the collision decreases, diminishing the flow again. The flow is thus maximized at some intermediate bombarding energy. In the case with a phase transition the decrease of the flow is much more rapid than for the purely hadronic fluid. The reason for this is not that $`c_s`$, the isentropic velocity of sound, vanishes but rather geometry: The compactness and tilt-angle $`\mathrm{\Theta }`$ are different in the calculation with phase transition, and this leads to a different initial condition for the subsequent expansion (see below and ). After passing through a local minimum at $`E_{\mathrm{Lab}}^{\mathrm{kin}}5A`$ GeV, the directed in-plane momentum reaches a second local maximum around $`E_{\mathrm{Lab}}^{\mathrm{kin}}1020A`$ GeV. This is the point where the compressed matter first becomes hot enough (over a large volume) to “respond” with small pressure gradients.
Due to non-equilibrium effects in the early stage of the reaction, which delay the build-up of transverse pressure , the flow in the three-fluid model is reduced as compared to the one-fluid calculation in the AGS energy range. Furthermore, the minimum in the excitation function of the directed flow shifts to higher energies. The case without dynamical unification yields the least amount of stopping and energy deposition, while the one-fluid calculation has instantaneous full stopping and maximum energy deposition. The three-fluid model with dynamical unification lies between these two limits; it accounts for the limited stopping power of nuclear matter in the early stages of the collision and mutual equilibration of the different fluids in the later stages. Most importantly, the three-fluid calculations predict an increase of $`p_x^{\mathrm{dir}}/N`$ towards $`E_{\mathrm{Lab}}^{\mathrm{kin}}40A`$ GeV, if indeed a phase coexistence with small $`c_s`$ occurs. Data at that energy has recently been taken, and should prove very useful to pin down the onset (or absence) of a first-order phase transition in the AGS-SPS energy domain.
First order phase transitions “soften” the EoS , i.e. $`P`$ increases slower with $`e`$ and $`\rho `$ than in the case without phase transition. This corresponds to a reduction of the isentropic speed of sound, $`c_s`$, as compared to that in the interacting hadron fluid. However, as shown in Fig. 2, this happens only if the entropy per net baryon in the central region is large, i.e. if the ratio $`T/\mu _B`$ is not too small . In the three-fluid model that is due mainly to the kinetic equilibration of the decelerating baryon dense projectile/target fluids with the midrapidity fluid of secondary particles, which leads to considerably larger $`s/\rho `$ than in one-fluid hydrodynamics. However, at the energy corresponding to the minimum of the directed flow the specific entropy is rather small, $`s/\rho 10`$; i.e. in the AGS energy domain matter is rather baryon dense but not very hot. Consequently, the EoS is not soft (i.e. the isentropic velocity of sound is not small), even if mixed phase matter does occur.
Fig. 3 shows the time-like component of the net baryon current in momentum space ($`p_xp_{long}`$ plane) for Pb+Pb-collisions at $`b=3`$ fm. One clearly observes the directed in-plane flow (before the collision, there is no matter at $`p_x0`$). However, in the left panel there is almost no momentum of baryons in the upper left or bottom right quadrants, where $`p_xp_{long}<0`$ (except for two “jets”, see below). The central region passed through the phase coexistence region at high $`s/\rho `$, with a rather small average $`c_s`$, and isentropic expansion of the highly excited matter is inhibited. Note the difference to the expansion pattern observed for lower energies, right panel of Fig. 3, where $`c_s`$ is not small .
The slope of $`p_x/N`$ at midrapidity is shown in Fig. 4 as a function of beam energy. Experimental Data are shown as well . One observes a steady decrease of $`F_y=\mathrm{d}(p_x/N)/\mathrm{d}y`$ up to about top BNL-AGS energy, where the flow around midrapidity even becomes negative due to preferred expansion towards $`p_xp_{long}<0`$. The “overshoot” towards negative slope is due to the small incompressibility in the top AGS energy region, and the rather early fluid unification employed here. However, such a behavior can not be observed in the data. A less steep decrease of $`F_y`$ could be achieved in the three-fluid model by a more stringent unification criterion (i.e. later unification) or early kinetic decoupling on the hadronization hypersurface. At higher energy, $`E_{Lab}40A`$ GeV, we encounter the expansion pattern depicted in the left panel of Fig. 3: flow towards $`p_xp_{long}<0`$ can not build up! Consequently, $`F_y`$ increases rapidly towards $`E_{Lab}=2040A`$ GeV, decreasing again at even higher energy because of the more forward-backward peaked kinematics. Note that the increase of the slope is due to the absence of the “anti-flow”, see Fig. 3. In any case, Fig. 4 shows that it will be difficult to see the effect of the possible phase transition in $`F_y`$. The double-differential in-plane cross section, Fig. 3, appears more useful.
## 4 Compactness
Measurement of the compactness is a promising new tool to observe the onset of the phase transition. It relies on measuring the shape of the source, which is uniquely related to the pressure and density of the system in the compression and expansion stage of the nucleus-nucleus collision. The compactness can be identified via interferometry: The illuminiation of the baryon source by the pion radiation is subject to experimental scrutiny via pion interferometry . Fig. 5 illustrates the basic idea. It shows the baryon density in the reaction plane for the EoS without (i) and with (ii) phase transition, respectively. One clearly observes the higher compression in the case with phase transition. As indicated above, the onset of the transition to quark matter at a given incident energy $`E_{lab}^{kin}`$ leads to higher compression $`\rho /\rho _0`$ than for the case without transition. Now, as $`\rho V\pi R_A^2L\rho `$ must equal $`2A`$ by virtue of baryon number conservation, the longitudinal thickness $`L`$ of the compressed matter is approximately proportional to $`1/\rho `$. Thus, a transition to quark matter leads to a more compact system, just as quark matter stars are more compact than pure neutron stars . Of course, in heavy-ion collisions that expectation is based on the behavior of relativistic compression shocks rather than hydrostatic and gravitational equilibrium.
In particular, we study the compactness in the reactions Au+Au at impact parameter $`b=3`$ fm for the energy $`E_{lab}^{kin}=8A`$ GeV. The compactness is defined as the ratio of the smallest to the biggest in-plane eigenvalue of the configuration space sphericity tensor, which we define as the second moment of the net baryon current. On fixed-time hypersurfaces we have
$$F_{ij}=\mathrm{d}^3xx_ix_jN_B^0\mathrm{\Theta }\left(\rho (\stackrel{}{x})\rho _{cut}\right).$$
(3)
We apply an additional density cut $`\rho >\rho _{cut}`$ in the integral to discard spectator matter. In the future the cuts and the hypersurface will have to be adapted to the experimental conditions. However, this is not crucial for understanding the effect.
The three eigenvalues $`f_n`$ are the solutions of the cubic equation $`\mathrm{det}(F_{ij}f\delta _{ij})=0`$, and the eigenvectors $`\stackrel{}{e}_n`$ follow from solving the linear systems of equations $`(F_{ij}f_n\delta _{ij})e_n^j=0`$. In terms of the eigenvalues $`f_n`$ and orthogonal eigenvectors $`\stackrel{}{e}_n`$ of $`F`$, the compactness tensor can be written as $`F=f_1\stackrel{}{e}_1\stackrel{}{e}_1+f_2\stackrel{}{e}_2\stackrel{}{e}_2+f_3\stackrel{}{e}_3\stackrel{}{e}_3`$. In diagonal form, $`F`$ specifies an ellipsoid in configuration space with principal axis along $`\stackrel{}{e}_n`$ and radii $`\sqrt{f_n}`$. Cigar-like patterns, oriented along the z-axis, would lead to $`f_1>f_2=f_3`$, $`\stackrel{}{e}_1=\stackrel{}{e}_z`$, $`\stackrel{}{e}_2=\stackrel{}{e}_x`$, $`\stackrel{}{e}_3=\stackrel{}{e}_y`$. On the other hand, a “pancake”/“lensil” shape corresponds to $`f_1<f_2=f_3`$. The tilt angle $`\mathrm{\Theta }`$ is determined from the scalar product of $`\stackrel{}{e}_z`$ (the longitudinal direction in the lab frame) with the vector $`\stackrel{}{e}_n`$ corresponding to the biggest of the eigenvalues $`f_n`$. As already indicated in the introduction, we find very different eigenvalues for the two equations of state. The calculation with transition to quark matter corresponds to higher compactness of the baryon distribution. That is, the compactness tensor is much flatter (nearly a factor of two !) in the model with (ii) than in the model without (i) phase transition. Moreover, our (3+1)-dimensional expansion solutions show that after the compression stage the ratio of the in-plane radii $`\sqrt{f_2/f_1}`$ remains much smaller in the case with phase transition, cf. Fig. 6. Note also that the extremely flat shape of the baryon distribution means very small curvature of the surface, which in turn will lead to a strongly “bundled” emission of hadrons from the (almost planar) rarefaction wave or deflagration shock by which the dense droplet decays; see also the discussion in .
The eigenvalues of the compactness tensor allow to measure directly the density increase in the high density stage of the reaction, if a phase transition occurs. Care must be taken that the impact parameter range investigated constitutes a moderately small bin of centrality values. One should keep in mind that the compression factor is affected by the incompressibility $`P/e`$ evaluated on the shock adiabat (in the one-fluid model), not along a path of fixed specific entropy. Therefore, the incompressibility needs not be equal to the isentropic speed of sound. The latter is not much reduced in the presence of the phase transition at AGS energies , because of the high net baryon density and relatively low temperature.
Lisa et al. have recently proposed a new interferometry analysis, which could be used to observe this change in the compactness directly. The developed method is quite robust and incorporates other interesting information as the configuration space tilt angle, which nicely complement the momentum-space flow angles. It avoids cuts in tilted ellipsoids, which are not analysed in the appropriate rotated frame, and where the excentricity and the RMS-radii are much less distinct for the two different equations of state.
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# Explicit Formulas for the Multivariate Resultant
## 1. Introduction
Given $`n`$ homogeneous polynomials $`f_1,\mathrm{},f_n`$ in $`n`$ variables over an algebraically closed field $`k`$ with respective degrees $`d_1,\mathrm{},d_n`$, the resultant $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ is an irreducible polynomial in the coefficients of $`f_1,\mathrm{},f_n`$, which vanishes whenever $`f_1,\mathrm{},f_n`$ have a common root in projective space. The study of resultants goes back to classical work of Sylvester, Bézout, Cayley, Macaulay and Dixon. The use of resultants as a computational tool for elimination of variables as well as a tool for the study of complexity aspects of polynomial system solving in the last decade, has renewed the interest in finding explicit formulas for their computation (cf. , , , ,, ,, ,).
By a determinantal formula it is meant a matrix whose entries are polynomials in the coefficients of $`f_1,\mathrm{},f_n`$ and whose determinant equals the resultant $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$. Of course, the interest on such a formula is the computation of the resultant, and so it is implicit that the entries should be algorithmically computed from the inputs. It is also meant that all non-zero entries have degree strictly less than the degree of the resultant.
In case all $`d_i`$ have a common value $`d`$, all currently known determinantal formulas are listed by Weyman and Zelevinsky in . This list is short: if $`d2,`$ there exist determinantal formulas for all $`d`$ just for binary forms (given by the well known Sylvester matrix), ternary forms and quaternary forms; when $`n=5`$, the only possible values for $`d`$ are $`2`$ and $`3`$; finally, for $`n=6`$, there exists a determinantal formula only for $`d=2`$. We find similar strict restrictions on general $`n,d_1,\mathrm{},d_n`$ (cf. Lemma 5.3).
Given $`d_1,\mathrm{},d_n,`$ denote $`t_n:=_{i=1}^n(d_i1)`$ the critical degree. Classical Macaulay formulas describe the resultant $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ as an explicit quotient of two determinants. These formulas involve a matrix of size at least the number of monomials in $`n`$ variables of degree $`t_n+1`$, and a submatrix of it.
Macaulay’s work has been revisited and sharpened by Jouanolou in , where he proposes for each $`t0,`$ a square matrix $`M_t`$ of size
(1)
$$\rho \left(t\right):=\left(\genfrac{}{}{0pt}{}{t+n1}{n1}\right)+i(t_nt)$$
whose determinant is a nontrivial multiple of $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ (cf. , $`\mathrm{3.11.19.7}`$). Here, $`i(t_nt)`$ denotes the dimension of the $`k`$vector space of elements of degree $`t_nt`$ in the ideal generated by a regular sequence of $`n`$ polynomials with degrees $`d_1,\mathrm{},d_n.`$ Moreover, Jouanolou shows that the resultant may be computed as the ratio between the determinant of $`M_{t_n}`$ and the determinant of one of its square submatrices. (cf. , Corollaire $`\mathrm{3.9.7.7}`$).
In this paper, we explicitly find the extraneous factor in Jouanolou’s formulation, i.e. the polynomial $`det(M_t)/\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$, for all $`t0`$ which again happens to be the determinant of a submatrix $`𝔼_t`$ of $`M_t`$ for every $`t`$, and this allows us to present new resultant formulas à la Macaulay for the resultant, i.e. as a quotient of two determinants
(2)
$$\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)=\frac{det(M_t)}{det(𝔼_t)}.$$
For $`t>t_n,`$ we recover Macaulay’s classical formulas. For $`tt_n,`$ the size of the matrix $`M_t`$ is considerably smaller.
In order to give explicit examples, we need to recall the definition of the Bezoutian associated with $`f_1,\mathrm{},f_n`$ (cf. , , , and under the name “Formes de Morley”). Let $`(f_1,\mathrm{},f_n)`$ be a sequence of generic homogeneous polynomials with respective degrees $`d_1,\mathrm{},d_n`$
$$f_i:=\underset{|\alpha _i|=d_i}{}a_{\alpha _i}X^{\alpha _i}A[X_1,\mathrm{},X_n],$$
where $`A`$ is the factorial domain $`A:=\left[a_{\alpha _i}\right]_{|\alpha _i|=d_i,i=1,\mathrm{},n}.`$
Introduce two sets of $`n`$ variables $`X,Y`$ and for each pair $`(i,j)`$ with $`1i,jn`$, write $`\mathrm{\Delta }_{ij}(X,Y)`$ for the incremental quotient
(3)
$$\frac{f_i(Y_1,\mathrm{},Y_{j1},X_j,\mathrm{},X_n)f_i(Y_1,\mathrm{},Y_j,X_{j+1},\mathrm{},X_n)}{X_jY_j}.$$
Note that $`f_i(X)f_i(Y)=_{j=1}^n\mathrm{\Delta }_{ij}(X,Y)(X_jY_j).`$
The determinant
(4)
$$\mathrm{\Delta }(X,Y):=det(\mathrm{\Delta }_{ij}(X,Y))_{1i,jn}=\underset{\left|\gamma \right|t_n}{}\mathrm{\Delta }_\gamma \left(X\right).Y^\gamma .$$
is a representative of the Bezoutian associated with $`(f_1,\mathrm{},f_n).`$ It is a homogeneous polynomial in $`A[X,Y]`$ of degree $`t_n.`$
Recall also that
$$\mathrm{deg}\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)=\underset{i=1}{\overset{n}{}}d_1\mathrm{}d_{i1}d_i\mathrm{}d_n.$$
As a first example, let $`n=3,(d_1,d_2,d_3)=(1,1,2),`$ and let
$$\begin{array}{ccc}f_1& =& a_1X_1+a_2X_2+a_3X_3\\ f_2& =& b_1X_1+b_2X_2+b_3X_3\\ f_3& =& c_1X_1^2+c_2X_2^2+c_3X_3^2+c_4X_1X_2+c_5X_1X_3+c_6X_2X_3\end{array}$$
be generic polynomials of respective degrees $`1,1,2`$. Here, $`t_3=1.`$ Macaulay’s classical matrix $`M_2`$ looks as follows:
$$\left(\begin{array}{cccccc}a_1& 0& 0& 0& 0& c_1\\ 0& a_2& 0& b_2& 0& c_2\\ 0& 0& a_3& 0& b_3& c_3\\ a_2& a_1& 0& b_1& 0& c_4\\ a_3& 0& a_1& 0& b_1& c_5\\ 0& a_3& a_2& b_3& b_2& c_6\end{array}\right)$$
and its determinant equals $`a_1Res_{1,1,2}.`$ The extraneous factor is the $`1\times 1`$ minor formed by the element in the fourth row, second column.
On the other hand, because of Lemma 5.3, we can exhibit a determinantal formula for $`\pm Res_{1,1,2},`$ and it is given by Proposition 5.6 for $`t=\left[\frac{t_3}{2}\right]=0`$ by the determinant of
$$\left(\begin{array}{ccc}\mathrm{\Delta }_{(1,0,0)}& a_1& b_1\\ \mathrm{\Delta }_{(0,1,0)}& a_2& b_2\\ \mathrm{\Delta }_{(0,0,1)}& a_3& b_3\end{array}\right),$$
where $`\mathrm{\Delta }_\gamma `$ are coefficients of the Bezoutian (4). Explicitly, we have
$$\mathrm{\Delta }_{(1,0,0)}=c_1(a_2b_3a_3b_2)c_4(a_1b_3a_3b_1)+c_5(a_1b_2a_2b_1),$$
$$\mathrm{\Delta }_{(0,1,0)}=c_6(a_1b_2a_2b_1)c_2(a_1b_3b_1a_3)$$
and
$$\mathrm{\Delta }_{(0,0,1)}=c_3(a_1b_2b_1a_2).$$
This is the matrix $`M_0`$ corresponding to the linear transformation $`\mathrm{\Psi }_0`$ which is defined in (9).
Take now $`n=4,`$ and $`(d_1,d_2,d_3,d_4)=(1,1,2,3).`$ The critical degree is $`3.`$ Macaulay’s classical matrix, $`M_4,`$ has size $`35\times 35.`$ Because the degree of $`\mathrm{Res}_{1,1,2,3}`$ is $`2+3+6+6=17,`$ we know that its extraneous factor must be a minor of size $`18\times 18.`$ By Proposition 5.6, we can find the smallest possible matrix for $`t=1`$ or $`t=2.`$ Set $`t=2.`$ We get the following $`12\times 12`$ matrix
$$\left(\begin{array}{cccccccccccc}\mathrm{\Delta }_{(2,0,0,0)}^1& \mathrm{\Delta }_{(2,0,0,0)}^2& \mathrm{\Delta }_{(2,0,0,0)}^3& \mathrm{\Delta }_{(2,0,0,0)}^4& a_1& 0& 0& 0& 0& 0& 0& c_1\\ \mathrm{\Delta }_{(0,2,0,0)}^1& \mathrm{\Delta }_{(0,2,0,0)}^2& \mathrm{\Delta }_{(0,2,0,0)}^3& \mathrm{\Delta }_{(0,2,0,0)}^4& 0& a_2& 0& 0& b_2& 0& 0& c_2\\ \mathrm{\Delta }_{(0,0,2,0)}^1& \mathrm{\Delta }_{(0,0,2,0)}^2& \mathrm{\Delta }_{(0,0,2,0)}^3& \mathrm{\Delta }_{(0,0,2,0)}^4& 0& 0& a_3& 0& 0& b_3& 0& c_3\\ \mathrm{\Delta }_{(0,0,0,2)}^1& \mathrm{\Delta }_{(0,0,0,2)}^2& \mathrm{\Delta }_{(0,0,0,2)}^3& \mathrm{\Delta }_{(0,0,0,2)}^4& 0& 0& 0& a_4& 0& 0& b_4& c_4\\ \mathrm{\Delta }_{(1,1,0,0)}^1& \mathrm{\Delta }_{(1,1,0,0)}^2& \mathrm{\Delta }_{(1,1,0,0)}^3& \mathrm{\Delta }_{(1,1,0,0)}^4& a_2& a_1& 0& 0& b_1& 0& 0& c_5\\ \mathrm{\Delta }_{(1,0,1,0)}^1& \mathrm{\Delta }_{(1,0,1,0)}^2& \mathrm{\Delta }_{(1,0,1,0)}^3& \mathrm{\Delta }_{(1,0,1,0)}^4& a_3& 0& a_1& 0& 0& b_1& 0& c_6\\ \mathrm{\Delta }_{(1,0,0,1)}^1& \mathrm{\Delta }_{(1,0,0,1)}^2& \mathrm{\Delta }_{(1,0,0,1)}^3& \mathrm{\Delta }_{(1,0,0,1)}^4& a_4& 0& 0& a_1& 0& 0& b_1& c_7\\ \mathrm{\Delta }_{(0,1,1,0)}^1& \mathrm{\Delta }_{(0,1,1,0)}^2& \mathrm{\Delta }_{(0,1,1,0)}^3& \mathrm{\Delta }_{(0,1,1,0)}^4& 0& a_3& a_2& 0& b_3& b_2& 0& c_8\\ \mathrm{\Delta }_{(0,1,0,1)}^1& \mathrm{\Delta }_{(0,1,0,1)}^2& \mathrm{\Delta }_{(0,1,0,1)}^3& \mathrm{\Delta }_{(0,1,0,1)}^4& 0& a_4& 0& a_2& b_4& 0& b_2& c_9\\ \mathrm{\Delta }_{(0,0,1,1)}^1& \mathrm{\Delta }_{(0,0,1,1)}^2& \mathrm{\Delta }_{(0,0,1,1)}^3& \mathrm{\Delta }_{(0,0,1,1)}^4& 0& 0& a_4& a_3& 0& b_4& b_3& c_{10}\\ a_1& a_2& a_3& a_4& 0& 0& 0& 0& 0& 0& 0& 0\\ b_1& b_2& b_3& b_4& 0& 0& 0& 0& 0& 0& 0& 0\end{array}\right)$$
where
$$\begin{array}{ccc}f_1& =& a_1X_1+a_2X_2+a_3X_3+a_4X_4\\ f_2& =& b_1X_1+b_2X_2+b_3X_3+b_4X_4\\ f_3& =& c_1X_1^2+c_2X_2^2+c_3X_3^2+c_4X_4^2+c_5X_1X_2+c_6X_1X_3\\ & & +c_7X_1X_4+c_8X_2X_3+c_9X_2X_4+c_{10}X_3X_4,\end{array}$$
$`f_4`$ is a homogeneous generic polynomial of degree $`3`$ in four variables, and for each $`\gamma ,|\gamma |=2,`$ we write
$$\mathrm{\Delta }_\gamma (X)=\underset{j=1}{\overset{4}{}}\mathrm{\Delta }_\gamma ^jX_j,$$
which has degree $`1`$ in the coefficients of each $`f_i,i=1,\mathrm{},4.`$ The determinant of this matrix is actually $`\pm a_1\mathrm{Res}_{1,1,2,3}.`$ Here, the extraneous factor is the minor $`1\times 1`$ of the matrix obtained by taking the element in the fifth row, sixth column.
In the following table, we display the minimal size of the matrices $`M_t`$ and the size of classical Macaulay matrix for several values of $`n,d_1,\mathrm{},d_n.`$
$$\overline{)\begin{array}{cccccccc}n& & (d_1,\mathrm{},d_n)& & \text{ min size }& & \text{ classical }& \\ & & & & & & & \\ 2& & (10,70)& & 70& & 80& \\ 2& & (150,200)& & 200& & 350& \\ 3& & (1,1,2)& & 3& & 6& \\ 3& & (1,2,5)& & 14& & 28& \\ 3& & (2,2,6)& & 21& & 45& \\ 4& & (1,1,2,3))& & 12& & 35& \\ 4& & (2,2,5,5)& & 94& & 364& \\ 4& & (2,3,4,5)& & 90& & 364& \\ 5& & (4,4,4,4,4)& & 670& & 4845& \\ 7& & (2,3,3,3,3,3,3)& & 2373& & 38760& \\ 10& & (3,3,\mathrm{},3)& & 175803& & 14307150& \\ 20& & (2,2,\mathrm{},2)& & 39875264& & 131282408400& \end{array}}.$$
We give in section 4 an estimate for the ratio between these sizes. However, it should be noted that the number of coefficients of the Bezoutian that one needs to compute increases when the size of the matrix $`M_t`$ decreases. We refer to and for complexity considerations on the computation of Bezoutians. In particular, this computation can be well parallelized. Also, the particular structure of the matrix and the coefficients could be used to improve the complexity estimates; this problem is studied for $`n=2`$ and $`n=3`$ in .
Our approach combines Macaulay’s original ideas , expanded by Jouanolou in , with the expression for the resultant as the determinant of a Koszul complex inspired by the work of Cayley . We also use the work , of Chardin on homogeneous subresultants, where a Macaulay style formula for subresultants is presented. In fact, we show that the proposed determinants are explicit non-zero minors of a bigger matrix which corresponds to one of the morphisms in a Koszul resultant complex which in general has many non zero terms, and whose determinant is $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ (cf. Theorem 5.1). These are the complexes considered in , in the equal degree case, built from the spectral sequence associated with a twisted Koszul complex at the level of sheaves.
We give explicit expressions for the morphisms in these complexes in terms of the Bezoutian associated with $`f_1,\mathrm{},f_n`$ for degrees under critical degree, addressing in this manner a problem raised by Weyman and Zelevinsky in (cf. also \[15, 13.1.C\].
In the last sections, we show that different classical formulas can be viewed as special cases of the determinantal formulas that we present here (cf. ,). In particular, we also recover in this setting the “affine” Dixon formulas considered in and we classify in particular all such determinantal formulas.
## 2. Notations and some preliminary statements
Let $`S_u`$ denote the $`A`$-free module generated by the monomials in $`A[X]`$ with degree $`u.`$ If $`u<0,`$ then we set $`S_u=0.`$ Define also the following free submodules $`E^{t,j}S^{t,j}S_{td_j},`$ for all $`j=1,\mathrm{},n:`$
(5)
$$S^{t,j}:=X^\gamma ,|\gamma |=td_j,\gamma _1<d_1,\mathrm{},\gamma _{j1}<d_{j1}$$
(6)
$$E^{t,j}:=X^\gamma S^{t,j},\text{there exists }ij:\gamma _id_i.$$
Note that $`E^{t,n}=0,`$ and $`S^{t,1}=S_{td_1}t_0.`$
Let $`j_u:S_uS_u^{}`$ be the isomorphism associated with the monomial bases in $`S_u`$ and denote by $`T_\gamma :=j_u(X^\gamma )`$ the elements in the dual basis.
Convention. All spaces that we will consider have a monomial basis, or a dual monomial basis. We shall suppose all these bases have a fixed order. This will allow us to define matrices “in the monomial bases”, with no ambiguity.
Let $`\psi _{1,t}`$ be the $`A`$-linear map
$$\psi _{1,t}:S_{t_nt}^{}S_t$$
which sends
(7)
$$T_\gamma \mathrm{\Delta }_\gamma \left(X\right),$$
where the polynomial $`\mathrm{\Delta }_\gamma \left(X\right)`$ is defined in (4). Let $`\mathrm{\Delta }_t`$ denote the matrix of $`\psi _{1,t}`$ in the monomial bases.
###### Lemma 2.1.
For suitable orders of the monomial bases in $`S_t`$ and $`S_{t_nt},`$ we have that
$${}_{}{}^{𝐭}\mathrm{\Delta }_{t}^{}=\mathrm{\Delta }_{t_nt}.$$
###### Proof.
It holds that $`\mathrm{\Delta }(X,Y)=\mathrm{\Delta }(Y,X)`$ by the symmetry property of Bezoutians (cf. \[17, 3.11.8\]). This implies that
$$\underset{|\gamma |=t_nt}{}\mathrm{\Delta }_\gamma \left(X\right)Y^\gamma =\underset{|\lambda |=t}{}\mathrm{\Delta }_\lambda \left(Y\right)X^\lambda =\underset{|\gamma |=t_nt,|\lambda |=t}{}c_{\gamma \lambda }X^\lambda Y^\gamma ,$$
with $`c_{\gamma \lambda }A.`$ It is easy to see that if $`\mathrm{\Delta }_t=\left(c_{\gamma \lambda }\right)_{|\gamma |=t_nt,|\lambda |=t}`$ then $`\mathrm{\Delta }_{t_nt}=\left(c_{\gamma \lambda }\right)_{|\lambda |=t,|\gamma |=t_nt}.`$
Let us consider also the Sylvester linear map $`\psi _{2,t}:`$
(8)
$$\begin{array}{cccccccc}\psi _{2,t}:& S^{t,1}& & \mathrm{}& & S^{t,n}& & S_t\\ & (g_1& ,& \mathrm{}& ,& g_n)& & _{i=1}^ng_if_i,\end{array}$$
and denote by $`D_t`$ its matrix in the monomial bases. As usual, $`\psi _{2,t_nt}^{}`$ denotes the dual mapping of (8) in degree $`t_nt.`$
Denote
(9)
$$\mathrm{\Psi }_t:S_{t_nt}^{}\left(S^{t,1}\mathrm{}S^{t,n}\right)S_t\left(S^{t_nt,1}\mathrm{}S^{t_nt,n}\right)^{}$$
the $`A`$-morphism defined by
(10)
$$(T,g)(\psi _{1,t}(T)+\psi _{2,t}(g),\psi _{2,t_nt}^{}(T)),$$
and call $`M_t`$ the matrix of $`\mathrm{\Psi }_t`$ in the monomial bases.
Denote also by $`E_t`$ the submatrix of $`M_t`$ whose columns are indexed by the monomials in $`E^{t,1}\mathrm{}E^{t,n1},`$ and whose rows are indexed by the monomials $`X^\gamma `$ in $`S_t`$ for which there exist two different indices $`i,j`$ such that $`\gamma _id_i,\gamma _jd_j.`$ With these choices it is not difficult to see that $`M_t`$ and $`E_t`$ (when defined) are square matrices.
###### Remark 2.2.
Observe that $`E_t`$ is actually a submatrix of $`D_t.`$ In fact, $`E_t`$ is transposed of the square submatrix named $`\left(t\right)`$ in , and whose determinant is denoted by $`\mathrm{\Delta }(n,t)`$ in \[21, Th. 6\].
###### Lemma 2.3.
$`M_t`$ is a square matrix of size $`\rho (t),`$ where $`\rho `$ is the function defined in (1).
###### Proof.
The assignment which sends a monomial $`m`$ in $`S^{t,i}`$ to $`x_i^{d_i}m`$ injects the union of the monomial bases in each $`S^{t,i}`$ onto the monomials of degree $`t`$ which are divisible by some $`x_i^{d_i}.`$ It is easy to see that the cardinality of the set of complementary monomials of degree $`t`$ is precisely $`H_d(t)`$, where $`H_d(t)`$ denotes the dimension of the $`t`$-graded piece of the quotient of the polynomial ring over $`k`$ by the ideal generated by a regular sequence of homogeneous polynomials with degrees $`d_1,\mathrm{},d_n`$ (cf. \[17, 3.9.2\]). Moreover, using the assignment $`(\gamma _1,\mathrm{},\gamma _n)(d_11\gamma _1,\mathrm{},d_n1\gamma _n),`$ it follows that
(11)
$$H_d(t)=H_d(t_nt).$$
We can compute explicitly this Hilbert function by the following formula (cf. \[21, §2\]):
(12)
$$\frac{_{i=1}^n\left(1Y^{d_i}\right)}{\left(1Y\right)^n}=\underset{t=0}{\overset{\mathrm{}}{}}H_d(t).Y^t.$$
Then,
$$\mathrm{rk}(S^{t,1}\mathrm{}S^{t,n})=\mathrm{rk}S_tH_d(t).$$
Similarly,
$$\mathrm{rk}\left(S^{t_nt,1}\mathrm{}S^{t_nt,n}\right)^{}=\mathrm{rk}(S_{t_nt})^{}H_d(t_nt).$$
Therefore, $`M_t`$ is square of size $`\mathrm{rk}S_tH_d(t_nt)+\mathrm{rk}S_{t_nt}.`$ Since $`i(t_nt)=\mathrm{rk}S_{t_nt}H_d(t_nt),`$ the size of $`M_t`$ equals $`\mathrm{rk}S_t+i(t_nt)=\rho (t).`$
###### Remark 2.4.
Ordering properly the monomial bases, $`M_t`$ is the transpose of the matrix which appears in \[17, 3.11.19.7\]. It has the following structure:
(13)
$$\left[\begin{array}{cc}\mathrm{\Delta }_t& D_t\\ {}_{}{}^{𝐭}D_{t_nt}^{}& 0\end{array}\right].$$
###### Remark 2.5.
Because $`\psi _{2,t}=0`$ if and only if $`t<\mathrm{min}\{d_i\},`$ we have that $`\mathrm{\Psi }_t=\psi _{2,t}+\psi _{1,t}`$ if $`t>t_n\mathrm{min}\{d_i\},`$ and $`\mathrm{\Psi }_t=\psi _{2,t}`$ if $`t>t_n.`$
Finally, denote $`𝔼_t`$ the square submatrix of $`M_t`$ which has the following structure:
(14)
$$𝔼_t=\left[\begin{array}{cc}& E_t\\ {}_{}{}^{𝐭}E_{t_nt}^{}& 0\end{array}\right].$$
It is clear from the definition that $`det(𝔼_t)=\pm det(E_t)det(E_{t_nt}).`$
###### Remark 2.6.
Dualizing (10) and using lemma 2.1 with a careful inspection at (13) and (14), we have that ordering properly their rows and columns,
$${}_{}{}^{𝐭}M_{t}^{}=M_{t_nt}\text{and}^𝐭𝔼_t=𝔼_{t_nt}.$$
## 3. Generalized Macaulay formulas
We can extend the map $`\psi _{2,t}`$ in (8) to the direct sum of all homogeneous polynomials with degrees $`td_1,\mathrm{},td_n,`$ and the map $`\psi _{2,t_nt}`$ to the direct sum of all homogeneous polynomials with degrees $`t_ntd_1,\mathrm{},t_ntd_n,`$ to get a map
$$\stackrel{~}{\mathrm{\Psi }}_t:\left(S_{t_nt}\right)^{}\left(S_{td_1}\mathrm{}S_{td_n}\right)S_t\left(S_{t_ntd_1}\mathrm{}S_{t_ntd_n}\right)^{}.$$
We can thus see the matrix $`M_t`$ of $`\mathrm{\Psi }_t`$ in (9) as a choice of a square submatrix of $`\stackrel{~}{\mathrm{\Psi }}_t.`$ We will show that its determinant is a non zero minor of maximal size.
###### Proposition 3.1.
Let $`M_t^{}`$ be a square matrix over $`A`$ of the form
(15)
$$M_t^{}:=\left[\begin{array}{cc}\mathrm{\Delta }_t& F_t\\ {}_{}{}^{𝐭}F_{t_nt}^{}& 0\end{array}\right]$$
where $`F_t`$ has $`i(t)`$ columns and corresponds to a restriction of the map
$$\begin{array}{ccc}S_{td_1}\mathrm{}S_{td_n}& & S_t\\ (g_1,\mathrm{},g_n)& & _{i=1}^ng_if_i;\end{array}$$
and similarly for $`F_{t_nt}`$ in degree $`t_nt.`$ Then, $`det(M_t^{})`$ is a multiple of $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ (probably zero).
###### Proof.
It is enough to mimic for the matrix $`M_t^{}`$ the proof performed by Jouanolou in \[17, Prop. 3.11.19.10\] to show that the determinant of the matrix $`M_t^{}`$ is an inertia form of the ideal $`f_1,\mathrm{},f_n`$ (i.e. a multiple of the resultant). We include this proof for the convenience of the reader.
Let $`N:=_{i=1}^n\mathrm{\#}\{\alpha _i^n:|\alpha _i|=d_i\}.`$ Given an algebraically closed field $`k,`$ and $`a=(a_{\alpha _i})_{|\alpha _i|=d_i,i=1,\mathrm{},n},`$ a point in $`k^N,`$ we denote by $`f_1(a),\mathrm{},f_n(a)k[X]`$ the polynomials obtained from $`f_1,\mathrm{},f_n`$ when the coefficients are specialized to $`a`$, and similarly for the coefficients of the Bezoutian. Because of the irreducibility of $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n),`$ it is enough to show that for all $`ak^N`$ such that $`f_1(a),\mathrm{},f_n(a)`$ have a non trivial solution in $`k^n,`$ the determinant of the specialized matrix $`M_t^{}(a)`$ is equal to $`0.`$
Suppose that this is case, and let $`(p_1,\mathrm{},p_n)`$ be a non trivial solution. Without loss of generality, we can suppose $`p_10.`$ One of the rows of $`M_t^{}(a)`$ is indexed by $`X_1^t.`$ Replace all the elements in that row as follows:
1. if the element belongs to a column indexed by a monomial $`X^\gamma ,|\gamma |=t_nt,`$ then replace it with $`\mathrm{\Delta }_\gamma (a);`$
2. if it belongs to a column indexed by a monomial $`X^\gamma S_{td_i},`$ replace it with $`X^\gamma f_i(a).`$
It is easy to check that, the determinant of the modified matrix is equal to $`X_1^tdet(M_t^{}(a)).`$ Now, we claim that under the specialization $`X_ip_i,`$ the determinant of the modified matrix will be equal to zero if and only if $`det(M_t^{}(a))=0.`$
In order to prove this, we will show that the following submatrix of size $`\left(i(t_nt)+1\right)\times \left(\genfrac{}{}{0pt}{}{n+t1}{n1}\right)`$ has rank less or equal than $`i(t_nt):`$
$$\left[\begin{array}{ccc}\mathrm{\Delta }_{\gamma _1}(a)(p)& \mathrm{}& \mathrm{\Delta }_{\gamma _s}(a)(p)\\ & {}_{}{}^{𝐭}F_{t_nt}^{}(a)& \end{array}\right].$$
This, combined with a Laplace expansion of the determinant of the modified matrix, gives the desired result.
If the rank of the block $`\left[{}_{}{}^{𝐭}F_{t_nt}^{}(a)\right]`$ is less than $`i(t_nt),`$ then the claim follows straightforwardly. Suppose this is not the case. Then the family $`\{X^\gamma f_i(a),X^\gamma S_{t_ntd_i}\}`$ is a basis of the piece of degree $`t_nt`$ of the generated ideal $`I(a):=f_1(a),\mathrm{},f_n(a).`$ We will show that in this case the polynomial $`_{|\gamma |=t_nt}\mathrm{\Delta }_\gamma (a)(p)X^\gamma `$ belongs to $`I(a)`$, which proves the claim.
Because of (3) and (4), the polynomial $`\left(X_1Y_1\right)\mathrm{\Delta }(a)(X,Y)`$ lies in the ideal $`f_1(a)(X)f_1(a)(Y),\mathrm{},f_n(a)(X)f_n(a)(Y)`$. Specializing $`Y_ip_i,`$ we deduce that $`(X_1p_1)_{j=0}^{t_n}(_{|\gamma |=j}\mathrm{\Delta }_\gamma (a)(p)X^\gamma )`$ is in the graded ideal $`I(a).`$ This, combined with the fact that $`p_10,`$ proves that $`_{|\gamma |=j}\mathrm{\Delta }_\gamma (a)(p)X^\gamma I(a)`$ for all $`j.`$
In particular, $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ divides $`det(M_t).`$ We describe the extraneous factor explicitly in the following theorem, which is the main result in this section. Before stating it, we set the following convention: if the matrix $`𝔼_t`$ is indexed by an empty set, we define $`det\left(𝔼_t\right)=1.`$
###### Theorem 3.2.
For any $`t0,`$ $`det\left(M_t\right)0`$ and $`det(𝔼_t)0.`$
Moreover, we have the following formula à la Macaulay :
$$\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)=\pm \frac{det(M_t)}{det(𝔼_t)}.$$
For the proof of Theorem 3.2, we will need the following auxiliary lemma. Let $`D_t`$ and $`E_t`$ be the matrices defined in §2 before Lemma 2.3.
###### Lemma 3.3.
Let $`t0,`$ and $`\mathrm{\Lambda }`$ a ring which contains $`A.`$ Suppose we have a square matrix $`M`$ with coefficients in $`\mathrm{\Lambda }`$ which has the following structure:
$$M=\left[\begin{array}{cc}M_1& D_t\\ M_2& 0\end{array}\right],$$
where $`M_1,M_2`$ are rectangular matrices. Then, there exists an element $`m\mathrm{\Lambda }`$ such that
$$det\left(M\right)=m.det\left(E_t\right)$$
###### Proof.
$`D_t`$ is square if and only if $`t>t_n.`$ (cf. \[21, §3\]). In this case,
$$det(M)=\pm det(M_2)det(D_t);$$
because of Macaulay’s formula (cf. \[21, Th. 5\]), we have that the right hand side equals
$$\pm det(M_2)det(E_t)\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n),$$
and the conclusion follows easily.
Suppose now $`0tt_n.`$ As in the introduction, let $`i(t)`$ denote the dimension of the $`k`$-vector space of elements of degree $`t`$ in the ideal generated by a regular sequence of $`n`$ polynomials with degrees $`d_1,\mathrm{},d_n.`$ Then $`D_t`$ has $`i(t)+H_d\left(t\right)`$ rows and $`i(t)`$ columns, and there is a bijection between the family $``$ of $`H_d\left(t\right)`$ monomials of degree $`t,`$ and the maximal minors $`m_{}`$ of $`D_t.`$ Namely, $`m_{}`$ is the determinant of the square submatrix made by avoiding all rows indexed by monomials in $`.`$
It is not hard to check that $`m_{}`$ is the determinant $`\varphi _{}^{}`$ which is used in , for computing the subresultant associated with the family $`\{X^\gamma \}_\gamma .`$
Now, using the generalized Macaulay’s formula for the subresultant (cf. ), we have that
$$m_{}=\pm det(E_t).\mathrm{\Delta }_{}^t,$$
where $`\mathrm{\Delta }_{}^t`$ is the subresultant associated with the family $`.`$ It is a polynomial in $`A`$ which vanishes under a specialization of the coefficients $`f_1(a),\mathrm{},f_n(a)`$ if and only if the family $`\{X^\gamma \}_\gamma `$ fails to be a basis of the $`t`$ graded piece of the quotient $`k[X_1,\mathrm{},X_n]/f_1(a),\mathrm{},f_n(a)`$ (cf. ).
Let $`m_{}^c`$ be the complementary minor of $`m_{}`$ in $`M`$ (i.e. the determinant of the square submatrix of $`M`$ which is made by deleting all rows and columns that appear in $`m_{}`$). By the Laplace expansion of the determinant, we have that
$$det\left(M\right)=\underset{}{}s_{}m_{}m_{}^c=det\left(E_t\right)\left(\underset{}{}s_{}m_{}^c\mathrm{\Delta }_{}^t\right)$$
with $`s_{}=\pm 1.`$ Setting $`m=_{}s_{}m_{}^c\mathrm{\Delta }_{}^t\mathrm{\Lambda },`$ we have the desired result. ∎
We now give the proof of Theorem 3.2.
###### Proof.
In it is shown that $`det\left(E_t\right)0,t0.`$ This implies that $`det\left(𝔼_t\right)0.`$ In order to prove that $`det\left(𝔼_t\right)=det(E_t)det(E_{t_nt})`$ divides $`det(M_t),`$ we use the following trick: consider the ring $`B:=\left[b_{\alpha _i}\right]_{|\alpha _i|=d_i,i=1,\mathrm{},n},`$ where $`b_{\alpha _i}`$ are new variables, and the polynomials
$$f_{b,i}:=\underset{|\alpha _i|=d_i}{}b_{\alpha _i}X^{\alpha _i}B[X_1,\mathrm{},X_n].$$
Let $`D_t^b`$ the matrix of the linear transformation $`\psi _{2,t}^b`$ determined by the formula (8) but associated with the sequence $`f_{b,1},\mathrm{},f_{b,n}`$ instead of $`f_1,\mathrm{},f_n.`$ Set $`\mathrm{\Lambda }:=[a_{\alpha _i},b_{\alpha _i}]`$, and consider the matrix $`M(a,b)`$ with coefficients in $`\mathrm{\Lambda }`$ given by
$$M(a,b)=\left[\begin{array}{cc}\mathrm{\Delta }_t& D_t\\ {}_{}{}^{𝐭}D_{t_nt}^{b}& 0\end{array}\right].$$
It is easy to see that $`M(a,a)=M_t`$, and because of Lemma 3.3, we have that $`det\left(E_t\right)`$ divides $`det\left(M(a,b)\right)`$ in $`\mathrm{\Lambda }.`$ Transposing $`M(a,b)`$ and using a symmetry argument, again by the same lemma, we can conclude that $`det\left(E_{t_nt}^b\right)`$ divides $`det\left(M(a,b)\right)`$ in $`\mathrm{\Lambda },`$ where $`E_{t_nt}^b`$ has the obvious meaning.
The ring $`\mathrm{\Lambda }`$ is a factorial domain and $`det(E_t)`$ and $`det(E_{t_nt}^b)`$ have no common factors in $`\mathrm{\Lambda }`$ because they depend on different sets of variables. So, we have
$$det\left(M(a,b)\right)=p(a,b)det(E_t)det(E_{t_nt}^b)$$
for some $`p\mathrm{\Lambda }.`$ Now, specialize $`b_{\alpha _i}a_{\alpha _i}.`$ The fact that $`det\left(M_t\right)`$ is a multiple of the resultant has been proved in Proposition 3.1 (see also \[17, Prop. 3.11.19.21\]) for $`0tt_n,`$ and in for $`t>t_n.`$ On the other side, since $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ is irreducible and depends on all the coefficients of $`f_1,\mathrm{},f_n`$ while $`det(E_t)`$ and $`det(E_{t_nt})`$ do not depend on the coefficients of $`f_n,`$ we conclude that $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ divides $`p(a,a).`$ Moreover, the following lemma shows that they have the same degree. Then, their ratio is a rational number $`\lambda .`$ We can see that $`\lambda =\pm 1,`$ considering the specialized family $`X_1^{d_1},\mathrm{},X_n^{d_n}.`$
###### Lemma 3.4.
For each $`i=1,\mathrm{},n`$ the degree $`\mathrm{deg}_{(a_{\alpha _i})}\left(M_t\right)`$ of $`M_t`$ in the coefficients of $`f_i`$ equals
$$\begin{array}{ccc}\mathrm{deg}_{(a_{\alpha _i})}\left(\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)\right)+\mathrm{deg}_{(a_{\alpha _i})}\left(E_t\right)+\mathrm{deg}_{(a_{\alpha _i})}\left(E_{t_nt}\right)& =& \\ d_1\mathrm{}d_{i1}.d_{i+1}\mathrm{}d_n+\mathrm{deg}_{(a_{\alpha _i})}\left(E_t\right)+\mathrm{deg}_{(a_{\alpha _i})}\left(E_{t_nt}\right)& & \end{array}$$
###### Proof.
Set $`J_u(i):=\{X^\gamma S_u,\gamma _id_i,\gamma _j<d_jji\},u=t,t_nt.`$ From the definitions of $`\psi _{2,t}`$ and $`E_t,`$ it is easy to check that, if $`\delta _t`$ is a maximal minor of $`D_t,`$
$$\mathrm{deg}_{(a_{\alpha _i})}\left(\delta _t\right)\mathrm{deg}_{(a_{\alpha _i})}\left(E_t\right)=\mathrm{\#}J_t(i).$$
Using Laplace expansion, it is easy to see that $`det\left(M_t\right)`$ may be expanded as follows
$$det\left(M_t\right)=\underset{\delta _t,\delta _{t_nt}}{}s_\delta m_\delta \delta _t\delta _{t_nt}$$
where $`s_\delta =\pm 1`$, $`\delta _{t_nt}`$ is a maximal minor of $`{}_{}{}^{𝐭}D_{t_nt}^{}`$ and $`m_\delta `$ is a minor of size $`H_d(t)`$ in $`\mathrm{\Delta }_t.`$
As each entry of $`\mathrm{\Delta }_t`$ has degree $`1`$ in the coefficients of $`f_i,`$ the lemma will be proved if we show that
(16)
$$\mathrm{\#}J_t(i)+\mathrm{\#}J_{t_nt}(i)+H_d(t)=d_1\mathrm{}d_{i1}.d_{i+1}\mathrm{}d_n.$$
Now, as already observed in the proof of Lemma 2.3, $`H_d(t)`$ can be computed as the cardinality of the following set:
(17)
$$H_{d,t}:=\{X^\gamma S_t,\gamma _j<d_jj\},$$
and $`d_1\mathrm{}d_{i1}.d_{i+1}\mathrm{}d_n`$ is the cardinality of
$$\mathrm{\Gamma }_i:=\{X_1^{\gamma _1}\mathrm{}X_{i1}^{\gamma _{i1}}X_{i+1}^{\gamma _{i+1}}\mathrm{}X_n^{\gamma _n},\gamma _j<d_jj\}.$$
In order to prove (16) it is enough to exhibit a bijection between $`\mathrm{\Gamma }_i`$ and the disjoint union $`J_t(i)J_{t_nt}(i)H_{d,t}.`$ This is actually a disjoint union for all $`t,`$ unless $`t_nt=t.`$ But what follows shows that the bijection is well defined even in this case.
Let $`X^{\widehat{\gamma }}\mathrm{\Gamma }_i,\widehat{\gamma }=(\gamma _1,\mathrm{},\gamma _{i1},\gamma _{i+1},\mathrm{},\gamma _n)`$ with $`\gamma _j<d_jji.`$ If $`|\widehat{\gamma }|t,`$ then there exists a unique $`\gamma _i`$ such that $`\gamma :=(\gamma _1,\mathrm{},\gamma _n)_0^n`$ verifies $`|\gamma |=t.`$ If $`\gamma _i<d_i,`$ then we send $`X^{\widehat{\gamma }}`$ to $`X^\gamma H_{d,t}.`$ Otherwise, we send it to $`X^\gamma J_t(i).`$
If $`|\widehat{\gamma }|>t,`$ let $`\widehat{\gamma }^{}`$ denote the multiindex
$$(d_11\gamma _1,\mathrm{},d_{i1}1\gamma _{i1},d_{i+1}1\gamma _{i+1},\mathrm{},d_n1\gamma _n).$$
Then, $`|\widehat{\gamma }^{}|<t_nt,`$ and there exists a unique $`\gamma _i`$ such that the multiindex $`\gamma `$ defined by
$$(d_11\gamma _1,\mathrm{},d_{i1}1\gamma _{i1},\gamma _i,d_{i+1}1\gamma _{i+1},\mathrm{},d_n1\gamma _n)$$
has degree $`t_nt.`$ We can send $`X^{\widehat{\gamma }}`$ to $`X^\gamma J_{t_nt}(i)`$ provided that $`\gamma _id_i.`$ Suppose this last statement does not happen, this implies that the monomial with exponent
$$\gamma ^{}:=(\gamma _1,\mathrm{},\gamma _{i1},d_i1\gamma _i,\gamma _{i+1},\mathrm{},d_n)$$
has degree $`t`$ contradicting the fact that $`|\widehat{\gamma }|>t.`$
With these rules, it is straightforward to check that we obtain the desired bijection. ∎
Changing the order of the sequence $`(f_1,\mathrm{},f_n)`$, and applying Theorem (3.2), we deduce that
###### Corollary 3.5.
$`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)=\mathrm{gcd}\{\text{maximal minors of }\stackrel{~}{\mathrm{\Psi }}_t\}.`$
## 4. Estimating the size of $`M_t`$
We have, for each integer $`t0,`$ a matrix $`M_t`$ of size $`\rho (t),`$ where $`\rho `$ was defined in (1), whose determinant is a nontrivial multiple of the resultant, and such that, moreover, its extraneous factor is a minor of it. We want to know which is the smallest matrix we can have.
We can write $`\rho `$ as
$$\rho (t)=\left(\genfrac{}{}{0pt}{}{n+t1}{n1}\right)+\left(\genfrac{}{}{0pt}{}{n+t_nt1}{n1}\right)H_d(t_nt).$$
It is straightforward to check that $`\left(\genfrac{}{}{0pt}{}{n+t1}{n1}\right)+\left(\genfrac{}{}{0pt}{}{n+t_nt1}{n1}\right)`$ is the restriction to the integers of a polynomial $`\varphi (t)`$ in a real variable $`t,`$ symmetric with respect to $`\frac{t_n}{2}`$ (i.e. $`\varphi (\frac{t_n}{2}+t)=\varphi (\frac{t_n}{2}t)`$ for all $`t`$). Moreover, $`\varphi `$ reaches its minimum over $`[0,t_n]`$ at $`t=\frac{t_n}{2}.`$ Since
(18)
$$\rho (t)=\varphi (t)H_d(t)=\varphi (t_nt)H_d(t_nt)=\rho (t_nt),$$
in order to study the behaviour of $`\rho `$ we need to understand how $`H_d(t)`$ varies with $`t`$. We denote as usual the integer part of a real number $`x`$ by the symbol $`\left[x\right].`$
###### Proposition 4.1.
$`H_d(t)`$ is non decreasing on (the integer points of ) the interval $`[0,\left[\frac{t_n}{2}\right]].`$
###### Proof.
We will prove this result by induction on $`n`$. The case $`n=1`$ is obvious since $`t_1=d1`$ and $`H_d(t)=1`$ for any $`t=0,\mathrm{},d1.`$ Suppose then that the statement holds for $`n`$ variables and set
$$\widehat{d}:=(d_1,\mathrm{},d_{n+1})_0^{n+1},$$
$$d:=(d_1,\mathrm{},d_n).$$
Let $`t<t+1\left[\frac{t_{n+1}}{2}\right].`$ We want to see that $`\phi (t):=H_{\widehat{d}}(t+1)H_{\widehat{d}}(t)`$ is non negative. Recall from (17) that, for every $`t_0`$, $`H_{\widehat{d}}(t)`$ equals the cardinality of the set
$$\{\gamma _0^{n+1}:|\gamma |=t,\mathrm{\hspace{0.17em}0}\gamma _id_i1,i=1,\mathrm{},n+1\}.$$
Then, it can also be computed as
$$\underset{j=0}{\overset{d_{n+1}1}{}}\mathrm{\#}\{\widehat{\gamma }_0^n:|\widehat{\gamma }|=tj,\mathrm{\hspace{0.17em}0}\widehat{\gamma }_id_i1,i=1,\mathrm{},n\},$$
which gives the equality $`H_{\widehat{d}}(t)=_{j=0}^{d_{n+1}1}H_d(tj).`$ It follows that $`\phi (t)=H_d(t+1)H_d(t+1d_{n+1}).`$
If $`t+1\left[\frac{t_n}{2}\right],`$ we deduce that $`\phi (t)0`$ by inductive hypothesis. Suppose then that $`t+1`$ is in the range $`\left[\frac{t_n}{2}\right]<t+1\left[\frac{t_{n+1}}{2}\right].`$ As $`H_d(t+1)=H_d(t_nt1),`$ it is enough to show that $`t_nt1t+1d_{n+1}`$ and $`t_nt1\left[\frac{t_n}{2}\right],`$ which can be easily checked, and the result follows again by inductive hypothesis. ∎
###### Corollary 4.2.
The size $`\rho \left(t\right)`$ of the matrix $`M_t`$ is minimal over $`_0`$ when $`t=\left[\frac{t_n}{2}\right].`$
###### Proof.
By (18), $`\rho `$ has a maximum at $`\left[\frac{t_n}{2}\right]`$ over $`[0,t_n]`$ because $`\varphi `$ has a maximum and $`H_d`$ has a minimum. If $`t>t_n,`$ we have that $`\rho \left(t\right)=\left(\genfrac{}{}{0pt}{}{n+t1}{n1}\right).`$ For $`t`$ in this range, it is easy to check that $`\rho \left(t_n\right)=\left(\genfrac{}{}{0pt}{}{n+t_n1}{n1}\right)1<\rho \left(t\right).`$ Then, $`\rho (t)>\rho (t_n)\rho \left(\left[\frac{t_n}{2}\right]\right).`$
###### Remark 4.3.
Note that when $`t_n`$ is odd, $`\rho (\left[\frac{t_n}{2}\right])=\rho (\left[\frac{t_n}{2}\right]+1)`$, and then the size of $`M_t`$ is also minimal for $`t=\left[\frac{t_n}{2}\right]+1`$ in this case.
Denote $`p:=\frac{_{i=1}^nd_i}{n}`$ the average value of the degrees, and set $`q:=\frac{p+1}{2p}.`$ Note that except in the linear case when all $`d_i=1,`$ it holds that $`p>1`$ and $`q<1.`$
###### Proposition 4.4.
Assume $`p>1.`$ The ratio between the size of the smallest matrix $`M_t`$ and the classical Macaulay matrix $`M_{t_n+1}`$ can be bounded by
$$\frac{\rho \left(\left[t_n/2\right]\right)}{\rho \left(t_n+1\right)}2q^{n1}.$$
In particular, it tends to zero exponentially in $`n`$ when the number of variables tends to infinity and $`p`$ remains bigger that a constant $`c>1.`$.
###### Proof.
When $`t_n`$ is even, $`t_n[t_n/2]=[t_n/2]`$ and when $`t_n`$ is odd, $`t_n[t_n/2]=[t_n/2]+1.`$ In both cases,
$$\frac{\rho \left(\left[t_n/2\right]\right)}{\rho \left(t_n+1\right)}\frac{2\left(\genfrac{}{}{0pt}{}{n+\left[t_n/2\right]}{n1}\right)}{\left(\genfrac{}{}{0pt}{}{n+t_n}{n1}\right)}=2\frac{\left(\left[t_n/2\right]+n\right)\mathrm{}\left(\left[t_n/2\right]+2\right)}{\left(t_n+n\right)\mathrm{}\left(t_n+2\right)}=$$
$$=2\left(\frac{\left[t_n/2\right]+n}{t_n+n}\right)\left(\frac{\left[t_n/2\right]+n1}{t_n+n1}\right)\mathrm{}\left(\frac{\left[t_n/2\right]+2}{t_n+2}\right)$$
$$2\left(\frac{\left[t_n/2\right]+n}{t_n+n}\right)^{n1}.$$
Since $`t_n=npn,`$ we deduce that
$$\frac{\left[t_n/2\right]+n}{t_n+n}\frac{\frac{np}{2}+\frac{n}{2}}{np}=\frac{1}{2}+\frac{1}{2p}=q,$$
as wanted. ∎
## 5. Resultant complexes
In this section we consider Weyman’s complexes (cf. , ) and we make explicit the morphisms in these complexes, which lead to polynomial expressions for the resultant via determinantal formulas in the cases described in Lemma 5.3.
We will consider a complex which is a “coupling” of the Koszul complex $`𝐊^{}(t;f_1,\mathrm{},f_n)`$ associated with $`f_1,\mathrm{},f_n`$ in degree $`t`$ and the dual of the Koszul complex $`𝐊^{}(t_nt,f_1,\mathrm{},f_n)^{}`$ associated with $`f_1,\mathrm{},f_n`$ in degree $`t_nt.`$ This complex arises from the spectral sequence derived from the Koszul complex of sheaves on $`^{n1}`$ associated with $`f_1,\mathrm{},f_n`$ twisted by $`𝒪_{^{n1}}(t).`$ Here, $`𝒪_{^{n1}}(t)`$ denotes as usual the $`t`$-twist of the sheaf of regular functions over the $`(n1)`$-projective space $`^{n1}`$ (see for instance \[15, p. 34\]). Its space of global sections can be identified with the space of homogeneous polynomials in $`n`$ variables of degree $`t`$. We make explicit in terms of the Bezoutian the map $`_0`$ (see (10) below) produced by cohomology obstructions. In fact, the non-trivial contribution is given in terms of the mapping $`\psi _{1,t}`$ defined in (7).
Precisely, let $`𝐊^{}(t;f_1,\mathrm{},f_n)`$ denote the complex
(19)
$$\{0K(t)^n\stackrel{\delta _{(n1)}}{}\mathrm{}\stackrel{\delta _1}{}K(t)^1\stackrel{\delta _0}{}K(t)^0\},$$
where
$$K(t)^j=\underset{i_1<\mathrm{}<i_j}{}S_{td_{i_1}\mathrm{}d_{i_j}}$$
and $`\delta _j`$ are the standard Koszul morphisms.
Similarly, let $`𝐊^{}(t_nt;f_1,\mathrm{},f_n)^{}`$ denote the complex
(20)
$$\{K(t_nt)^0\stackrel{\delta _0^{}}{}K(t_nt)^1\stackrel{\delta _1^{}}{}\mathrm{}\stackrel{\delta _n^{}}{}K(t_nt)^n\},$$
where
$$K(t_nt)^j=\underset{i_1<\mathrm{}<i_j}{}S_{t_ntd_{i_1}\mathrm{}d_{i_j}}^{}$$
and $`\delta _j^{}`$ are the duals of the standard Koszul morphisms. Note that in fact $`K(t_nt)^n=0`$ for any $`t0.`$
Now, define $`𝐂^{}(t;f_1,\mathrm{},f_n)`$ to be the following coupled complex
(21)
$$\{0C^n\stackrel{_{(n1)}}{}\mathrm{}\stackrel{_1}{}C^1\stackrel{_0}{}C^0\stackrel{_1}{}\mathrm{}\stackrel{_{n1}}{}C^{n1}0\},$$
where
(22)
$$\begin{array}{cccc}C^j\hfill & =& K(t)^j,\hfill & j=2,\mathrm{},n\hfill \\ C^j\hfill & =& K(t_nt)^{j+1},\hfill & j=1,\mathrm{},n1\hfill \\ C^1\hfill & =& K(t_nt)^0K(t)^1\hfill & \\ C^0\hfill & =& K(t)^0K(t_nt)^1\hfill & \end{array}$$
and the morphisms are defined by
(23)
$$\begin{array}{cccc}_j\hfill & =& \delta _j,\hfill & j=2,\mathrm{},n1\hfill \\ _j\hfill & =& \delta _j^{},\hfill & j=2,\mathrm{},n1\hfill \\ _1\hfill & =& 0\delta _1\hfill & \\ _0\hfill & =& (\psi _{1,t}+\delta _0)\delta _0^{}\hfill & \\ _1\hfill & =& 0+\delta _1^{}\hfill & \end{array}$$
More explicitly, $`_0(T,(g_1,\mathrm{},g_n))=(\psi _{1,t}(T)+\delta _0(g_1,\mathrm{},g_n),\delta _0^{}(T))`$ and $`_1(h,(T_1,\mathrm{},T_n))=\delta _1^{}(T_1,\mathrm{},T_n).`$ Observe that $`_0`$ is precisely the mapping we called $`\stackrel{~}{\mathrm{\Psi }}_t`$ in the previous section.
As in the proof of Proposition 3.1, given an algebraically closed field $`k,`$ and $`a=(a_{\alpha _i})_{|\alpha _i|=d_i,i=1,\mathrm{},n},`$ a point in $`k^N,`$ we denote by $`f_1(a),\mathrm{},`$ $`f_n(a)`$ the polynomials $`k[X]`$ obtained from $`f_1,\mathrm{},f_n`$ when the coefficients are specialized to $`a`$. For any particular choice of coefficients in (21) we get a complex of $`k`$-vector spaces. We will denote the specialized modules and morphisms by $`K(t)^1(a),\delta _0(a),`$ etc. Let $`D`$ denote the determinant (cf. \[15, Appendix A\], ) of the complex of $`A`$-modules (21) with respect to the monomial bases of the $`A`$-modules $`C^{\mathrm{}}`$. This is an element in the field of fractions of $`A`$.
We now state the main result in this section.
###### Theorem 5.1.
The complex (21) is generically exact, and for each specialization of the coefficients it is exact if and only if the resultant does not vanish. For any positive integer $`t`$ we have that
(24)
$$D=\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n),$$
and moreover, $`D`$ equals the greatest common divisor of all maximal minors of a matrix representing the $`A`$-module map $`_0`$.
###### Proof.
For $`t>t_n,`$ we get the Koszul complex in degree $`t,`$ and so the specialized complex at a point $`ak^N`$ is exact if and only if $`f_1(a),\mathrm{},f_n(a)`$ is a regular sequence, i.e. if and only if the resultant does not vanish. The fact that the determinant of this complex equals the resultant goes back to ideas of Cayley; for a proof see , or .
Suppose $`0tt_n.`$ Since $`\delta _0\delta _1=\delta _1^{}\delta _0^{}=0,`$ it is easy to see that (21) is a complex.
Set
$$U:=\{a=(a_{\alpha _i})k^N,i=1,\mathrm{},n,|\alpha _i|=d_i:det(M_t(a))0\}.$$
Note that the open set $`U`$ is non void because the vector of coefficients of $`\{X_1^{d_1},\mathrm{},X_n^{d_n}\}`$ lies in $`U,`$ since in this case $`detM_t=\pm 1.`$ For any choice of homogeneous polynomials $`f_1(a),\mathrm{},f_n(a)k[X]`$ with respective degrees $`d_1,\mathrm{},d_n`$ and coefficients $`a`$ in $`U`$, the resultant does not vanish by Theorem 3.2 and then the specialized Koszul complexes in (19) and (20) are exact.
Then, the dimension $`dim\mathrm{Im}(\delta _0(a))`$ of the image of $`\delta _0(a)`$ equals $`i(t)=dim<f_1(a),\mathrm{},f_n(a)>_t.`$ Similarly, $`dim(\mathrm{ker}(\delta _0^{}(a))=i(t_nt).`$ Therefore,
$`dim\mathrm{ker}(_0(a))dim\mathrm{Im}(_1(a))=dim\mathrm{Im}(\delta _1(a))=`$
$`=dim\mathrm{ker}(\delta _0(a))=dimK(t)^1(a)i(t).`$
On the other side, the fact that $`M_t(a)`$ is non singular of size $`\rho (t)`$ implies that
$`dim\mathrm{ker}(_0(a))dimC^1(a)\rho (t)=`$
$`=dimK(t)^1(a)+dimK(t_nt)^0(a)\rho (t)=`$
$`=\mathrm{Im}K(t)^1(a)+dimS_{t_nt}(a)\rho (t)=`$
$`=dimK(t)^1(a)i(t).`$
Therefore, $`dim\mathrm{Im}(_1(a))=dim\mathrm{ker}(_0(a))`$ and the complex is exact at level $`1.`$
In a similar way, we can check that the complex is exact at level $`0`$, and so the full specialized complex (21) is exact when the coefficients $`a`$ lie in $`U.`$
In order to compute the determinant of the complex in this case, we can make suitable choices of monomial subsets in each term of the complex starting from the index sets that define $`M_t(a)`$ to the left and to the right. Then,
$$D(a)=\frac{detM_t(a)}{p_1(a)p_2(a)},$$
where $`p_1(a)`$ (resp. $`p_2(a)`$) is a quotient of product of minors of the morphisms on the left (resp. on the right).
Taking into account (19) and (20), it follows from that
$$p_1(a)=det(E_t(a)),p_2(a)=det(E_{t_nt}(a)),$$
and so by Theorem 3.2 we have
$`D(a)=\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)(a){\displaystyle \frac{det(𝔼_t)(a)}{det(E_t(a))det(E_{t_nt}(a))}}=`$
$`=\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)(a)`$
for all families of homogeneous polynomials with coefficients $`a`$ in the dense open set $`U`$, and since $`D`$ and the resultant are rational functions, this implies (24), as wanted. Moreover, it follows that the complex is exact if and only if the resultant does not vanish.
The fact that $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ is the greatest common divisor of all maximal minors of the matrix representing $`_0`$ has been proved in Corollary 3.5. ∎
We remark that from the statement of Theorem 5.1 plus a close look at the map at level $`0`$, it is not hard to deduce that for a given specialization of $`f_1,\mathrm{},f_n`$ in $`k`$ with non vanishing resultant, the specialized polynomials $`\mathrm{\Delta }_\gamma (a),|\gamma |=t_nt`$ generate the quotient of the polynomial ring $`k[X]`$ by the ideal $`I(a)=f_1(a),\mathrm{},f_n(a)`$ in degree $`t`$. We can instead use the known dualizing properties of the Bezoutian in case the polynomials define a regular sequence, to provide an alternative proof of Theorem 5.1. This is a consequence of Proposition 5.2 below. We refer to ,\[, Appendix F\], and for the relation between the Bezoutian and the residue (i.e. an associated trace) and we simply recall the properties that we will use.
Assume $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1(a),\mathrm{},f_n(a))`$ is different from zero. This implies that $`f_1(a),\mathrm{},f_n(a)`$ is a regular sequence and their zero locus consists of the single point $`\mathrm{𝟎}k^n.`$ Then, there exists a dualizing $`k`$-linear operator
$$R_0:k[Y]/f_1(a)(Y),\mathrm{},f_n(a)(Y)k,$$
called the residue or trace operator, which verifies
1. $`h(X)=R_0\left(h(Y)\mathrm{\Delta }(a)(X,Y)\right)`$ in the quotient ring $`k[X]/I(a).`$
2. If $`h`$ is homogeneous of degree $`t`$ with $`tt_n,`$ $`R_0(h)=0`$
Then, for every polynomial $`h(X)k[X]`$ of degree $`t,`$ it holds that
(25)
$$h(X)=\underset{|\gamma |=t_nt}{}R_0\left(h(Y)Y^\gamma \right)\mathrm{\Delta }_\gamma (a)(X)\text{ mod }I(a),$$
where $`\mathrm{\Delta }(a)(X,Y)=_{|\gamma |=t_nt}\mathrm{\Delta }_\gamma (a)(X)Y^\gamma `$ as in (4). As a consequence, the family $`\{\mathrm{\Delta }_\gamma (a)(X)\}_{|\gamma |=t_nt},`$ (resp. $`|\gamma |=t`$) generates the graded piece of the quotient in degree $`t`$ (resp. $`t_nt`$). Moreover, it is easy to verify that for any choice of polynomials $`p_i(X,Y),q_i(X,Y)k[X,Y],i=1,\mathrm{},n,`$ the polynomial $`\stackrel{~}{\mathrm{\Delta }}_a(X,Y)`$ defined by
(26)
$$\stackrel{~}{\mathrm{\Delta }}_a(X,Y):=\mathrm{\Delta }(a)(X,Y)+\underset{i=1}{\overset{n}{}}p_i(X,Y)f_i(a)(X)+q_i(X,Y)f_i(a)(Y).$$
has the same dualizing properties as $`\mathrm{\Delta }(a)(X,Y)`$.
We are ready to prove a kind of “converse” to Proposition 3.1.
###### Proposition 5.2.
If $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1(a),\mathrm{},f_n(a))0,`$ it is possible to extract a square submatrix $`M_t^{}`$ of $`\stackrel{~}{\mathrm{\Psi }}_t`$ as in (15) such that $`det\left(M_t^{}(a)\right)0.`$
###### Proof.
Since $`f_1(a)(X),\mathrm{},f_n(a)(X)`$ is a regular sequence in $`k[X],`$ the dimensions of the graded pieces of the quotient $`k[X]/I(a)`$ in degrees $`t`$ and $`t_nt`$ are $`i(t)`$ and $`i(t_nt)`$ respectively.
We can then choose blocks $`F_t`$ and $`F_{t_nt}`$ as in (15) such that $`F_t(a)`$ and $`F_{t_nt}(a)`$ have maximal rank. Suppose without loss of generality that the blocks $`F_t`$ and $`F_{t_nt}`$ have respectively the form $`\left[\begin{array}{c}Q_t\\ R_t\end{array}\right]`$ and $`\left[\begin{array}{c}Q_{t_nt}\\ R_{t_nt}\end{array}\right]`$, where $`Q_t(a)`$ and $`Q_{t_nt}(a)`$ are square invertible matrices of maximal size. We are going to prove that, with this choice, the matrix $`M_t^{}(a)`$ is invertible.
Our specialized matrix will look as follows:
$$M_t^{}(a)=\left[\begin{array}{cc}\mathrm{\Delta }_t(a)& \begin{array}{c}Q_t(a)\\ R_t(a)\end{array}\\ {}_{}{}^{𝐭}Q_{t_nt}^{}(a)^𝐭R_{t_nt}(a)& 0\end{array}\right].$$
Applying linear operations in the rows and columns of $`M_t^{}(a),`$ it can be transformed into:
$$\left[\begin{array}{ccc}0& 0& Q_t(a)\\ 0& \stackrel{~}{\mathrm{\Delta }}_{t,a}& R_t(a)\\ {}_{}{}^{𝐭}Q_{t_nt}^{}(a)& {}_{}{}^{𝐭}R_{t_nt}^{}(a)& 0\end{array}\right],$$
where the block $`[\stackrel{~}{\mathrm{\Delta }}_{t,a}]`$ is square and of size $`H_d(t).`$
But it is easy to check that this $`\stackrel{~}{\mathrm{\Delta }}_{t,a}`$ corresponds to the components in degree $`t`$ of another Bezoutian $`\stackrel{~}{\mathrm{\Delta }}_a(X,Y)`$ (in the sense of (26)). This is due to the fact that each of the linear operations performed on $`M_t^{}(a),`$ when applied to the block $`\mathrm{\Delta }_{t,a},`$ can be read as a polynomial combination of $`f_i(a)(X)`$ and $`f_i(a)(Y)`$ applied to the bezoutian $`\mathrm{\Delta }(a)(X,Y).`$
Using the fact that the polynomials $`\stackrel{~}{\mathrm{\Delta }}_{\gamma ,a}(X)`$ read in the columns of $`\stackrel{~}{\mathrm{\Delta }}_{t,a}`$ generate the quotient in degree $`t_nt`$ and they are as many as its dimension, we deduce that they are a basis and so
$$det\left(\begin{array}{cc}0& \stackrel{~}{\mathrm{\Delta }}_t\\ {}_{}{}^{𝐭}Q_{t_nt}^{}(a)& {}_{}{}^{𝐭}R_{t_nt}^{}(a)\end{array}\right)0,$$
which completes the proof of the claim. ∎
We could then avoid the consideration of the open set $`U`$ in the proof of Theorem 5.1, and use Proposition 5.2 to show directly that the complex is exact outside the zero locus of the resultant. In fact, this is not surprising since for all specializations such that the resultant is non zero, the residue operator defines a natural duality between the $`t`$-graded piece of the the quotient of the ring of polynomials with coefficients in $`k`$ by the ideal $`I(a)`$ and the $`t_nt`$ graded piece of the quotient, and we can read dual residue bases in the Bezoutian.
We characterize now those data $`n,d_1,\mathrm{},d_n`$ for which we get a determinantal formula.
###### Lemma 5.3.
Suppose $`d_1d_2\mathrm{}d_n.`$ The determinant of the resultant complex provides a determinantal formula for the resultant $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ if and only if the following inequality is verified
(27)
$$d_3+\mathrm{}+d_nn<d_1+d_21.$$
Moreover, when (27) holds, there exists a determinantal formula given by the resultant complex for each $`t`$ such that
(28)
$$d_3+\mathrm{}+d_nn<t<d_1+d_2.$$
###### Remark 5.4.
When all $`d_i`$ have a common value $`d`$, (27) reads
$$(n2)d<2d+n1,$$
which is true for any $`d`$ for $`n4`$, for $`d=1,2,3`$ in case $`n=5`$, for $`d=1,2`$ in case $`n=6`$, and never happens for $`n>7`$ unless $`d=1,`$ as we quoted in the introduction.
###### Proof.
The determinant of the resultant complex provides a determinantal formula for $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ precisely when $`C^2=C_1=0.`$ This is respectively equivalent to the inequalities
$$t<d_1+d_2$$
and
$$t_nt=d_1+\mathrm{}d_nnt<d_1+d_2,$$
from which the lemma follows easily. We have decreased the right hand side of (27) by a unit in order to allow for a natural number $`t`$ satisfying (28). ∎
###### Corollary 5.5.
For all $`n7`$ there exists a determinantal formula only if $`d_1=d_2=d_3=1`$ and $`n3d_4+\mathrm{}+d_n<n,`$ which forces all $`d_i`$ to be $`1`$ or at most, all of them equal $`1`$ except for two of them which equal $`2`$, or all of them equal $`1`$ except for one of them which equals $`3`$.
The proof of the corollary follows easily from the inequality (27). In any case, if a determinantal formula exists, we have a determinantal formula for $`t=[t_n/2],`$ as the following proposition shows.
###### Proposition 5.6.
If a determinantal formula given by the resultant complex exists, then $`M_{[t_n/2]}`$ is square and of the smallest possible size $`\rho (\left[\frac{t_n}{2}\right]).`$
###### Proof.
In order to prove that $`M_{[t_n/2]}`$ is square, we need to check by Lemma 5.3 that
(29)
$$d_3+\mathrm{}+d_nn<\left[\frac{t_n}{2}\right]<d_1+d_2.$$
If there exists a determinantal formula, then the inequalily (27) holds, from which it is straightforward to verify that
$$d_3+\mathrm{}+d_nn<\frac{t_n}{2}<d_1+d_2.$$
To see that in fact (29) holds, it is enough to check that
$$d_3+\mathrm{}+d_nn+1/2\frac{t_n}{2}=\frac{d_1+\mathrm{}+d_nn}{2}.$$
But if the equality holds, we would have that $`d_3+\mathrm{}+d_n=d_1+d_2+n1,`$ which is a contradiction. According to Corollary 4.2, we also know that $`M_{[t_n/2]}`$ has the smallest possible size. ∎
## 6. Dixon formulas
We prove in this section that “affine” Dixon formulas can in fact be recovered in this setting. We first recall classical Dixon formulas to compute the resultant of three bivariate affine polynomials of degree $`d.`$ We will make a slight change of notation in what follows. The input affine polynomials (having monomials of degree at most $`d`$ in two variables $`(X_1,X_2)`$) will be denoted $`f_1,f_2,f_3`$ and we will use capital letters $`F_1,F_2,F_3`$ to denote the homogeneous polynomials in three variables given by their respective homogenizations (with homogeneizing variable $`X_3`$). Dixon (cf. ) proposed the following determinantal formula to compute the resultant $`\mathrm{Res}_{d,d,d}(f_1,f_2,f_3)`$ $`=\mathrm{Res}_{d,d,d}(F_1,F_2,F_3)`$:
Let $`\mathrm{Bez}(X_1,X_2,Y_1,Y_2)`$ denote the polynomial obtained by dividing the following determinant by $`(X_1Y_1)(X_2Y_2)`$:
$$det\left(\begin{array}{ccc}f_1(X_1,X_2)& f_2(X_1,X_2)& f_3(X_1,X_2)\\ f_1(Y_1,X_2)& f_2(Y_1,X_2)& f_3(Y_1,X_2)\\ f_1(Y_1,Y_2)& f_2(Y_1,Y_2)& f_3(Y_1,Y_2).\end{array}\right)$$
Note that by performing row operations we have that $`\mathrm{Bez}(X_1,X_2,Y_1,Y_2)`$ equals the determinant of the matrix
$$det\left(\begin{array}{ccc}\mathrm{\Delta }_{11}& \mathrm{\Delta }_{21}& \mathrm{\Delta }_{31}\\ \mathrm{\Delta }_{12}& \mathrm{\Delta }_{22}& \mathrm{\Delta }_{32}\\ f_1(Y_1,Y_2)& f_2(Y_1,Y_2)& f_3(Y_1,Y_2),\end{array}\right)$$
where $`\mathrm{\Delta }_{ij}`$ are as in (3). Write
$$\mathrm{Bez}(X_1,X_2,Y_1,Y_2)=\underset{|\beta |2d2}{}B_\beta (X_1,X_2)Y_1^{\beta _1}Y_2^{\beta _2}.$$
Set $`A:=[a],`$ where $`a`$ denotes one indeterminate for each coefficient of $`f_1,f_2,f_3.`$ Let $`S`$ denote the free module over $`A`$ with basis $``$ given by all monomials in two variables of degree less or equal than $`d2,`$ which has an obvious isomorphism with the free module $`S^{}`$ over $`A`$ with basis $`^{}`$ given by all monomials in three variables of degree equal to $`d2.`$ The monomial basis of all polynomials in two variables of degree less or equal than $`2d2`$ will be denoted by $`𝒞.`$
Let $`M`$ be the square matrix of size $`2d^2d`$ whose columns are indexed by $`𝒞`$ and whose rows contain consecutively the expansion in the basis $`𝒞`$ of $`mf_1,`$ of $`mf_2,`$ and of $`mf_3,`$ where $`m`$ runs in the three cases over $``$, and finally, the expansion in the basis $`𝒞`$ of all $`B_\beta ,|\beta |d1.`$ Then, Dixon’s formula says that
$$\mathrm{Res}_{d,d,d}(f_1,f_2,f_3)=\pm detM.$$
Here, $`d_1=d_2=d_3=d`$ and $`n=3,`$ so that (27) holds and by (28) there is a determinantal formula for each $`t`$ such that $`d3<t<2d.`$ So, one possible choice is $`t=2d2.`$ Then, $`t_3t=d1<d,`$ which implies $`<F_1,F_2,F_3>_{t_3t}=0.`$ Also, $`td=d2<d,`$ and therefore $`S^{t,i}=S^{},`$ for all $`i=1,2,3.`$
Let $`\mathrm{\Delta }(X_1,X_2,X_3,Y_1,Y_2,Y_3)=_{|\gamma |3d3}\mathrm{\Delta }_\gamma (X)Y^\gamma `$ be the Bezoutian associated with the homogeneous polynomials $`F_1,F_2,F_3.`$ We know that $`\mathrm{Res}_{d,d,d}(F_1,F_2,F_3)=\pm detM_{2d2}.`$ In this case, the transposed matrix $`M_{2d2}^t`$ is a square matrix of the same size as $`M`$, and it is obvious that their $`3d(d1)/2`$ first rows coincide (if the columns are ordered conveniently). According to (7), the last $`(d+1)d/2`$ rows of $`M_{2d2}^t`$ contain the expansion in the basis $`^{}`$ of all $`\mathrm{\Delta }_\gamma ,|\gamma |=d1.`$
###### Proposition 6.1.
The “affine” matrix $`M`$ and the “homogeneous” matrix $`M_{2d2}^t`$ coincide.
###### Proof.
Denote $`P(X_1,X_2,X_3,Y_1,Y_2,t)`$ the homogeneous polynomial of degree $`3d2`$ in $`6`$ variables obtained by dividing the following determinant by $`(X_1Y_1)(X_2Y_2):`$
$$det\left(\begin{array}{ccc}\mathrm{\Delta }_{1,1}(F)& \mathrm{\Delta }_{2,1}(F)& \mathrm{\Delta }_{3,1}(F)\\ \mathrm{\Delta }_{1,2}(F)& \mathrm{\Delta }_{2,2}(F)& \mathrm{\Delta }_{3,2}(F)\\ F_1(Y_1,Y_2,t)& F_2(Y_1,Y_2,t)& F_3(Y_1,Y_2,t)\end{array}\right),$$
where
$$\mathrm{\Delta }_{i,1}(F):=F_i(X_1,X_2,X_3)F_i(Y_1,Y_2,X_3),i=1,2,3$$
and
$$\mathrm{\Delta }_{i,2}(F):=F_i(Y_1,X_2,X_3)F_i(Y_1,Y_2,X_3),i=1,2,3.$$
It is easy to check that
(30)
$$(X_3Y_3)\mathrm{\Delta }(x,y)=P(X_1,X_2,X_3,Y_1,Y_2,X_3)P(X_1,X_2,X_3,Y_1,Y_2,Y_3)$$
and that
(31)
$$P(X_1,X_2,1,Y_1,Y_2,1)=\mathrm{Bez}(X_1,X_2,Y_1,Y_2)$$
We are looking for the elements in $`\mathrm{Bez}(X_1,X_2,Y_1,Y_2)`$ of degree less or equal than $`d1`$ in the variables $`Y_1,Y_2.`$ But it is easy to check that $`\mathrm{deg}_y\left(P(X_1,X_2,X_3,Y_1,Y_2,Y_3)\right)d.`$ This, combined with the equality given in (30), implies that, for each $`1jd1:`$
$$X_3\underset{|\gamma |=j}{}\mathrm{\Delta }_\gamma (X)Y^\gamma Y_3\underset{|\gamma |=j1}{}\mathrm{\Delta }_\gamma (X)Y^\gamma $$
is equal to the piece of degree $`j`$ in the variables $`Y_i`$ of the polynomial $`P(X_1,X_2,X_3,Y_1,Y_2,Y_3).`$
Besides, this polynomial does not depend on $`Y_3,`$ so the following formula holds for every pair $`\gamma ,\stackrel{~}{\gamma }`$ such that $`\gamma =\stackrel{~}{\gamma }+(0,0,k),|\gamma |=j:`$
(32)
$$X_3^k\mathrm{\Delta }_\gamma (X)=\mathrm{\Delta }_{\stackrel{~}{\gamma }}(X).$$
This allows us to compute $`\mathrm{\Delta }_\gamma (X)`$ for every $`|\gamma |=d1,`$ in terms of the homogeneization of $`B_{(\gamma _1,\gamma _2)}.`$ From equation (32), the claim follows straightforwardly. ∎
We conclude that Dixon’s formula can be viewed as a particular case of the determinantal expressions that we addressed. Moreover, Proposition 6.1 can be extended to any number of variables and all Dixon matrices as in \[14, §3.5\] can be recovered in degrees $`t`$ such that $`\psi _{2,t_nt}^{}=0,`$ i.e. such that $`t_nt>t_n\mathrm{min}\{d_1,\mathrm{},d_n\}.`$ As we have seen, all one can hope in general is the explicit quotient formula we give in Theorem 3.2. In fact, we have the following consequence of Lemma 5.3
###### Lemma 6.2.
There exists a determinantal Dixon formula if and only if $`n=2,`$ or $`n=3`$ and $`d_1=d_2=d_3,`$ i.e. in the case considered by Dixon.
###### Proof.
Assume $`d_1d_2\mathrm{}d_n.`$ If inequality (28) is verified for $`t>t_nd_1,`$ we deduce that
(33)
$$(n2)d_1nd_3+\mathrm{}+d_nn<d_12,$$
and so $`(n3)d_1<n2.`$ This equality cannot hold for any natural number $`d_1`$ unless $`n3.`$ It is easy to check that for $`n=2`$ there exist a determinantal Dixon formula for any value of $`d_1,d_2.`$ In case $`n=3`$, (33) implies that $`d_3<d_1+1.`$ Then, $`d_1=d_2=d_3,`$ as claimed. ∎
## 7. Other known formulas and some extensions
We can recognize other well known determinantal formulas for resultants in this setting.
### 7.1. Polynomials in one variable
Let
$$f_1(x)=\underset{j=0}{\overset{d_1}{}}a_jx^j,f_2(x)=\underset{j=0}{\overset{d_2}{}}b_jx^j$$
be generic univariate polynomials (or their homogenizations in two variables) of degrees $`d_1d_2.`$ In this case, inequality (28) is verified for all $`t=0,\mathrm{},d_1+d_21`$ and so we have a determinantal formula for all such $`t`$. Here, $`t_2=d_1+d_22.`$ When $`t=d_1+d_21=t_2+1`$ we have the classical Sylvester formula.
Assume $`d_1=d_2=d`$ and write
$$\frac{f_1(x)f_2(y)f_1(y)f_2(x)}{xy}=\underset{i,j=0}{\overset{d}{}}c_{ij}x^iy^j.$$
Then, the classical Bézout formula for the resultant between $`f_1`$ and $`f_2`$ says that
$$\mathrm{Res}_{d,d}(f_1,f_2)=det(c_{ij}).$$
It is easy to see that we obtain precisely this formulation for $`t=d1.`$ For other values of $`t`$ we get formulas interpolating between Sylvester and Bézout as in \[15, Ch. 12\], even in case $`d_1d_2`$. It is easy to check that the smallest possible matrix has size $`d_2.`$
Suppose for example that $`d_1=1,d_2=2.`$ In this case, $`\left[t_2/2\right]=\left[1/2\right]=0,`$ and $`M_0`$ is a $`2\times 2`$ matrix representing a map from $`S_1^{}`$ to $`S_0S_0^{},`$ whose determinant equals the resultant
$$\mathrm{Res}_{1,2}(f_1,f_2)=a_1^2b_0a_0a_1b_1+b_2a_0^2.$$
If we write $`f_1(x)=0x^2+a_1x+a_0`$ and we use the classical Bezout formula for $`d=2,`$ we would also get a $`2\times 2`$ matrix but whose determinant equals $`b_2\mathrm{Res}_{1,2}(f_1,f_2).`$ The exponent $`1`$ in $`b_2`$ is precisely the difference $`d_2d_1.`$
### 7.2. Sylvester formula for three ternary quadrics
Suppose that $`n=3,d_1=d_2=d_3=2`$ and $`20.`$ Let $`J`$ denote the Jacobian determinant associated with the homogeneous polynomials $`f_1,f_2,f_3.`$ A beautiful classical formula due to Sylvester says that the resultant $`\mathrm{Res}_{2,2,2}(f_1,f_2,f_3)`$ can be obtained as $`1/512`$ times the determinant of the $`6\times 6`$ matrix whose columns are indexed by the monomials in $`3`$ variables of degree $`2`$ and whose rows correspond to the expansion in this monomial basis of $`f_1,f_2,f_3,\frac{J}{X_1},\frac{J}{X_2}`$ and $`\frac{J}{X_3}.`$ In this case, $`\left[t_3/2\right]=\left[3/2\right]=1,`$ and by Lemma 5.3 we have a determinantal formula in this degree since $`23<1<4.`$ From Euler equations
$$2f_i(X)=\underset{j=1}{\overset{3}{}}X_j\frac{f_i(X)}{X_j},$$
we can write
$$\begin{array}{ccc}2\left(f_i(X)f_i(Y)\right)& =& _{j=1}^3\left(X_j\frac{f_i(X)}{X_j}Y_j\frac{f_i(Y)}{Y_j}\right)\hfill \\ & =& _{j=1}^3(X_jY_j)\frac{f_i(X)}{X_j}+Y_j\left(\frac{f_i(X)}{X_j}\frac{f_i(Y)}{Y_j}\right)\hfill \\ & =& _{j=1}^3\left((X_jY_j)\frac{f_i(X)}{X_j}+Y_j_{l=1}^3\frac{^2f_i(X)}{X_jX_l}(X_lY_l)\right).\hfill \end{array}$$
Because of (26), we can compute the Bezoutian using
$$\mathrm{\Delta }_{ij}(X,Y):=\frac{1}{2}\left(\frac{f_i(X)}{X_j}+\underset{l=1}{\overset{3}{}}\frac{^2f_i(X)}{X_lX_j}Y_l\right).$$
Using this formulation, it is not difficult to see that we can recover Sylvester formula from the equality $`\mathrm{Res}_{2,2,2}(f_1,f_2,f_3)=detM_1.`$
### 7.3. Jacobian formulations
When $`t=t_n,`$ one has $`H_d(t)=1,`$ and via the canonical identification of $`S_0^{}`$ with $`A,`$ the complex (22) reduces to the following modified Koszul Complex:
(34)
$$0K(t)^n\stackrel{\delta _{(n1)}}{}\mathrm{}\stackrel{\delta _1}{}AK(t)^1\stackrel{\delta _0}{}K(t)^00,$$
where $`\delta _0`$ is the following map:
$$\begin{array}{ccccc}A& & S_{t_nd_1}\mathrm{}S_{t_nd_n}& & S_{t_n}\\ (\lambda & ,& g_1,\mathrm{},g_n)& & \lambda \mathrm{\Delta }_0+_{i=1}^ng_if_i,\end{array}$$
and $`\mathrm{\Delta }_0:=\mathrm{\Delta }(X,0).`$ As a corollary of Theorem 5.1 we get that, for every specialization of the coefficients, $`\mathrm{\Delta }_0`$ is a non-zero element of the quotient if the resultant does not vanish.
Assume that the characteristic of $`k`$ does not divide the product $`d_1\mathrm{}d_n.`$ It is a well-known fact that the jacobian determinant $`J`$ of the sequence $`(f_1,\mathrm{},f_n)`$ is another non-zero element of degree $`t_n,`$ which is a non-zero element of the quotient whenever the resultant does not vanish (cf. for instance ). In fact, one can easily check that
(35)
$$J=d_1\mathrm{}d_n\mathrm{\Delta }_0modf_1,\mathrm{},f_n.$$
In , the same complex is considered in a more general toric setting, but using $`J`$ instead of $`\mathrm{\Delta }_0`$ . Because of (35), we can recover their results in the homogeneous case.
###### Theorem 7.1.
Consider the modified complex (34) with $`J`$ instead of $`\mathrm{\Delta }_0.`$ Then, for every specialization of the coefficients, the complex is exact if and only if the resultant does not vanish. Moreover, the determinant of the complex equals $`d_1\mathrm{}d_n\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n).`$
We can also replace $`\mathrm{\Delta }_0`$ by $`J`$ in Macaulay’s Formula (Theorem 3.2), and have the following result:
###### Theorem 7.2.
Consider the square submatrix $`\stackrel{~}{M}_{t_n}`$ which is extracted from the matrix of $`\delta _0`$ in the monomial bases, choosing the same rows and columns of $`M_{t_n}.`$ Then, $`det(\stackrel{~}{M}_{t_n})0,`$ and we have the following formula à la Macaulay:
$$d_1\mathrm{}d_n\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)=\frac{det(\stackrel{~}{M}_{t_n})}{det(E_{t_n})}.$$
We end the paper by addressing two natural questions that arise:
### 7.4. Different choices of monomial bases
Following Macaulay’s original ideas, one can show that there is some flexibility in the choice of the monomial bases defining $`S^{t,i}`$ in order to get other non-zero minors, of $`\stackrel{~}{\mathrm{\Psi }}_t`$, i.e different square matrices $`M_t^{}`$ whose determinants are non-zero multiples of $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ with different extraneous factors $`det(E_t^{}),det(E_{t_nt}^{})`$, for appropiate square submatrices $`E_t^{},E_{t_nt}^{}`$ of $`M_t^{}.`$ Besides the obvious choices coming from a permutation in the indices of the variables, other choices can be made as follows.
For any $`i=1,\mathrm{},n,`$ set $`\widehat{d}_i:=(d_1,\mathrm{},d_{i1},d_{i+1},\mathrm{},d_n)`$ and define $`H_{\widehat{d}_i}(t)`$ for any positive integer $`t`$ by the equality
$$\frac{_{ji}\left(1Y^{d_j}\right)}{\left(1Y\right)^{n1}}=\underset{t=0}{\overset{\mathrm{}}{}}H_{\widehat{d}_i}(t).Y^t.$$
For each $`t_0,`$ set also $`\mathrm{\Lambda }_t:=\{X^\gamma S_t:\gamma _j<d_j,j=1,\mathrm{},n\}.`$
We then have the following result:
###### Proposition 7.3.
Let $`M_t^{}`$ a square submatrix of $`\stackrel{~}{\mathrm{\Psi }}_t`$ of size $`\rho (t).`$ Denote its blocks as in (15). Suppose that, for each $`i=1,\mathrm{},n`$, the block $`F_t`$ has exactly $`H_{\widehat{d}_i}(td_i)`$ of its columns corresponding to $`f_i`$ in common with the matrix $`D_t`$ defined in (8) and, also, the block $`F_{t_nt}`$ shares exactly $`H_{\widehat{d}_i}(t_ntd_i)`$ columns corresponding to $`f_i`$ with $`D_{t_nt}.`$ Then, if $`det\left(M_t^{}\right)`$ is not identically zero, the resultant $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ can be computed as the ratio $`\frac{det(M_t^{})}{det(𝔼_t^{})},`$ where $`𝔼_t^{}`$ is made by joining two submatrices $`E_t^{}`$ of $`F_t`$ and $`E_{t_nt}^{}`$ of $`F_{t_nt}.`$ These submatrices are obtained by omitting the columns in common with $`D_t`$ (resp. $`D_{t_nt}`$) and the rows indexed by all common monomials in $`D_t`$ (resp. $`D_{t_nt}`$) and all monomials in $`\mathrm{\Lambda }_t`$ (resp. $`\mathrm{\Lambda }_{t_nt}`$).
We omit the proof which is rather technical, and based in \[21, 6a\], and ,.
### 7.5. Zeroes at infinity
Given a non-homogeneous system of polynomial equations $`\stackrel{~}{f}_1,\mathrm{},\stackrel{~}{f}_n`$ in $`n1`$ variables with respective degrees $`d_1,\mathrm{},d_n,`$ we can homogenize these polynomials and consider the resultant $`\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1,\mathrm{},f_n)`$ associated with their respective homogenizations $`f_1,\mathrm{},f_n.`$ However, this resultant may vanish due to common zeros of $`f_1,\mathrm{},f_n`$ at infinity in projective space $`^{n1}`$ even when there is no affine common root to $`\stackrel{~}{f}_1,=\mathrm{}=\stackrel{~}{f}_n=0.`$ We can in this case extend Canny’s construction of the Generalised Characteristic Polynomial (GCP) for classical Macaulay’s matrices to the matrices $`M_t`$ for any natural number $`t.`$ In fact, when we specialize $`f_i`$ to $`X_i^{d_i}`$ for all $`i=1,\mathrm{},n,`$ the Bezoutian is given by
$$\underset{j_1=0}{\overset{d_11}{}}\mathrm{}\underset{j_n=0}{\overset{d_n1}{}}X_1^{d_11j_1}\mathrm{}X_n^{d_n1j_n}Y_1^{j_1}\mathrm{}Y_n^{j_n},$$
and the specialized matrix $`M_t(e)`$ of $`M_t`$ has a single non zero entry on each row and column which is equal to $`1,`$ so that $`det(M_t(e))=\pm 1.`$ We order the columns in such a way that $`M_t(e)`$ is the identity matrix. With this convention, define the polynomial $`C_t(s)`$ by
$$C_t(s):=\frac{\text{Charpoly }(M_t)(s)}{\text{Charpoly }(𝔼_t)(s)},$$
where $`s`$ denotes a new variable and Charpoly means characteristic polynomial. We then have by the previous observation that
$$C_t(s)=\mathrm{Res}_{d_1,\mathrm{},d_n}(f_1sx_1^{d_1},\mathrm{},f_nsx_1^{d_n}).$$
Moreover, this implies that $`C_t(s)`$ coincides with Canny’s GCP $`C(s),`$ but involves matrices of smaller size. Canny’s considerations on how to compute more efficiently the GCP also hold in this case. Of course, it is in general much better to find a way to construct “tailored” residual resultants for polynomials with a generic structure which is not dense, as in the case of sparse polynomial systems (, ).
Acknowledgements: We are grateful to J. Fernández Bonder, G. Massaccesi and J.M. Rojas for helpful suggestions. We are also grateful to David Cox for his thorough reading of the manuscript.
Author’s addresses:
Departamento de Matemática, F.C.E y N., UBA, (1428) Buenos Aires, Argentina.
cdandrea@dm.uba.ar alidick@dm.uba.ar
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# Untitled Document
McGILL-00-03
UA/NPPS-04-00
Next-to-Leading Order Corrections to Heavy Flavour Production in Longitudinally Polarized Photon-Nucleon Collisions
Z. Merebashvili<sup>a,∗,1</sup>, A.P. Contogouris<sup>a,b,2</sup> and G. Grispos<sup>b,3</sup>
a. Department of Physics, McGill University, Montreal, Qc., H3A 2T8, Canada
b. Nuclear and Particle Physics, University of Athens, Athens 15771, Greece
ABSTRACT
A complete next-to-leading order calculation of longitudinally polarized heavy quark photoproduction is presented. All results of the purturbative calculation are given in detail. For reactions and energies of interest cross sections differential in the transverse momentum and rapidity of the heavy quark, total cross sections and the corresponding asymmetries are given. Errors in the asymmetries are estimated and the possibility to distinguish between various scerarios of the polarized gluon distribution is discussed. Our results are compared with other related publications.
<sup>*</sup><sup>*</sup>footnotetext: Present address: High Energy Physics Institute, Tbilisi State University, University St. 9, 380086 Tbilisi, Republic of Georgia.<sup>1</sup><sup>1</sup>footnotetext: e-mail: mereb@sun20.hepi.edu.ge<sup>2</sup><sup>2</sup>footnotetext: e-mail: apcont@physics.mcgill.ca, acontog@cc.uoa.gr<sup>3</sup><sup>3</sup>footnotetext: e-mail: ggrispos@cc.uoa.gr
I. INTRODUCTION
Deep inelastic scattering of longitudinally polarized particles has provided important information on the spin structure of the nucleon. However, the size and shape of the polarized gluon distribution $`\mathrm{\Delta }g`$ in the proton remains an essential problem. Significant progress requires experiments on reactions with longitudinally polarized particles dominated by subprocesses with initial gluons. Such a reaction is
$$\stackrel{}{\gamma }+\stackrel{}{p}Q\overline{(Q)}+X,$$
(1.1)
where $`Q\overline{(Q)}`$ denotes heavy quark (antiquark); this is dominated by
$$\stackrel{}{\gamma }+\stackrel{}{g}Q+\overline{Q}$$
(1.2)
An experiment closely related to (1.1) is soon going to take place $`[`$1$`]`$ and there are more than one proposals $`[`$2$`]`$.
At the Born level, (1.1) has been studied long ago $`[`$3,4$`]`$. However, the importance of knowing the next-to-leading order corrections (NLOC) cannot be overemphasized. This work presents detailed results on a NLOC calculation.
It should be noted that NLOC for (1.1) have already been published $`[`$5$`]`$. We believe, however, that in view of the importance of (1.1), an independent determination of NLOC in a different regularization approach (see below) is in order. Extensive comparisons with the calculation of $`[`$5$`]`$, as well as certain differences in our view regarding certain questions will be also reported.
At NLO, apart from the loop and gluon Bremsstrahlung (Brems) contributions to (1.2), the subprocesses
$$\stackrel{}{\gamma }+\stackrel{}{q}(\stackrel{}{\overline{q}})Q+\overline{Q}+q(\overline{q}),$$
(1.3)
where $`q`$ denotes a light quark, should also be taken into account.
We note that the Abelian part of NLOC for (1.4) provides the corrections to
$$\stackrel{}{\gamma }+\stackrel{}{\gamma }Q+\overline{Q}$$
(1.4)
This part has already been determined $`[`$6,7$`]`$. NLOC to (1.4) are of interest in themselves in connection with Higgs boson searches when the Higgs mass is in the range of 90 to 160 GeV.
The loop and $`23`$ parton graphs involved in NLOC introduce ultraviolet (UV), infrared (IR) and collinear singularities, which are eliminated by working in $`n=42\epsilon `$ dimensions. For polarized reactions this requires extension of the Dirac matrix $`\gamma _5`$ in $`n4`$ dimensions. Unless otherwise stated, we work in the scheme of dimensional reduction (RD), which simplifies the calculation of the traces. Certain subtleties of RD have been discussed in $`[`$6$`]`$ and are mentioned below. Furthermore, we use parton distributions whose evolution, via 2-loop anomalous dimensions, is determined in a scheme different from RD. This necessitates the addition to our perturbative results of certain conversion terms.
In all the above contributions the photon interacts in a direct way. In addition, there are also resolved contributions, in which it interacts through its partonic constituents; in fact, strictly speaking, at NLO, scheme independent cross sections arise only by adding them. At this moment a complete calculation of the resolved contributions is not possible, and we will be limited in giving an estimate.
The paper is organized as follows. Sect. II contains our general procedures, Sect. III discusses the loop contributions to the photon-gluon fusion subprocess and Sect. IV the corresponding Brems ones. Sect. V presents analytic results on the subprocess (1.3). In Sect. VI we derive the nesessary formulas for calculating various physical observables. Sect. VII presents our numerical results and discusses the possibility to distinguish between three sets differing essentially in the polarized gluon distribution function $`\mathrm{\Delta }g`$. Sect. VIII deals with our comparison with $`[`$5$`]`$, as well as with $`[`$8$`]`$. Sect. IX presents our conclusions. Finally, in three Appendices we present results completing our determination of NLOC.
II. GENERAL PROCEDURES
The Born and the loop contributions to $`\gamma +gQ+\overline{Q}`$ are shown in Fig. 1. With the 4-momenta $`p_i,i=1,\mathrm{},4,`$ as indicated and with $`m`$ the heavy quark mass we define:
$$s(p_1+p_2)^2,tTm^2(p_1p_3)^2m^2,uUm^2(p_2p_3)^2m^2$$
(2.1)
Let $`M_i(\lambda _1,\lambda _2)`$ the amplitude of any of the contributing graphs, where $`\lambda _1,\lambda _2`$ the helicities of the initial partons; our polarized cross sections correspond to the quantities:
$$\frac{1}{2}\mathrm{\Sigma }[M_i(++)\underset{j}{\overset{}{M}}(++)M_i(+)\underset{j}{\overset{}{M}}(+)]$$
(2.2)
where $`\mathrm{\Sigma }`$ denotes summation over the helicities and colors of the final particles and average over the colors of the initial. For the determination of the asymmetries we need also the unpolarized cross sections, which correspond to the average of $`M_i(++)\underset{j}{\overset{}{M}}(++)`$ and $`M_i(+)\underset{j}{\overset{}{M}}(+)`$.
We also introduce
$$v1+t/swu/(s+t)$$
(2.3)
To reduce the length of the subsequent expressions we will make use of the results presented in $`[`$6$`]`$. Thus our leading-order (LO) polarized and unpolarized cross sections are
$$[\mathrm{\Delta }]\frac{d\sigma _{\mathrm{LO}}^{\gamma g}}{dvdw}=\kappa C_F[\mathrm{\Delta }]\frac{d\sigma _{\mathrm{LO}}}{dvdw}$$
(2.4)
where $`\kappa \alpha _s/8\alpha e_Q^2`$ and $`[\mathrm{\Delta }]d\sigma _{\mathrm{LO}}/dvdw`$ the corresponding \[polarized\] unpolarized cross sections for $`\gamma \gamma Q\overline{Q}`$ (Eq. (9) of $`[`$6$`]`$). For later use we note that $`[\mathrm{\Delta }]d\sigma _{\mathrm{LO}}/dvdw`$ are proportional to:
$$\mathrm{\Delta }B(s,t,u)=\frac{1}{s}\left(\frac{t^2+u^2}{tu}+2\frac{sm^2}{tu}\left(\frac{s^2}{tu}2\right)\right)$$
and (see also $`[`$8$`]`$)
$$B(s,t,u)=\frac{1}{s}\left(\frac{t^2+u^2}{tu}+4\frac{sm^2}{tu}\left(1\frac{sm^2}{tu}\right)\right)$$
(2.5)
In determining the loop contributions, the renormalization of the heavy quark mass and wave function were carried on shell, as in $`[`$6$`]`$, i.e. the renormalized heavy quark self-energy $`\mathrm{\Sigma }_r(p)`$ was taken to satisfy at $`p^2=m^2`$:
$$\mathrm{\Sigma }_r(p)=0\frac{}{p}\mathrm{\Sigma }_r(p)=0$$
(2.6)
This determines the mass and wave function renormalization constants $`Z_m`$ and $`Z_2`$ $`[`$6$`]`$.
Dimensional reduction does not automatically satisfy the Ward identity
$$Z_1=Z_2,$$
where $`Z_1`$ is the renormalization constant for the vertex of the graph Fig. 1(d). This requires the introduction of proper finite counterterm, of which the form is given in $`[`$6$`]`$.
In the present case charge renormalization is also required. Defining
$$C_\epsilon (m)\frac{\mathrm{\Gamma }(1+\epsilon )}{(4\pi )^2}\left(\frac{4\pi \mu ^2}{m^2}\right)^\epsilon ,$$
(2.7)
let $`g_0(g)`$ be the bare (renormalized) coupling, $`Z_g=g_0/g`$ the charge renormalization constant and $`b=(11N_C2N_{lf})/6`$, where $`N_{lf}`$ is the number of light flavors. We take
$$Z_g=1\frac{g^2}{\epsilon }\left\{C_\epsilon (M)b\frac{1}{3}C_\epsilon (m)\right\}$$
(2.8)
where $`M`$ is a regularization mass. In this scheme the contribution of a heavy quark loop in the gluon self-energy is subtracted out, i.e. the heavy quark is decoupled $`[`$9,8$`]`$. This is consistent with parton distributions $`\mathrm{\Delta }F_{a/p}(x,Q^2)`$ of which the evolution is determined from split functions involving only light quarks, as is the case of $`\mathrm{\Delta }F_{a/p}`$ used below.
Finally, the renormalization of the $`gQ\overline{Q}`$ vertex was carried using the Slavnov-Taylor identities $`[`$10$`]`$.
III. LOOP CONTRIBUTIONS
The loop graphs contributing to (1.2) are depicted in Fig. 1. The integrals for the Abelian type of graphs (a)-(e) were calculated in $`[`$6$`]`$. The non-Abelian graphs (g) and (h) introduce tensor integrals of the form
$$\frac{d^nq}{(2\pi )^n}\frac{q^\mu ,q^\mu q^\nu }{q^2(qp_2)^2[(q+p_4p_2)^2m^2]}$$
and
$$\frac{d^nq}{(2\pi )^n}\frac{q^\mu ,q^\mu q^\nu ,q^\mu q^\nu q^\rho }{q^2(qp_2)^2[(q+p_4p_2)^2m^2][(qp_3)^2m^2]}$$
As in $`[`$6$`]`$, using Passarino-Veltman techniques $`[`$11$`]`$, we reduce them to scalar ones; those can be found in $`[`$12$`]`$.
The contributions presented below include the $`tu`$ crossing symmetric of Fig. 1 plus UV counterterms; thus they contain no UV singularities.
The graphs (a)-(e) give:
$$\frac{d\sigma _{ae}^{\gamma g}}{dvdw}=\kappa C_F\frac{d\sigma _{\mathrm{vse}}}{dvdw}\frac{N_C}{2}\frac{d\sigma _{ae}}{dv}\delta (1w)$$
(3.1)
where $`d\sigma _{\mathrm{vse}}/dvdw`$ is given in Eq.(16) of $`[`$6$`]`$ and
$`{\displaystyle \frac{d\sigma _{ae}}{dv}}`$ $`=`$ $`K_L\{2\stackrel{~}{A}_1[2(\zeta (2)\mathrm{Li}_2({\displaystyle \frac{T}{m^2}}))(1+3{\displaystyle \frac{m^2}{t}})\mathrm{ln}({\displaystyle \frac{t}{m^2}})(1+{\displaystyle \frac{m^2}{T}})+2]+\stackrel{~}{A}_2\mathrm{ln}({\displaystyle \frac{t}{m^2}})`$ (3.2)
$`+\stackrel{~}{A}_3(\mathrm{Li}_2({\displaystyle \frac{T}{m^2}})\zeta (2))+\stackrel{~}{A}_4+(tu)\}`$
with
$$K_L\frac{1}{8s}\alpha \alpha _s^2e_Q^2$$
Here and subsequently the polarized cross sections are given by (3.1) and (3.2) with $`\mathrm{\Delta }d\sigma /dvdw`$, $`\mathrm{\Delta }d\sigma /dv`$ and $`\mathrm{\Delta }\stackrel{~}{A}_1`$, i=1,…,4, replacing the corresponding unpolarized quantities. The $`[\mathrm{\Delta }]\stackrel{~}{A}_i`$ are given in Appendix A.
Graph (f) contributes:
$$\frac{d\sigma _f}{dvdw}=\kappa (C_F\frac{N_C}{2})\frac{d\sigma _{\mathrm{box}}}{dvdw}$$
(3.3)
with $`d\sigma _{\mathrm{box}}/dvdw`$ in Eq.(22) of $`[`$6$`]`$.
Turning to the non-Abelian graphs, (g) gives
$$\frac{d\sigma _g^{\gamma g}}{dvdw}=2\pi \alpha _sC_\epsilon (m)N_C\frac{d\sigma _{\mathrm{LO}}^{\gamma g}}{dvdw}\left(\frac{1}{\epsilon ^2}+\frac{4}{\epsilon }\right)\frac{N_C}{2}\left(\frac{d\stackrel{~}{\sigma }_g}{dv}+\frac{d\sigma _g}{dv}\right)\delta (1w),$$
(3.4)
where
$$\frac{d\stackrel{~}{\sigma }_g}{dv}=2F(\epsilon )\{A_1\frac{1}{\epsilon }\mathrm{ln}(\frac{t}{m^2})+A_{}^{}{}_{1}{}^{}[\frac{1}{\epsilon ^2}\frac{2}{\epsilon }\frac{2}{\epsilon }\mathrm{ln}(\frac{t}{m^2})]+(tu)\}$$
(3.5)
with
$$F(\epsilon )K_L\mu ^{2\epsilon }\left(\frac{4\pi \mu }{m}\right)^{2\epsilon }\frac{\mathrm{\Gamma }(1+\epsilon )}{\mathrm{\Gamma }(1\epsilon )}\left(\frac{s\mu ^2}{tusm^2}\right)^\epsilon $$
(3.6)
and
$`{\displaystyle \frac{d\sigma _g}{dv}}=K_L\{2A_1[\mathrm{Li}_2({\displaystyle \frac{T}{m^2}})+\mathrm{ln}^2({\displaystyle \frac{t}{m^2}})2]+A_{}^{}{}_{1}{}^{}[4\mathrm{L}\mathrm{i}_2({\displaystyle \frac{T}{m^2}})+4\mathrm{ln}^2({\displaystyle \frac{t}{m^2}})]`$
$`+A_{}^{}{}_{2}{}^{}\mathrm{ln}({\displaystyle \frac{t}{m^2}})+A_{}^{}{}_{3}{}^{}+(tu)\}`$ (3.7)
In (3.5) and (3.7), $`[\mathrm{\Delta }]A_1`$ are given in App. B of $`[`$6$`]`$ and $`[\mathrm{\Delta }]A_{}^{}{}_{i}{}^{}`$, i=1,2,3, in App. A of this paper.
The contribution of the graph (h) is:
$$\frac{d\sigma _h^{\gamma g}}{dvdw}=\frac{N_C}{2}\left\{4\pi \alpha _sC_\epsilon (m)\frac{3}{\epsilon ^2}\frac{d\sigma _{\mathrm{LO}}^{\gamma g}}{dvdw}+\left(\frac{d\stackrel{~}{\sigma }_h}{dv}+\frac{d\sigma _h}{dv}\right)\delta (1w)\right\},$$
(3.8)
where
$$\frac{d\stackrel{~}{\sigma }_h}{dv}=2F(\epsilon )\{A_1\frac{1}{\epsilon }[\mathrm{ln}(\frac{t}{m^2})+2\mathrm{ln}(\frac{u}{m^2})]+A_{}^{}{}_{1}{}^{}[\frac{1}{\epsilon ^2}+\frac{2}{\epsilon }+\frac{2}{\epsilon }\mathrm{ln}(\frac{t}{m^2})]+(tu)\}$$
(3.9)
and
$`{\displaystyle \frac{d\sigma _h}{dv}}=K_L\{A_1[{\displaystyle \frac{35}{4}}\zeta (2)\mathrm{Li}_2({\displaystyle \frac{T}{m^2}})+4\mathrm{ln}({\displaystyle \frac{t}{m^2}})\mathrm{ln}({\displaystyle \frac{u}{m^2}})\mathrm{ln}^2({\displaystyle \frac{t}{m^2}})]+B_{}^{}{}_{1}{}^{}\mathrm{Li}_2({\displaystyle \frac{T}{m^2}})`$
$`+(B_{}^{}{}_{2}{}^{}+A_{}^{}{}_{1}{}^{})\zeta (2)+B_{}^{}{}_{3}{}^{}\mathrm{ln}^2({\displaystyle \frac{t}{m^2}})+B_{}^{}{}_{4}{}^{}\mathrm{ln}({\displaystyle \frac{t}{m^2}})+B_{}^{}{}_{5}{}^{}\mathrm{ln}({\displaystyle \frac{t}{m^2}})\mathrm{ln}({\displaystyle \frac{u}{m^2}})`$
$`+B_{}^{}{}_{6}{}^{}+(tu)\}`$ (3.10)
The coefficients $`[\mathrm{\Delta }]B_{}^{}{}_{i}{}^{}`$, i=1,…,6, are given in App. A.
Finally, after cancellation of the UV singularities, graph (i) does not contribute.
We remark that regarding the terms $`1/\epsilon ^2`$, the contributions of the graphs (g) and (h) taken separately are not proportional to the Born $`d\sigma _{\mathrm{LO}}^{\gamma g}/dvdw`$; only their sum is proportional to the Born. The same holds regarding the terms $`1/\epsilon `$.
IV. GLUON BREMS CONTRIBUTIONS
In this chapter we present complete analytic results for the NLOC arising from Brems. To the best of our knowledge, in relation with heavy quark production, such results have not so far been presented.
With $`k`$ the 4-momentum of the emitted gluon we introduce also
$$s_2(k+p_4)^2m^2=s+t+u=sv(1w)$$
(4.1)
The Brems graphs contributing to the NLOC of (1.2) are shown in Fig. 2A. The squared sum of the corresponding amplitudes (plus those obtained via $`p_1p_2`$) after summing over final spins and colors and averaging over initial colors is given by
$$4m^2|M_{23}^{\gamma g}|^2=K_B(\epsilon )\left(\frac{C_F}{2}G^{\gamma \gamma }\frac{N_C}{16}G^{\gamma g}\right)$$
(4.2)
where $`G^{\gamma \gamma }`$ is the quantity in the square bracket of Eq.(24) of $`[`$6$`]`$ (plus $`p_1p_2`$), $`G^{\gamma g}`$ has the expansion
$`G^{\gamma g}=e_1+{\displaystyle \frac{e_2}{p_2p_4}}+{\displaystyle \frac{e_3}{p_1p_4^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}}}+{\displaystyle \frac{e_4}{p_1p_4}}+{\displaystyle \frac{\stackrel{~}{e}_5}{p_3k}}+{\displaystyle \frac{\stackrel{~}{e}_6}{p_1p_4p_3k}}+{\displaystyle \frac{e_7}{p_1p_4p_2p_4}}+e_8{\displaystyle \frac{p_2k}{p_1p_4}}`$ (4.3)
$`+`$ $`{\displaystyle \frac{\stackrel{~}{e}_9}{p_2p_4p_3k}}+e_{10}{\displaystyle \frac{p_2k}{p_3k}}+{\displaystyle \frac{e_{11}}{p_1p_4^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}p_3k}}+f_1{\displaystyle \frac{p_3k}{p_2p_4}}+f_2{\displaystyle \frac{p_3k^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}}{p_2p_4}}+\stackrel{~}{f}_3{\displaystyle \frac{p_3k^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}}{p_2k}}+\stackrel{~}{f}_4{\displaystyle \frac{p_3k}{p_2k}}`$
$`+`$ $`{\displaystyle \frac{\stackrel{~}{f}_5}{p_2k}}+{\displaystyle \frac{\stackrel{~}{f}_6}{p_1p_4p_2k}}+{\displaystyle \frac{\stackrel{~}{f}_7}{p_2k^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}}}+{\displaystyle \frac{\stackrel{~}{f}_8}{p_1p_4p_2k^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}}}+{\displaystyle \frac{\stackrel{~}{f}_9}{p_2kp_3k}}+{\displaystyle \frac{\stackrel{~}{f}_{10}}{p_1p_4^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}p_2k}}`$
$`+`$ $`{\displaystyle \frac{\stackrel{~}{f}_{11}}{p_1p_4^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}p_2k^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}}}`$
and
$$K_B(\epsilon )=(4\pi )^3\alpha \alpha _s^2e_Q^2\mu ^{6\epsilon }$$
(4.4)
As in Sect. III, $`\mathrm{\Delta }|M_{23}^{\gamma g}|^2`$ is given by (4.2) and (4.3) with $`\mathrm{\Delta }G^{\gamma \gamma }`$, $`\mathrm{\Delta }G^{\gamma g}`$, $`\mathrm{\Delta }e_i`$ and $`\mathrm{\Delta }f_i`$, i=1,2,…,13, replacing the corresponding unpolarized quantities. The coefficients $`[\mathrm{\Delta }]e_i`$, $`[\mathrm{\Delta }]f_i`$ of (4.3) are given in App. B.
The Brems contribution to $`[\mathrm{\Delta }]d\sigma /dvdw`$ is obtained by working in the Gottfried-Jackson frame of $`\overline{Q}(Q)`$ and gluon (c.m. system of $`p_4`$ and $`k`$). Details are given in $`[`$6$`]`$. The terms with coefficients $`[\mathrm{\Delta }]\stackrel{~}{e}_i`$, $`[\mathrm{\Delta }]\stackrel{~}{f}_i`$ in (4.3) give contributions singular at $`s_2=0(w=1)`$ and must be integrated in $`n4`$ dimensions. In view of the fact that the $`23`$ particle phase space is proportional to $`s_2^{12\epsilon }`$ (Eq. (26) of $`[`$6$`]`$), the remaining terms can be integrated in 4 dimensions. The arising integrals are given in $`[`$12$`]`$. Certain terms of special interest not given in $`[`$6$`]`$ are determined in App. C.
Corresponding to the second term in (4.2), with $`\overline{y}\sqrt{(t+u)^24m^2s}`$, $`S_2=s_2+m^2`$ and $`x=(1\beta )/(1+\beta )`$, where $`\beta =\sqrt{14m^2/s}`$, the final result is
$`{\displaystyle \frac{d\sigma _{\mathrm{Br}}^{\gamma g}}{dvdw}}={\displaystyle \frac{K_B(0)}{(4\pi )^3}}{\displaystyle \frac{N_C}{16}}2\pi {\displaystyle \frac{vs_2}{8\pi S_2}}\{e_1+{\displaystyle \frac{2S_2}{s_2(s+u)}}\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}e_2+{\displaystyle \frac{4S_2}{m^2(s+t)^2}}e_3+{\displaystyle \frac{2S_2}{s_2(s+t)}}\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}e_4`$ (4.5)
$`+`$ $`e_7I_8+e_8I_{10}+e_{10}I_{16}+e_{11}I_{13}(tu)+f_1F_1+f_2F_2\}`$
$``$ $`{\displaystyle \frac{1}{(1w)_+}}{\displaystyle \frac{s_2^2}{S_2}}{\displaystyle \frac{N_C}{16}}F(0)\{{\displaystyle \frac{2S_2}{s_2\overline{y}}}\mathrm{ln}{\displaystyle \frac{T+U\overline{y}}{T+U+\overline{y}}}\stackrel{~}{e}_5+\stackrel{~}{e}_6I_{11}(tu)+\stackrel{~}{e}_9I_{11}\}`$
$``$ $`{\displaystyle \frac{vss_2}{S_2}}{\displaystyle \frac{N_C}{16}}F(0)\{\stackrel{~}{f}_3F_3^c+\stackrel{~}{f}_4F_4^c+\stackrel{~}{f}_7F_7^c+\stackrel{~}{f}_8F_8^c+\stackrel{~}{f}_{10}F_{10}^c+\stackrel{~}{f}_{11}F_{11}^c(2\mathrm{ln}{\displaystyle \frac{s_2}{m^2}}+\mathrm{ln}{\displaystyle \frac{m^2}{S_2}})(\stackrel{~}{f}_3F_3^s`$
$`+`$ $`\stackrel{~}{f}_4F_4^s+\stackrel{~}{f}_8F_8^s+\stackrel{~}{f}_{10}F_{10}^s+\stackrel{~}{f}_{11}F_{11}^s)\}`$
$``$ $`{\displaystyle \frac{1}{(1w)_+}}{\displaystyle \frac{s_2^2}{S_2}}{\displaystyle \frac{N_C}{16}}F(0)\{\stackrel{~}{f}_6F_6^c+\stackrel{~}{f}_9F_9^c(2\mathrm{ln}{\displaystyle \frac{sv}{m^2}}+\mathrm{ln}{\displaystyle \frac{m^2}{S_2}})(\stackrel{~}{f}_5F_5^s+\stackrel{~}{f}_6F_6^s+\stackrel{~}{f}_9F_9^s)\}`$
$`+`$ $`2L_+{\displaystyle \frac{s_2^2}{S_2}}{\displaystyle \frac{N_C}{16}}F(0)\{\stackrel{~}{f}_5F_5^s+\stackrel{~}{f}_6F_6^s+\stackrel{~}{f}_9F_9^s\}`$
$`+`$ $`8\pi \alpha _sN_CC_\epsilon (m)[\mathrm{\Delta }]{\displaystyle \frac{d\sigma _{\mathrm{LO}}^{\gamma g}}{dvdw}}\{2\mathrm{ln}^2{\displaystyle \frac{sv}{m^2}}\mathrm{ln}^2(x)+{\displaystyle \frac{1}{2}}\mathrm{ln}^2({\displaystyle \frac{t}{u}}x)+2\mathrm{ln}{\displaystyle \frac{u}{t}}\mathrm{ln}{\displaystyle \frac{sv}{m^2}}+\mathrm{Li}_2(1{\displaystyle \frac{1}{x}}{\displaystyle \frac{u}{t}})`$
$``$ $`\mathrm{Li}_2(1{\displaystyle \frac{1}{x}}{\displaystyle \frac{t}{u}})2\zeta (2){\displaystyle \frac{2m^2s}{s\beta }}[(2\mathrm{ln}{\displaystyle \frac{sv}{m^2}}\mathrm{ln}(x))\mathrm{ln}(x)\mathrm{Li}_2\left({\displaystyle \frac{4\beta }{(1\beta )^2}}\right)]\}`$
In Eq. (4.5), the integrals $`I_i`$ are given in the App. C of $`[`$6$`]`$ and the integrals $`F_i`$ in the App. C of this paper. Also, $`L_+(\mathrm{ln}(1w)/(1w))_+`$, which enters through the relation
$$(1w)^{12\epsilon }=\frac{1}{2\epsilon }\delta (1w)+\frac{1}{(1w)_+}2\epsilon L_+,$$
(4.6)
where the so called ”plus” distributions are defined in a usual way:
$$_0^1𝑑z\frac{f(w)}{(1w)_+}_0^1𝑑z\frac{f(w)f(1)}{1w}.$$
(4.7)
For the second term of (4.2), we give the terms $`1/\epsilon ^2`$ and $`1/\epsilon `$, as well:
$`{\displaystyle \frac{d\sigma _{\mathrm{Br}}^{\gamma g,\epsilon }}{dvdw}}`$ $`=`$ $`{\displaystyle \frac{N_C}{\epsilon ^2}}8\pi \alpha _sC_\epsilon (m)[\mathrm{\Delta }]{\displaystyle \frac{d\sigma _{\mathrm{LO}}^{\gamma g}}{dvdw}}{\displaystyle \frac{N_C}{\epsilon }}8\pi \alpha _sC_\epsilon (m)[\mathrm{\Delta }]{\displaystyle \frac{d\sigma _{\mathrm{LO}}^{\gamma g}}{dvdw}}\{3\mathrm{ln}({\displaystyle \frac{u}{m^2}})\mathrm{ln}({\displaystyle \frac{t}{m^2}})\}`$ (4.8)
$`+`$ $`{\displaystyle \frac{N_C}{\epsilon }}8\pi \alpha _sC_\epsilon (m)[\mathrm{\Delta }]{\displaystyle \frac{d\sigma _{\mathrm{LO}}^{\gamma g}}{dvdw}}{\displaystyle \frac{2m^2s}{s\beta }}\mathrm{ln}(x)`$
$``$ $`{\displaystyle \frac{1}{\epsilon }}{\displaystyle \frac{2sv}{1vw}}F(\epsilon )[\mathrm{\Delta }]P_{gg}^f(x_2)[\mathrm{\Delta }]B(x_2s,t,x_2u)`$
where $`[\mathrm{\Delta }]P_{gg}^f(x)`$ is the 4-dimensional $`ggg`$ split function without the $`\delta (1w)`$ part, $`F(\epsilon )`$, $`[\mathrm{\Delta }]B(s,t,u)`$ are given by (3.6) and (2.5) and
$$x_2=\frac{1v}{1vw}$$
(4.9)
Addition of loop and Brems contributions cancels the singularities $`1/\epsilon ^2`$ and part of the $`1/\epsilon `$. The remaining $`1/\epsilon `$ are cancelled by a factorization counterterm corresponding to the final gluon emitted collinearly with the initial one (Fig. 2A, graph (d)). In the $`\overline{MS}`$ scheme this counterterm gives:
$$[\mathrm{\Delta }]\frac{d\sigma _{\mathrm{ct}}}{dvdw}=\frac{1}{\epsilon }\frac{2sv}{1vw}F(\epsilon )[\mathrm{\Delta }]P_{gg}(x_2)[\mathrm{\Delta }]B(x_2s,t,x_2u)\left(\frac{m^2}{M_F^2}\right)^\epsilon ,$$
(4.10)
$`[\mathrm{\Delta }]P_{gg}(x)`$ the $`ggg`$ split function and $`M_F`$ the factorization scale.
Our cross sections will be convoluted with parton distributions evolved via two-loop split functions. In $`n`$ dimensions the split functions have the form
$$[\mathrm{\Delta }]P_{ba}^n(x,\epsilon )=[\mathrm{\Delta }]P_{ba}(x)+\epsilon [\mathrm{\Delta }]P_{ba}^\epsilon (x)$$
(4.11)
The polarized split functions have been determined $`[`$13,14$`]`$ in the t’Hooft-Veltman scheme $`[`$15$`]`$ modified so that $`\mathrm{\Delta }P_{qq}^n(x,\epsilon )=P_{qq}^n(x,\epsilon )`$. In this scheme
$$\mathrm{\Delta }P_{gg}^\epsilon (x)=4N_C(1x)+\frac{1}{6}N_{lf}\delta (1x)$$
(4.12)
However, our calculations were carried in dimensional reduction (Sect. I), where
$$\mathrm{\Delta }P_{ab}^\epsilon (x)=P_{ab}^\epsilon (x)=0$$
(4.13)
Thus, a conversion term $`\mathrm{\Delta }d\sigma _{\mathrm{conv}}/dvdw`$ should be added to our $`\mathrm{\Delta }d\sigma ^{\gamma g}/dvdw`$. Conversion terms are determined from the difference of $`[\mathrm{\Delta }]P_{ab}^\epsilon (x)`$ in the two schemes $`[`$16$`]`$: In the present case
$$\mathrm{\Delta }\frac{d\sigma _{\mathrm{conv}}}{dvdw}=\frac{2sv}{1vw}F(0)\mathrm{\Delta }P_{gg}^\epsilon (x_2)\mathrm{\Delta }B(x_2s,t,x_2u)$$
(4.14)
with $`\mathrm{\Delta }P_{gg}^\epsilon (x)`$ given by (4.11).
The unpolarized parton distributions we use were evolved in the $`\overline{MS}`$ scheme, where $`P_{gg}^\epsilon (x)=0`$. Thus conversion term is not required.
V. SUBPROCESS $`\gamma qQ\overline{Q}q`$
The graphs contributing to this subprocess are shown in Fig. 2B. The squared sum of the corresponding amplitudes, after summing over spins and colors and averaging over initial colors, is given by
$$4m^2|M_{23}^{\gamma q}|^2=\frac{2}{N_C}(4\pi )^3\alpha \alpha _s^2(e_Q^2Q_1+e_q^2Q_2+e_Qe_qQ_3)$$
(5.1)
where $`e_q`$ the charge of the light quark $`q`$. $`|M_{23}^{\gamma \overline{q}}|^2`$ corresponding to $`\gamma \overline{q}Q\overline{Q}\overline{q}`$ is given by the same expression with an opposite sign of the last term. The quantity $`Q_1`$ is given by an expansion similar to (4.3). Next we introduce
$$s_{34}p_3p_4+m^2$$
(5.2)
Then $`Q_2`$ and $`Q_3`$ are of the form:
$`Q_{2,3}=e_1+{\displaystyle \frac{e_4}{p_1p_4}}+e_8{\displaystyle \frac{p_2k}{p_1p_4}}+\stackrel{~}{f}_4{\displaystyle \frac{p_3k}{p_2k}}+{\displaystyle \frac{\stackrel{~}{f}_5}{p_2k}}+{\displaystyle \frac{\stackrel{~}{f}_6}{p_1p_4p_2k}}+f_{12}{\displaystyle \frac{p_1k}{s_{34}}}+f_{13}{\displaystyle \frac{p_1k}{s_{34}^2}}`$ (5.3)
$`+`$ $`{\displaystyle \frac{f_{14}}{s_{34}}}+{\displaystyle \frac{\stackrel{~}{f}_{15}}{p_1k}}+{\displaystyle \frac{\stackrel{~}{f}_{16}}{s_{34}p_1k}}+{\displaystyle \frac{\stackrel{~}{f}_{17}}{s_{34}^2p_1k}}+{\displaystyle \frac{f_{18}}{s_{34}p_1p_4}}+{\displaystyle \frac{\stackrel{~}{f}_{19}}{p_1kp_2k}}+\stackrel{~}{f}_{20}{\displaystyle \frac{p_3k}{p_1k}}`$
$`+`$ $`{\displaystyle \frac{\stackrel{~}{f}_{21}}{s_{34}p_2k}}`$
As before, $`\mathrm{\Delta }|M_{23}^{\gamma g}|^2`$ is given by (5.1) and (5.3) with $`\mathrm{\Delta }Q_r`$, r=1,2,3, $`\mathrm{\Delta }e_i`$ and $`\mathrm{\Delta }f_j`$ replacing $`Q_r`$, $`e_i`$ and $`f_j`$. The coefficients $`[\mathrm{\Delta }]e_i`$ and $`[\mathrm{\Delta }]f_j`$ are given in the last part of App. B.
The contribution to $`[\mathrm{\Delta }]d\sigma /dvdw`$ is obtained by working as in Sect. 4 (c.m. frame of $`\overline{Q}(Q)`$ and final light quark). Again the terms with coefficients $`[\mathrm{\Delta }]\stackrel{~}{e}_i`$ and $`[\mathrm{\Delta }]\stackrel{~}{f}_j`$ must be integrated in $`n`$ dimensions.
After phase space integrations, we get the following results for the sets $`Q_1,Q_2`$ and $`Q_3`$:
$`{\displaystyle \frac{d\sigma _{\mathrm{Br}}^{\gamma q,Q_1}}{dvdw}}`$ $`=`$ $`Le_Q^2\{e_1+{\displaystyle \frac{4S_2}{m^2(s+t)^2}}e_3+{\displaystyle \frac{2S_2}{s_2(s+t)}}\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}e_4+e_8I_{10}+\stackrel{~}{f}_4F_4^c+\stackrel{~}{f}_6F_6^c`$ (5.4)
$`+`$ $`\stackrel{~}{f}_8F_8^c+\stackrel{~}{f}_{10}F_{10}^c+\stackrel{~}{f}_{11}F_{11}^c\}`$
$$\frac{d\sigma _{\mathrm{Br}}^{\gamma q,Q_2}}{dvdw}=Le_q^2\{e_1+f_{12}F_{12}+f_{13}F_{13}+f_{14}F_{14}+\stackrel{~}{f}_{16}F_{16}^c+\stackrel{~}{f}_{17}F_{17}^c+\stackrel{~}{f}_{20}F_{20}^c\}$$
$`{\displaystyle \frac{d\sigma _{\mathrm{Br}}^{\gamma q,Q_3}}{dvdw}}`$ $`=`$ $`Le_Qe_q\{e_1+{\displaystyle \frac{2S_2}{s_2(s+t)}}\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}e_4+e_8I_{10}+\stackrel{~}{f}_4F_4^c+\stackrel{~}{f}_6F_6^c+f_{12}F_{12}`$ (5.5)
$`+`$ $`f_{14}F_{14}+\stackrel{~}{f}_{16}F_{16}^c+f_{18}F_{18}+\stackrel{~}{f}_{19}F_{19}^c+\stackrel{~}{f}_{20}F_{20}^c+\stackrel{~}{f}_{21}F_{21}^c\}`$
with
$$L=\alpha \alpha _s^2\frac{1}{N_C}\frac{v}{8\pi }\frac{\stackrel{~}{s}_2}{S_2}$$
We do not write down expressions containing $`1/\epsilon `$ poles coming from sets $`Q_1`$ and $`Q_2`$ as they are equal with an opposite sign to the corresponding counterterms with $`(S_2m^2/s_2^2)^\epsilon `$ instead of $`(m^2/M_F^2)^\epsilon `$ (see below).
The singularities arise when the final light quark is collinear with the initial one ($`kp_2=0`$, Fig. 2B, graphs (a), (b)) as well as when the photon is collinear with the light quark ($`kp_1=0`$, Fig. 2B, graphs (c), (d)). To eliminate them we introduce two counterterms. In the second case the counterterm involves the Born cross section for $`\stackrel{}{q}\stackrel{}{\overline{q}}Q\overline{Q}`$, which is proportional to $`[`$17$`]`$:
$$\mathrm{\Delta }B_{q\overline{q}}(s,t,u)=B_{q\overline{q}}(s,t,u)=\frac{1}{s}\left(\frac{t^2+u^2}{s^2}+2\frac{m^2}{s}\right)$$
(5.6)
Moreover, in $`n`$ dimensions, in the t’Hooft-Veltman scheme:
$$\mathrm{\Delta }P_{gq}^n(x,\epsilon )=C_F\{2x+2\epsilon (1x)\}P_{gq}^n(x,\epsilon )=C_F\{\frac{1+(1x)^2}{x}\epsilon x\}$$
and
$$\mathrm{\Delta }P_{q\gamma }^n(x,\epsilon )=x\frac{1}{2}\epsilon (1x)P_{q\gamma }^n(x,\epsilon )=\frac{x^2+(1x)^2}{2}\epsilon x(1x)$$
Thus, in the $`\overline{MS}`$ scheme, the first counterterm gives:
$$[\mathrm{\Delta }]\frac{d\sigma _{\mathrm{ct}}^{(1)}}{dvdw}=\frac{1}{\epsilon }\frac{2sv}{1vw}F(\epsilon )[\mathrm{\Delta }]P_{gq}(x_2)[\mathrm{\Delta }]B(x_2s,t,x_2u)\left(\frac{m^2}{M_F^2}\right)^\epsilon ,$$
(5.7)
where $`F(\epsilon )`$, $`x_2`$ and $`[\mathrm{\Delta }]B`$ given by (3.6), (4.8) and (2.5), and the second gives:
$$[\mathrm{\Delta }]\frac{d\sigma _{\mathrm{ct}}^{(2)}}{dvdw}=\frac{1}{\epsilon }\frac{16s}{9}F(\epsilon )\frac{e_q^2}{e_Q^2}[\mathrm{\Delta }]P_{q\gamma }(w)[\mathrm{\Delta }]B_{q\overline{q}}(ws,wt,u)\left(\frac{m^2}{M_F^2}\right)^\epsilon $$
(5.8)
Although not necessary, it is now customary and even advantageous $`[`$8$`]`$ to average the unpolarized cross section over $`n2`$ spin degrees of freedom for every incoming boson. This convention is employed when fitting the unpolarized structure functions. As a result, for the unpolarized case, the r.h. side of (5.8) should be multiplied by $`(1+\epsilon )`$ and of (5.9) by $`(1\epsilon )`$.
We note that, upon integration, the singular terms in $`[\mathrm{\Delta }]Q_3`$ cancel out, as they should since there is no counterterm proportional to $`e_qe_Q`$.
Conversion terms are also needed in the present case. Along the lines of the previous section:
$$[\mathrm{\Delta }]\frac{d\sigma _{\mathrm{conv}}^{(1)}}{dvdw}=\frac{2sv}{1vw}F(0)[\mathrm{\Delta }]P_{gq}^\epsilon (x_2)[\mathrm{\Delta }]B(x_2s,t,x_2u)$$
(5.9)
and
$$[\mathrm{\Delta }]\frac{d\sigma _{\mathrm{conv}}^{(2)}}{dvdw}=\frac{16s}{9}F(0)\frac{e_q^2}{e_Q^2}[\mathrm{\Delta }]P_{q\gamma }^\epsilon (w)[\mathrm{\Delta }]B_{q\overline{q}}(ws,wt,u)$$
(5.10)
Finally, we have carried our analytical calculations using REDUCE $`[`$18$`]`$ and to some extent FORM $`[`$19$`]`$.
VI. PHYSICAL CROSS SECTIONS
Here we present the necessary formulas needed for calculation of the differential and total cross sections for the physical process $`\gamma pQ+X`$. This includes derivation of physical cross sections for both components of the reaction, i.e. pointlike and resolved, the latter to leading order (LO). Note we always observe a heavy quark in the final state. Capital letters in this chapter refer to the kinematic variables of the physical process and small letters to those of the subprocess. Starting with the pointlike component, the total cross section for the reaction (1.1) can be written
$$[\mathrm{\Delta }]\sigma _{\gamma p}(S)=_{x_{min}}^1𝑑x[\mathrm{\Delta }]f_{b/p}(x,Q^2)[\mathrm{\Delta }]\widehat{\sigma }_{\gamma b}(s);$$
(6.1)
$`b`$ denotes the corresponding parton, $`f_{b/p}`$ its probability distribution and
$$s=xS,x_{min}=4m^2/S.$$
(6.2)
The total partonic cross section $`[\mathrm{\Delta }]\widehat{\sigma }_{\gamma b}(s)`$ can be calculated straightforwardly by
$$[\mathrm{\Delta }]\widehat{\sigma }_{\gamma b}(s)=_{v_{min}}^{v_{max}}𝑑v_{w_{min}}^1𝑑w[\mathrm{\Delta }]\frac{d\widehat{\sigma }_{\gamma b}(s,v,w)}{dvdw},$$
(6.3)
where
$$v_{max/min}=\frac{1}{2}(1\pm \beta ),w_{min}=\frac{m^2}{sv(1v)}.$$
(6.4)
To derive the transverse momentum differential cross section we note that the transverse momentum $`p_T`$ of a heavy quark is invariant under boosts along the beam axis; also, that our heavy quark rapidities are defined with respect to the photon. Consequently we have:
$$[\mathrm{\Delta }]\frac{d\sigma _{\gamma p}(S,p_T)}{dp_T}=_{x_{min}(p_T)}^1𝑑x[\mathrm{\Delta }]f_{b/p}(x,Q^2)[\mathrm{\Delta }]\frac{d\widehat{\sigma }_{\gamma b}(s,p_T)}{dp_T},$$
(6.5)
with
$$x_{min}(p_T)=4(p_T^2+m^2)/S,$$
(6.6)
$$[\mathrm{\Delta }]\frac{d\widehat{\sigma }_{\gamma b}(s,p_T)}{dp_T}=_{y_{min}}^{y_{max}}𝑑y\frac{2p_T}{sv}[\mathrm{\Delta }]\frac{d\widehat{\sigma }_{\gamma b}(s,v,w)}{dvdw}.$$
(6.7)
The integration limits on the c.m. rapidity $`y`$ are
$$y_{max}=y_{min}=\mathrm{ln}(\sqrt{w_m^1}+\sqrt{w_m^11}),$$
with $`w_m4(p_T^2+m^2)/s`$. Integration over rapidities in (6.7) is not well defined for ”plus” distributions (given in (4.7)) in the partonic cross section. The problem is solved with a change of variables; One needs to consider an integration contour for a heavy quark rapidity $`y`$ and split it into two parts that no overlapping (i.e. double counting) occurs. Formally one would have
$$_{y_{max}}^{y_{max}}𝑑yf(y)=_{y_0}^{y_{max}}𝑑yf(yy)+_{y_0}^{y_{max}}𝑑yf(y).$$
(6.8)
In particular, the splitting point $`y_0`$ must be the point where function $`w=f(y)`$ has a minimum. We find
$$y_0=\mathrm{ln}(2\sqrt{(p_T^2+m^2)/s})$$
and the general relation between the ”new” variable $`w`$ and the old variable $`y`$ is
$$\mathrm{e}_{1,2}^y=\frac{1}{2\sqrt{\frac{p_T^2+m^2}{s}}}\left\{1\pm \sqrt{1\frac{4(p_T^2+m^2)}{sw}}\right\}.$$
(6.9)
The correct sign in (6.9) is different in different integration regions, e.g. in the region $`[y_0,y_{max}]`$ the function $`\mathrm{e}^y`$ decreases when one goes from $`y_0`$ to $`y_{max}`$, thus minus sign in (6.9). Similarly, we find that for the first term of (6.8) the sign for $`\mathrm{e}^y`$ in (6.9) should be positive. The resulting expression for the $`p_T`$ differential cross section reads:
$$[\mathrm{\Delta }]\frac{d\widehat{\sigma }_{\gamma b}(s,p_T)}{dp_T}=_{w_m}^1\frac{dw}{w}\frac{2p_T}{s\sqrt{1w_m/w}}\left\{[\mathrm{\Delta }]\frac{d\widehat{\sigma }}{dvdw}(v=v_+)+[\mathrm{\Delta }]\frac{d\widehat{\sigma }}{dvdw}(v=v_{})\right\},$$
(6.10)
where
$$v_\pm =\frac{1}{2}(1\pm \sqrt{1w_m/x}).$$
(6.11)
However, even the expression (6.10) is not well suited for numerical integration. One notices that there is a numerically divergent (though analytically integrable) square root in the denominator. The singularity comes from the lower limit $`x_{min}(p_T)`$ of the $`x`$ integration. To avoid this minor problem one more change of variables is necessary. Instead of the old variables $`x,w`$ we introduce the new variables $`z,w^{}`$ through the relations
$$z=\sqrt{1\frac{x_{min}(p_T)}{wx}},w^{}=w.$$
(6.12)
To correctly define integration limits for the new variables one has to perform a nontrivial mapping. Finally we obtain:
$$[\mathrm{\Delta }]\frac{d\sigma _{\gamma b}(S,p_T)}{dp_T}=\frac{4p_T}{Sx_{min}(p_T)}_0^{z_m}𝑑z_{w_m^{}}^1𝑑w[\mathrm{\Delta }]F_{b/p}(x,Q^2)\left\{[\mathrm{\Delta }]\frac{d\widehat{\sigma }}{dvdw}(v_+)+[\mathrm{\Delta }]\frac{d\widehat{\sigma }}{dvdw}(v_{})\right\}.$$
(6.13)
$`F_{b/p}(x,Q^2)`$ is a momentum distribution and
$$z_m=\sqrt{1x_{min}(p_T)},w_m^{}=\frac{x_{min}(p_T)}{1z^2},x=\frac{w_m^{}}{w}.$$
(6.14)
For the rapidity $`Y`$ fixed one gets the following expression:
$$[\mathrm{\Delta }]\frac{d\sigma _{\gamma p}(S,Y)}{dY}=_{x_{min}(Y)}^1𝑑x[\mathrm{\Delta }]f_{b/p}(x,Q^2)[\mathrm{\Delta }]\frac{d\widehat{\sigma }_{\gamma b}(s,y)}{dy},$$
(6.15)
with
$$x_{min}(Y)=\mathrm{e}^Y/(\sqrt{S}/m\mathrm{e}^Y),$$
(6.16)
and
$$[\mathrm{\Delta }]\frac{d\widehat{\sigma }_{\gamma b}(s,y)}{dy}=_{w_{min}}^1\frac{2wdw}{(\mathrm{e}^y+w\mathrm{e}^y)^2}[\mathrm{\Delta }]\frac{d\widehat{\sigma }}{dvdw}$$
(6.17)
$$w_{min}=\frac{\mathrm{e}^y}{\sqrt{s}/m\mathrm{e}^y},v=\frac{1}{1+w\mathrm{e}^{2y}},y=Y+\frac{1}{2}\mathrm{ln}x.$$
Finally we turn to the resolved LO photon contributions. We define the doubly differential cross section $`d\sigma /dYdp_T`$ for the $`22`$ subprocess:
$$[\mathrm{\Delta }]\frac{d\sigma _{\gamma p}}{dYdp_T}=2p_T_{x_{1,min}}^1𝑑x_1\frac{[\mathrm{\Delta }]F_{a/\gamma }(x_1,Q^2)[\mathrm{\Delta }]F_{b/p}(x_2^0,Q^2)}{x_1\mathrm{e}^YA}[\mathrm{\Delta }]\widehat{\sigma }_{ab}(s,x_1,x_2^0),$$
(6.18)
where
$$x_2^0\frac{x_1\mathrm{e}^YA}{x_1\mathrm{e}^YA},A(\frac{p_T^2+m^2}{S})^{1/2},s=x_1x_2^0S.$$
(6.19)
The expressions (6.1) - (6.19) give all the formulas we have used.
VII. NUMERICAL RESULTS
We present results for Q=c-quark ($`m_c=1.5`$ GeV) at $`\sqrt{S_{\gamma p}}\sqrt{S}=10`$ GeV, relevant to the experiments $`[`$1$`]`$ and $`[`$2$`,(a)]`$ and $`\sqrt{S}=100`$ GeV, as well as for Q=b-quark ($`m_b=5`$ GeV) at $`\sqrt{S}=100`$ GeV; the later energy is relevant to HERA. Higher HERA energies are not considered as the cross sections become too small. The effect of changing $`m_c`$ is also considered.
We use the NLO sets of polarized parton distributions of $`[`$20$`]`$, which can be characterized in terms of the polarized gluon distribution $`\mathrm{\Delta }g\left(x\right)`$ as follows:
* Set A: $`\mathrm{\Delta }g\left(x\right)>0`$ and relatively large
* Set B: $`\mathrm{\Delta }g\left(x\right)>0`$ and small
* Set C: $`\mathrm{\Delta }g\left(x\right)`$ changing sign; $`\mathrm{\Delta }g\left(x\right)<0`$ for $`x>0.1`$.
Notice that in the presented results also the LO contribution is convoluted with NLO distributions; in this way we believe that e.g. the magnitude of K-factors more properly reflects the NLO subprocess terms. Also, we use throughout the NLO expression of $`\alpha _s\left(\mu \right)`$ with the values for the QCD scale $`\mathrm{\Lambda }`$, flavor thresholds and number of active flavors $`N_{lf}=N1`$ that match the definitions corresponding to heavy quark decoupling. We note that in $`[`$5$`]`$ the above values were taken to match the definitions for the respective parton distributions. However, we have explicitly verified that this amounts to a negligible change in the final numerical results. Note, in (2.8) we take $`M=m`$.
At this moment there is no experimental information on the polarized photon structure functions $`\mathrm{\Delta }F_{q/\gamma }`$ and $`\mathrm{\Delta }F_{g/\gamma }`$, which determine the resolved $`\gamma `$ contributions. To estimate them we have used the LO maximal and minimal saturation sets of $`[`$21$`]`$, as well as the sets of $`[`$22$`]`$, which belong to the class of the so-called asymptotic solutions. The two sets of $`[`$21$`]`$ give contributions differing little, with the maximal saturation one slightly exceeding; the results presented below correspond to this set. The largest resolved contributions come from $`[`$22$`]`$.
In Figs. 3I, II and III, at $`\sqrt{S}=10`$ and 100 GeV we present quantities related with the differential cross sections $`\mathrm{\Delta }d\sigma /dp_T`$, where $`p_T=p_{3T}`$ (Fig. 1), versus $`x_T2p_T/\sqrt{S}`$. Measurement of such cross sections at $`\sqrt{S}10`$ GeV may be carried in (a) of $`[`$2$`]`$. Here we use the renormalization and factorization scale $`\mu =M_f=\left(p_T^2+m^2\right)^{1/2}`$.
In the parts (a) of Figs. 3I, II and III we present the NLO and LO (denoted by a $``$) contributions to the physical differential cross section for sets A, B and C of $`[`$20$`]`$. For set B we also present the contribution of subprocess (1.3) and of the resolved photon.
In the parts (b) of the same Figs we present the asymmetries
$$A_{LL}(p_T)=\frac{\mathrm{\Delta }d\sigma /dp_T}{d\sigma /dp_T}$$
(7.1)
The unpolarized distributions are the most recent set, CTEQ5 $`[`$23$`]`$. In $`A_{LL}`$ the resolved $`\gamma `$ contributions have been left out since they are small and what is presently known does not permit a completely scheme independent calculation. The errors have been estimated using
$$\delta A_{LL}=\frac{1}{P_BP_T\sqrt{L\sigma ϵ}}$$
(7.2)
At $`\sqrt{S}=10`$ GeV we use the conditions of $`[`$1$`]`$ ($`P_B=80\%`$, $`P_T=25\%`$, $`L=2`$ $`fb^1`$, c-quark detection efficiency $`ϵ_c=0.014`$) and unpolarized cross section $`\sigma `$ integrated over a bin of $`x_T`$ corresponding to $`\mathrm{\Delta }p_T=0.5`$ GeV. At $`\sqrt{S}=100`$ GeV we use $`P_B=P_T=70\%`$, $`L=100`$ $`pb^1`$, $`ϵ_c=0.15`$, for b-quark $`ϵ_b=0.05`$ and $`\sigma `$ integrated over a bin corresponding to $`\mathrm{\Delta }p_T=5`$ GeV.
Figs. 3I(a) and 3II(a) show that between $`\sqrt{S}=10`$ and $`100`$ GeV the shape of the LO $`\mathrm{\Delta }d\sigma _{LO}/dp_T`$ and NLO $`\mathrm{\Delta }d\sigma /dp_T`$ varies dramatically; this also holds for the K-factor, $`K=\mathrm{\Delta }d\sigma /dp_T/\mathrm{\Delta }d\sigma _{LO}/dp_T`$.
Most impotrant is the possibility to distinguish between sets A, B and C. Fig. 3I(b) shows that at $`\sqrt{S}=10`$ GeV near $`x_T=0.3`$ one can distinguish A and C and perhaps all A, B, C. Figs. 3II(b) and 3III(b) show that at $`\sqrt{S}=100`$ GeV the best range is $`0.2x_T0.3`$; and for $`Q=c`$ one may distinguish all A, B, C, but for $`Q=b`$ only A and C.
In Figs. 4I, II and III we present rapidity distributions. Here we use $`\mu =M_f=2m`$. The presented differential cross sections are analogous to those of Figs. 3I, II and III and
$$A_{LL}(Y)=\frac{\mathrm{\Delta }d\sigma /dY}{d\sigma /dY}$$
(7.3)
The errors have been estimated using (7.2) where now the unpolarized cross sections $`\sigma `$ are integrated over a bin $`\mathrm{\Delta }Y=1`$.
Fig. 4I(b) shows that at $`\sqrt{S}=10`$ GeV the region $`1.25Y1.5`$ is the best to distinguish set C from A or B. Fig. 4II(b) shows that at $`\sqrt{S}=100`$ GeV for c-quark, $`A_{LL}(Y)`$ has become too small. At $`Y1`$ it seems one can distinguish all A, B, C, but $`\mathrm{\Delta }d\sigma /dY`$ is small for all sets (Fig. 4II(a)). Perhaps more promising is the range $`0Y1`$, where one can distinguish C from A or B. Finally Fig. 4III(b) shows that detection of b-quark is not useful due to large errors ($`ϵ_b`$ small).
Figs. 5I,II and III present integrated cross sections $`\mathrm{\Delta }\sigma `$ and the corresponding asymmetries $`A_{LL}=\mathrm{\Delta }\sigma /\sigma `$ versus the c.m. energy $`\sqrt{S}`$. The scale is again $`\mu =M=2m`$.
Comparison of Figs. 5I and 5II shows that at the two different ranges of $`\sqrt{S}`$ the changes in the shapes and signs of $`\mathrm{\Delta }\sigma `$ and $`A_{LL}`$ is again dramatic; clearly the same holds for the corresponding $`K`$-factors, $`K=\mathrm{\Delta }\sigma _{NLO}/\mathrm{\Delta }\sigma _{LO}`$.
In Fig. 5I(b) the error (at $`\sqrt{S}`$=10 GeV) is estimated using again in (7.2) the conditions of $`[`$1$`]`$. Under these conditions we conclude that sets A and C can be distinguished, but not sets A and B or B and C. The proposed SLAC experiment $`[`$2a$`]`$, which amounts to better conditions, and will give results at somewhat lower $`\sqrt{S}`$, may distinguish also B and C.
In Figs. 5II(b) and 5III(b) the errors (at $`\sqrt{S}`$=100 GeV) have been estimated using again the values of $`P_B,P_T,L,\epsilon _c`$ and $`\epsilon _b`$ stated after Eq. (7.2). For $`c`$-quark, $`A_{LL}`$ are very small due to relatively large unpolarized cross sections $`\sigma `$. For the same reason, however, the error $`\delta A_{LL}`$ is not very large, so set C can be distinguished from A or B. For $`b`$-quark, due to a combination of small $`\epsilon _b`$ and rather small $`\sigma `$, the error is very large and precludes any useful information on $`\mathrm{\Delta }g`$.
Finally, Fig. 6, for the integrated NLO cross sections $`\mathrm{\Delta }\sigma `$ and $`\sigma `$ and for the asymmetries $`A_{LL}=\mathrm{\Delta }\sigma /\sigma `$, shows the effect of changing the $`c`$-quark mass $`m_c`$ (part (a)) and the scales $`\mu ,M_f`$ (part (b)). The results refer to set B of $`[`$20$`]`$. E.g. regarding $`\mathrm{\Delta }\sigma `$, in Fig. 6(a) we define
$$R_m=\frac{\mathrm{\Delta }\sigma (m_c)\mathrm{\Delta }\sigma (1.5\mathrm{GeV})}{\mathrm{\Delta }\sigma (1.5\mathrm{GeV})},$$
(7.4)
and in Fig. 6(b), keeping $`\mu =M_f`$, we define
$$R_{SC}=\frac{\mathrm{\Delta }\sigma (\mu )\mathrm{\Delta }\sigma (2m_c)}{\mathrm{\Delta }\sigma (2m_c)};$$
(7.5)
similarly for $`\sigma `$ and $`A_{LL}`$. Fig. 6(a) shows that at the lower $`\sqrt{S}`$ the effect of changing $`m_c`$ is more pronounced.
VIII. COMPARISON WITH OTHER PUBLICATIONS
Figs. 3II(b) and 3III(b) show that at small $`x_T`$, $`A_{LL}(p_T)`$ is small; the same holds for $`A_{LL}(Y)`$ of Fig. 4II(b). This may lead one to conclude that HERA is rather useless in specifying $`\mathrm{\Delta }g`$ $`[`$5$`]`$. However, it may not be so. On the basis of Figs. 3II(b) and III(b), reconstruct events and select only those with, say, $`x_T>0.2`$, i.e. carry integrations of $`\mathrm{\Delta }d\sigma /dp_T`$ over some cut phase space. This may well enhance the resulting $`A_{LL}`$ $`[`$24$`]`$.. Of course, an estimate of the corresponding errors is required to reach a definite conclusion.
Finally, since we present analytic results for the unpolarized cross section as well, we will compare with similar results of $`[`$8$`]`$ (”soft” part, Eq. (2.24) of $`[`$8$`]`$); here the relevant part is the last three lines of Eq. (4.5). Ref. $`[`$8$`]`$ uses the phase space slicing method, which separates the soft and hard gluon parts via a cut parameter $`\mathrm{\Delta }`$. The formal relation with our approach is
$$\frac{sv}{m^2}\mathrm{\Delta };$$
(8.1)
the necessary framework to relate these two methods is developed in $`[`$25$`]`$. Now, concerning terms involving $`t`$ and $`u`$, we easily see that they are exactly the same, except that $`t`$ and $`u`$ are interchanged (c.f. our definition, Eq. (2.1) with that of $`[`$8$`]`$, Eq. (2.13)). The only difference seems to arise from the coefficient of $`\zeta (2)`$, which is $`2`$ in our case versus $`3/2`$ in $`[`$8$`]`$. Note, however, that our coefficient $`F(\epsilon )`$ in (3.6) contains $`\mathrm{\Gamma }(1+\epsilon )`$, which upon expansion in powers of $`\epsilon `$, gives a term $`\left(\epsilon ^2/2\right)\zeta (2)`$; this accounts for the difference. Since $`F(\epsilon )`$ appears both in our loop contribution (3.6) and in our Brems contribution (4.7), the overall result is unaffected. To verify our calculation we have evaluated numerically the NLO $`\overline{MS}`$ scaling functions for the partonic $`\gamma g`$ cross section, taking into account an additional ”mass” factorization term given in eq. (6.31) of $`[`$12$`]`$, and compared it to the corresponding curves of Fig. 5 of $`[`$8$`]`$. We found exact agreement. We have also explicitly verified that the sum of our non-Abelian loop contributions and the Brems ones, that are proportional to the Born contribution, equal analytically the corresponding ”virtual+soft” expression presented in $`[`$5$`]`$.
IX. CONCLUSIONS
In this paper we have presented the complete analytic results for the heavy flavor photoproduction for both, longitudinally polarized and unpolarized initial particles, in a closed form. These include the NLO contributions of the hard Brems due to the relevant partonic subprocesses (1.2) and (1.3) that are presented for the first time in analytic form. We have computed numerically various total and differential cross sections for the energy ranges of CERN, SLAC and HERA. We have discussed the possibilities to differentiate between various scenarios for the polarized gluon distribution $`\mathrm{\Delta }g`$ and have once more emphasized the way to enhance the asymmetries for HERA energies by measuring the differential cross sections with the help of certain acceptance cuts (see also our earlier Ref. $`(a)`$ of $`[`$24$`]`$ on this subject).
ACKNOWLEDGEMENTS
We thank I. Bojak for his kind collaboration in doing comparisons. Thanks are also due to G. Bunce, D. de Florian, B. Kamal and J. Körner for discussions, to W. Vogelsang for discussions and for providing us the sets of $`[`$19$`]`$, to P. Bosted for several communications, to A. Despande for useful information and remarks and to V. Spanos and G. Veropoulos for participating in part of the calculations. Z.M. would like to thank the Particle Theory group of the Institut für Physik, Universität Mainz, for hospitality, where the calculations of the final parts of this paper were carried out.
APPENDIX A
Here we present the coefficients of the loop contributions. In the following $`[\mathrm{\Delta }]A_i`$, i=1,3, are given in App. B of $`[`$6$`]`$. For $`[\mathrm{\Delta }]d\sigma _{ae}/dv`$ given in Eq. (3.2):
$`\mathrm{\Delta }\stackrel{~}{A}_1`$ $`=`$ $`\mathrm{\Delta }A_1;\mathrm{\Delta }\stackrel{~}{A}_2=4[2(7s^2/t^2+8s/t+6)m^2/u+11s^2/tu+24s/u+26t/u+`$
$`12t^2/su+st/uT2t^2/sT]m^2/T`$
$`\mathrm{\Delta }\stackrel{~}{A}_3`$ $`=`$ $`\mathrm{\Delta }A_3/2;\mathrm{\Delta }\stackrel{~}{A}_4=4[(2u/ts/u)m^2/t2s/t+2u/s]m^2/T`$
$`\stackrel{~}{A}_1`$ $`=`$ $`A_1;\stackrel{~}{A}_2=4[(24s/t2s^2/tT+2s/T+12t/T)m^4/tu+(s/t+6t/s+t/T+11)`$
$`m^2s/uT2t^2/T^2]`$
$`\stackrel{~}{A}_3`$ $`=`$ $`A_3/2;\stackrel{~}{A}_4=4[4m^4s/ut^2(2s^2/t^2s/t+t/T2)m^2/ut/T]`$ (A1)
For $`[\mathrm{\Delta }]d\stackrel{~}{\sigma }_g/dv`$ and $`[\mathrm{\Delta }]d\sigma _g/dv`$ given in Eqs. (3.5) and (3.7):
$`\mathrm{\Delta }A_{}^{}{}_{1}{}^{}`$ $`=`$ $`m^2s/t^2`$
$`\mathrm{\Delta }A_{}^{}{}_{2}{}^{}`$ $`=`$ $`4[2m^4s^2/Tut^2(8su/t^2+9u/T8t/T8t^2/sTst/T^2+2t^2u/sT^2)m^2/u]`$
$`\mathrm{\Delta }A_{}^{}{}_{3}{}^{}`$ $`=`$ $`4[m^4s/Tut2(s/t1t^2/sT)m^2/t]`$
$`A_{}^{}{}_{1}{}^{}`$ $`=`$ $`m^2s^2/t^2u1`$
$`A_{}^{}{}_{2}{}^{}`$ $`=`$ $`4[(16sT/t^2+2s^2/t^2+8u/T+t/T)m^4/uT+2(s/u2s/T)m^2/t]`$
$`A_{}^{}{}_{3}{}^{}`$ $`=`$ $`4[(4sT/t^2+s/t+2)m^4/uT+2m^2s^2/ut^2+3]`$ (A2)
For $`[\mathrm{\Delta }]d\stackrel{~}{\sigma }_h/dv`$ and $`[\mathrm{\Delta }]d\sigma _h/dv`$ given in Eqs. (3.9) and (3.10):
$`\mathrm{\Delta }B_{}^{}{}_{1}{}^{}`$ $`=`$ $`2[4(5/u+1/t)m^4/t+(5/u+7/t+4u/t^2)m^2+(4/u+3/t)(t^2+u^2)/s]`$
$`\mathrm{\Delta }B_{}^{}{}_{2}{}^{}`$ $`=`$ $`(15s^2/tu62)m^2s/tu4(s^2/tu+8)m^4/tu+13/4(s^2/tu2)`$
$`\mathrm{\Delta }B_{}^{}{}_{3}{}^{}`$ $`=`$ $`2[(3/u+4/t)(ut)m^2/t+(1/s1/t)(t^2+u^2)/u]`$
$`\mathrm{\Delta }B_{}^{}{}_{4}{}^{}`$ $`=`$ $`4[2(2/u+3/t+u/t^2)m^2(t^2/ut+9u+3u^2/t6tu/Tt^4/T^2u+t^3/T^2)/s]`$
$`\mathrm{\Delta }B_{}^{}{}_{5}{}^{}`$ $`=`$ $`8(s^2/tu3)m^2s/tu`$
$`\mathrm{\Delta }B_{}^{}{}_{6}{}^{}`$ $`=`$ $`2[2(1/T+1/U)m^2/s+2(t^2/u^2+u^2/t^2+2t/u+2u/t)/st/uUu/tT`$
$`(t^3/u2tu+u^3/t)/sTU]m^2`$
$`B_{}^{}{}_{1}{}^{}`$ $`=`$ $`2[4(1/t3/u)m^4/t(1u/t)(1/u+2/t)m^2+4t/u+2+3u/t]`$
$`B_{}^{}{}_{2}{}^{}`$ $`=`$ $`[4(29/u^26/tu+29/t^2)m^42(1/u^212/tu+1/t^2)m^2s13t/u+413u/t]/4`$
$`B_{}^{}{}_{3}{}^{}`$ $`=`$ $`2[4(1/t3/u)m^4/t+(2u/t^21/t5/u)m^22s/u+u/t]`$
$`B_{}^{}{}_{4}{}^{}`$ $`=`$ $`4[2(4/u+1/t)m^2s/t+4t/u3s/t2t^2/Tut/T2t^3/T^2u]`$ (A3)
$`B_{}^{}{}_{5}{}^{}`$ $`=`$ $`16(1/t^2+1/u^2)m^4`$
$`B_{}^{}{}_{6}{}^{}`$ $`=`$ $`4[(t/Uu^2+u/Tt^2)m^4(5s/tu1/U1/Tt^2/TUuu^2/TUt)m^2+2tu/TU]`$
APPENDIX B
In this Appendix we list the coefficients of the Brems contributions. For $`\mathrm{\Delta }G^{\gamma g}`$ given in Eq. (4.3) and $`\mathrm{\Delta }d\sigma _{\mathrm{Br}}^{\gamma g}/dvdw`$ given in Eq. (4.5), the coefficients $`\mathrm{\Delta }e_i`$ and $`\mathrm{\Delta }f_i`$ are
$`\mathrm{\Delta }e_1`$ $`=`$ $`16[4(s/s_2u1/u+1/t)m^2/t+3s/s_2u2/us/s_2ts_2/tuu/s_2t]`$
$`\mathrm{\Delta }e_2`$ $`=`$ $`8[8(s/u1+s/t)m^4/s_2t+2(2s/u2+4s/t3ss_2/t(s+u))m^2/s_2`$
$`4s^2/s_2u+4s/u2s/s_2u/s_2+u/t+2ss_2/tu+3u/(s+u)]`$
$`\mathrm{\Delta }e_3`$ $`=`$ $`0;\mathrm{\Delta }e_4=8[4(s/u+s/s_2t/s_2)m^2/t+ss_2/tu+su/s_2t2u/s_22]`$
$`\mathrm{\Delta }\stackrel{~}{e}_5`$ $`=`$ $`8[8(t/u+1)m^4/s_2t2(s^2/s_2u+s/u+4t/s_22s_2/u)m^2/t+st/s_2u+`$
$`4s/s_22+su/s_2t+2u/t]`$
$`\mathrm{\Delta }\stackrel{~}{e}_6`$ $`=`$ $`4[8(s/us_2/t+u/t)m^4/s_2+2(2st/s_2u+s/s_2s_2/u+2+2ss_2/tu`$
$`2s^2/s_2t)m^2+s(s^2+u^2)/s_2t]`$
$`\mathrm{\Delta }e_7`$ $`=`$ $`4[8(s_2/u1)m^4/t+2(2s/u+3s_2/s2s_2^2/tu+2s/t)m^2s/(s+u)`$
$`ss_2(s^2+s_2^2)/t(s+u)u];\mathrm{\Delta }e_8=16(s_2/uu/s_2)/t`$
$`\mathrm{\Delta }\stackrel{~}{e}_9`$ $`=`$ $`4[8(t/u+s/t)m^4/s_22(2s^2/s_2u2s/u+2s_2/u2+(st/s_2+2su/s_2`$
$`s_2)/(s+u))m^2st(2s/s_2u2/u+(s_2/u+u/s_2)/(s+u))]`$
$`\mathrm{\Delta }e_{10}`$ $`=`$ $`16(t/uu/t)/s_2;\mathrm{\Delta }e_{11}=16m^4s/u`$
$`\mathrm{\Delta }f_1`$ $`=`$ $`16[4(s/u1)m^2/s_2t+(3/s_2+1/t)s/u];\mathrm{\Delta }f_2=32(s+u)/ts_2u`$
$`\mathrm{\Delta }\stackrel{~}{f}_3`$ $`=`$ $`\mathrm{\Delta }f_2;\mathrm{\Delta }\stackrel{~}{f}_4=32[(2/s_2+2/ts/s_2u)m^2/t+1/t+1/s_2]`$
$`\mathrm{\Delta }\stackrel{~}{f}_5`$ $`=`$ $`8[2(s/s_2u+s/tu2u/t^2+s/t(s+u)4ss_2/t^2(s+u))m^2s^2/s_2u+`$
$`2s/u2s/s_2+22u/s_2+ss_2/tu+s_2/t+(u+2s_2u/t)/(s+u)]`$
$`\mathrm{\Delta }\stackrel{~}{f}_6`$ $`=`$ $`4[2(2s^2/s_2u+s/u+s_2/u+3s/s_2+(s+2su/s_23s_24s(s_2^2+`$
$`u^2)/tu)/(s+u))m^2+s^2/s_2+2su/s_2+2u^2/s_22s^2/t4su/t+`$
$`2s_2^2/t2u^2/t+(s_2^2+u^2)/(s+u)]`$
$`\mathrm{\Delta }\stackrel{~}{f}_7`$ $`=`$ $`0;\mathrm{\Delta }\stackrel{~}{f}_8=0`$
$`\mathrm{\Delta }\stackrel{~}{f}_9`$ $`=`$ $`4[2(2s/u2s_2^2/tu+s_2/t(ss_2/u+2su/ts_2u/t)/(s+u))m^2+s_2`$
$`2ts_2^2/tu^2/t+(s_2^2+u^2)/(s+u)]`$
$`\mathrm{\Delta }\stackrel{~}{f}_{10}`$ $`=`$ $`8(2s+t)m^2;\mathrm{\Delta }\stackrel{~}{f}_{11}=0`$ (B1)
For $`G^{\gamma g}`$ given in Eq. (4.3) and $`d\sigma _{\mathrm{Br}}^{\gamma g}/dvdw`$ given in Eq. (4.5), the coefficients $`e_i`$ and $`f_i`$ are
$`e_1`$ $`=`$ $`16[4(1/s_2u+1/s_2t1/t^2)m^23s/s_2u+2/u1/s_2s_2/tu1/t]`$
$`e_2`$ $`=`$ $`8[16(1/t+1/u)m^6/s_2t8(3t/s_2u+2/u+2/s_21/(s+u))m^4/t+`$
$`2(s^2/s_2u8s/u+4s/s_21+2u/s_23s_2/t+u/t)m^2/(s+u)+u/s_2`$
$`4st/s_2u2s/s_2+(2s_22u+2ss_2^2/tu+s_2u/t)/(s+u)]`$
$`e_3`$ $`=`$ $`16m^2`$
$`e_4`$ $`=`$ $`8[8m^4(1/u+1/s_2)2m^2(s/us/s_22s_2/u2u/s_2)+(s_2+u)(2+`$
$`2u/s_2+2s/s_2s/u+s/s_2)]/t`$
$`\stackrel{~}{e}_5`$ $`=`$ $`8[8(s/s_21)m^4/tu2m^2(s/s_2u+2/u5/s_2+3/tu/ts_2)+s^2/s_2us/u`$
$`s/s_2+4+su/ts_22s_2/t+2s^2/ts_2]`$
$`\stackrel{~}{e}_6`$ $`=`$ $`4[16m^6(1/su1/s_2u1/ts_2)+8(su(ss_2)^2(s_2u)^2)m^4/ts_2u+2m^2(t/u+`$
$`s/s_21+s/t)s(s^2+u^2)/s_2t]`$
$`e_7`$ $`=`$ $`4[16(s_2/u1)m^6/st+8(s^2+ss_2+s_2^2)m^4/tu(s+u1)+2m^2(3s/t+(ss_2/us_2`$
$`2ss_2/t)/(s+u))ss_2(s^2+s_2^2)/tu(s+u1)];e_8=16(s_2+u)^2/ts_2u`$
$`\stackrel{~}{e}_9`$ $`=`$ $`4[16(s/s_22+s_2/su/s)m^6/tu+8(s_2/u3st/s_2u2+u/s_2)m^4/(s+u)+2m^2\times `$
$`((ss_2)(1/u+1/s_2)1+2(s_2s)/(s+u))+(s_2/u2+u/s_22st/s_2u)st/(s+u)]`$
$`e_{10}`$ $`=`$ $`16(1/t1/u)(s/s_21);e_{11}=32m^6/u`$
$`f_1`$ $`=`$ $`16[8m^4/ts_2u+2(2/u+s/tu)m^2/s_2+s/u(1/t+3/s_2)2ss_2/t(s+u)u]`$
$`f_2`$ $`=`$ $`32(1/t+1/s_2)/u;\stackrel{~}{f}_3=f_2`$
$`\stackrel{~}{f}_4`$ $`=`$ $`32[4m^4/ts_2u+2(1/s_2+1/t)m^2/u1/s_2+s_2/t(s+u)]`$
$`\stackrel{~}{f}_5`$ $`=`$ $`8[8m^4(1/s_2u+1/tu+(1/t2s_2/t^2)/(s+u))4m^2(1/t+2s_2/t^2(st/s_2+`$
$`ss_2/t)/(s+u)u)t/u+2t/s_2s_2^2/tu+(s^2t/s_2u+2u2s_2u/t)/(s+u)]`$
$`\stackrel{~}{f}_6`$ $`=`$ $`4[8m^4(s^2/ts_2u1/s_24/t(s/ts_2s/uts_2u)/(s+u))+4(s^2/s_2+su/s_2`$
$`2s_2+sus_2(s+s_2)/(s+u))m^2/t(1(s_2^2+u^2)/(s+u1)t)(s+s_2+u)^2/s_2]`$
$`\stackrel{~}{f}_7`$ $`=`$ $`32(s+u)^2m^2/t^2;\stackrel{~}{f}_8=32(s+u)m^2u/t`$
$`\stackrel{~}{f}_9`$ $`=`$ $`4[8m^4(s^2/tu+s_2/u+s_2/t)+4m^2(ss_2/t+s_22s)+(s_2/t2)(s_2^2+u^2st`$ (B2)
$`ut)]/(s+u);\stackrel{~}{f}_{10}=16m^2(2m^2+u);\stackrel{~}{f}_{11}=8m^2u^2`$
Now we shall write down the coefficients $`[\mathrm{\Delta }]e_i`$ and $`[\mathrm{\Delta }]f_j`$ for the subprocess $`\gamma qQ\overline{Q}q`$.
For $`Q_1`$ we have:
$`\mathrm{\Delta }e_1`$ $`=`$ $`8(2m^2t)/t^2,\mathrm{\Delta }e_3=0,\mathrm{\Delta }e_4=8s/t,\mathrm{\Delta }e_8=0`$
$`\mathrm{\Delta }\stackrel{~}{f}_4`$ $`=`$ $`8(2m^2+t)/t^2,\mathrm{\Delta }\stackrel{~}{f}_5=4(2(s_2+s)m^2ut)/t^2`$
$`\mathrm{\Delta }\stackrel{~}{f}_6`$ $`=`$ $`2(2m^2+2s_2t),\mathrm{\Delta }\stackrel{~}{f}_7=0,\mathrm{\Delta }\stackrel{~}{f}_8=0,\mathrm{\Delta }\stackrel{~}{f}_{10}=2(2s+t)m^2,\mathrm{\Delta }\stackrel{~}{f}_{11}=0`$
$`e_1`$ $`=`$ $`8(2m^2+t)/t^2,e_3=4m^2,e_4=8(2u+s)/t,e_8=16/t`$
$`\stackrel{~}{f}_4`$ $`=`$ $`8/t,\stackrel{~}{f}_5=4(4m^4+4m^2s_2+ut)/t^2`$
$`\stackrel{~}{f}_6`$ $`=`$ $`2(4(s_2+u)m^2+2(2s+3t)s_2+8m^44s_2^22s^24st3t^2)/t`$
$`\stackrel{~}{f}_7`$ $`=`$ $`8(s_2t)^2m^2/t^2,\stackrel{~}{f}_8=8u(s_2t)m^2/t`$
$`\stackrel{~}{f}_{10}`$ $`=`$ $`4(m^2+u)m^2,\stackrel{~}{f}_{11}=2u^2m^2`$ (B3)
For $`Q_2`$:
$`\mathrm{\Delta }e_1`$ $`=`$ $`16/s,\mathrm{\Delta }f_{12}=8/s,\mathrm{\Delta }f_{13}=8m^2/s,\mathrm{\Delta }f_{14}=8(2m^2+u)/s`$
$`\mathrm{\Delta }\stackrel{~}{f}_{15}`$ $`=`$ $`8t/s,\mathrm{\Delta }\stackrel{~}{f}_{16}=2(2u(s+t)4m^2s+s^2+2ts_2)/s`$
$`\mathrm{\Delta }\stackrel{~}{f}_{17}`$ $`=`$ $`2m^2s,\mathrm{\Delta }\stackrel{~}{f}_{20}=0`$
$`e_1`$ $`=`$ $`16/s,f_{12}=8/s,f_{13}=8m^2/s,f_{14}=8(2m^2+u)/s`$
$`\stackrel{~}{f}_{15}`$ $`=`$ $`8(2m^2+u)/s,\stackrel{~}{f}_{16}=2(4m^2s4s_2^2+6s_2s+4s_2t3s^24st2t^2)/s`$
$`\stackrel{~}{f}_{17}`$ $`=`$ $`2m^2s,\stackrel{~}{f}_{20}=16/s`$ (B4)
And, finally, for $`Q_3`$:
$`\mathrm{\Delta }e_1`$ $`=`$ $`8(s_22s+t)/st,\mathrm{\Delta }e_4=8,\mathrm{\Delta }e_8=0`$
$`\mathrm{\Delta }\stackrel{~}{f}_4`$ $`=`$ $`8(s_2t)/st,\mathrm{\Delta }\stackrel{~}{f}_5=4(2(s_2t)m^2+s_2^2s_2tut)/st`$
$`\mathrm{\Delta }\stackrel{~}{f}_6`$ $`=`$ $`2(2(2s+t)m^22s_2t+t^2)u/(s+t)s,\mathrm{\Delta }f_{12}=8(s_2s)/st`$
$`\mathrm{\Delta }f_{14}`$ $`=`$ $`4(2(s_23s)m^22su+2s_2u+ts_2)/st,\mathrm{\Delta }\stackrel{~}{f}_{15}=8`$
$`\mathrm{\Delta }\stackrel{~}{f}_{16}`$ $`=`$ $`2(2(s+2t)s_22m^2ss^24st2t^2)u/(s+t)t`$
$`\mathrm{\Delta }f_{18}`$ $`=`$ $`2(2((2s+t)s_2+2s^2+2st)m^2(2u+t)s_2t)/(s+t)s`$
$`\mathrm{\Delta }\stackrel{~}{f}_{19}`$ $`=`$ $`2(2(s+t)s_22m^2ss^2+2ut)s_2/(s+t)t,\mathrm{\Delta }\stackrel{~}{f}_{20}=0,\mathrm{\Delta }\stackrel{~}{f}_{21}=0`$
$`e_1`$ $`=`$ $`8(s_22st)/st,e_4=8(2m^2+s_2+u+tu/(s+t))/s,e_8=16/s`$
$`\stackrel{~}{f}_4`$ $`=`$ $`8(s_2t)/st,\stackrel{~}{f}_5=4(2(s_2t)m^2+s_2^2s_2tut)/st`$
$`\stackrel{~}{f}_6`$ $`=`$ $`2(2(4s/t+t/(s+t))m^2+2s_2+2u+(2s_2t4s_2^2t^2)/(s+t))u/s`$
$`f_{12}`$ $`=`$ $`8(s_2s)/st,f_{14}=4(2(s_23s)m^2+2s_2u+s_2t2su)/st`$
$`\stackrel{~}{f}_{15}`$ $`=`$ $`8(2m^2+u+ss_2/(s+t))/t`$
$`\stackrel{~}{f}_{16}`$ $`=`$ $`2(2m^2s4s_2^2+6s_2s+4s_2t3s^24st2t^2)u/(s+t)t`$
$`f_{18}`$ $`=`$ $`2(2(4su/t+2ss_2t/(s+t))m^2+s_2(s_2+u((s2s_2)^2+2s_2s)/(s+t)))/s`$
$`\stackrel{~}{f}_{19}`$ $`=`$ $`2(2(s+2t)s_22m^2s4s_2^2s^22st2t^2)s_2/(s+t)t`$
$`\stackrel{~}{f}_{20}`$ $`=`$ $`16/t,\stackrel{~}{f}_{21}=16m^2s/t`$ (B5)
APPENDIX C
We give here the Brems integrals, $`F_i`$, $`i=1,\mathrm{},11`$, appearing in Eq. (4.5) and also in Eqs. (5.4) - (5.6). Define,
$$P_1=us_2s(t+2m^2),P_2=s((t+m^2)(us_2m^2s)m^2(s+u)^2)$$
(C1)
We may now write down the integrals:
$`F_1`$ $`=`$ $`{\displaystyle \frac{1}{(s+u)^2}}[P_1+{\displaystyle \frac{sS_2}{s_2}}(t+2m^2)\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}]`$
$`F_2`$ $`=`$ $`{\displaystyle \frac{1}{(s+u)^3}}[P_1({\displaystyle \frac{s_2(s+u)(t+u+2m^2)}{4S_2}}+s(t+2m^2))+P_2{\displaystyle \frac{s_2+2m^2}{2S_2}}+{\displaystyle \frac{S_2}{s_2}}`$
$`\times `$ $`({\displaystyle \frac{s^2(t+2m^2)^2}{2}}P_2{\displaystyle \frac{m^2}{S_2}})\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}]`$
$`F_3^s`$ $`=`$ $`{\displaystyle \frac{s_2S_2u^2}{2(s+u)^3}}`$
$`F_3^c`$ $`=`$ $`{\displaystyle \frac{s_2}{2S_2(s+u)}}[{\displaystyle \frac{(t+u)^2}{4}}m^2s{\displaystyle \frac{t+u+2m^2}{s+u}}P_1{\displaystyle \frac{3}{4}}{\displaystyle \frac{P_1^2}{(s+u)^2}}]`$
$`F_4^s`$ $`=`$ $`{\displaystyle \frac{S_2u}{(s+u)^2}},F_4^c={\displaystyle \frac{P_1}{(s+u)^2}},F_5^s={\displaystyle \frac{2S_2}{s_2(s+u)}},F_5^c=0`$
$`F_6^s`$ $`=`$ $`{\displaystyle \frac{4S_2}{s_2ut}},F_6^c={\displaystyle \frac{4S_2}{s_2ut}}[\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}+2\mathrm{ln}{\displaystyle \frac{ut}{ss_2+ut}}]`$
$`F_7^c`$ $`=`$ $`{\displaystyle \frac{8S_2^2}{s_2^2(s+u)^2}}(1\epsilon )`$
$`F_8^s`$ $`=`$ $`{\displaystyle \frac{16S_2^2}{s_2^2ut(s+u)}}{\displaystyle \frac{s_2}{2S_2}}{\displaystyle \frac{s+t}{u^2t^2}}(s+u)(ut2m^2s)`$
$`F_8^c`$ $`=`$ $`{\displaystyle \frac{16S_2^2}{s_2^2ut(s+u)}}[{\displaystyle \frac{s_2}{2S_2}}{\displaystyle \frac{s+t}{u^2t^2}}(s+u)(ut2m^2s)(\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}+2\mathrm{ln}{\displaystyle \frac{ut}{ss_2+ut}}){\displaystyle \frac{s_2^2}{S_2^2}}{\displaystyle \frac{2P_2}{u^2t^2}}1+\epsilon ]`$
$`F_9^s`$ $`=`$ $`{\displaystyle \frac{4S_2}{s_2^2u}},F_9^c={\displaystyle \frac{4S_2}{s_2^2u}}[\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}+\mathrm{ln}{\displaystyle \frac{u^2}{(s+u)^2}}]`$
$`F_{10}^s`$ $`=`$ $`{\displaystyle \frac{8S_2(s+u)}{s_2u^2t^2}},F_{10}^c={\displaystyle \frac{8S_2(s+u)}{s_2u^2t^2}}[\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}+2\mathrm{ln}{\displaystyle \frac{ut}{ss_2+ut}}+{\displaystyle \frac{s_2}{m^2}}{\displaystyle \frac{ut2m^2s}{ss_2+ut}}]`$
$`F_{11}^s`$ $`=`$ $`{\displaystyle \frac{32S_2^2}{s_2^2u^2t^2}}{\displaystyle \frac{s_2^2}{S_2^2u^2t^2}}[P_2+{\displaystyle \frac{S_2}{s_2}}(ss_2+ut)(ut2m^2s)]`$
$`F_{11}^c`$ $`=`$ $`{\displaystyle \frac{32S_2^2}{s_2^2u^2t^2}}[{\displaystyle \frac{s_2^2}{S_2^2u^2t^2}}(P_2+{\displaystyle \frac{S_2}{s_2}}(ss_2+ut)(ut2m^2s))(\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}+2\mathrm{ln}{\displaystyle \frac{ut}{ss_2+ut}}){\displaystyle \frac{s_2^2}{S_2^2}}{\displaystyle \frac{8P_2}{u^2t^2}}`$ (C2)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{s_2^2}{m^2S_2}}1+\epsilon ]`$
Note that parts of $`F_7^c`$, $`F_8^c`$ and $`F_{11}^c`$ proportional to $`1\epsilon `$ cancel out exactly in Eq. (4.5). The singular and finite parts of these integrals can be found in Appendix C of $`[`$12$`]`$; below we give the derivation of $`𝒪(\epsilon )`$ terms.
We use the momentum parametrizations of Appendix A of $`[`$6$`]`$ and, as in $`[`$12$`]`$, we denote
$$F_n^{(k,l)}𝑑\mathrm{\Omega }_n(a+b\mathrm{cos}\theta _1)^k(A+B\mathrm{cos}\theta _1+C\mathrm{sin}\theta _1\mathrm{cos}\theta _2)^l$$
(C3)
where
$$𝑑\mathrm{\Omega }_n_0^\pi 𝑑\theta _1\mathrm{sin}^{12\epsilon }\theta _1_0^\pi 𝑑\theta _2\mathrm{sin}^{2\epsilon }\theta _2.$$
(C4)
All the above integrals are proportional to $`1/a^2=1/\omega _k^2\omega _2^21/s_2^21/(1w)^2`$; since the $`23`$ particle phase space is proportional to $`(1w)^{12\epsilon }`$, in view of Eq. (4.6), the terms of $`𝒪(\epsilon )`$ give finite contributions proportional to $`\delta (1w)`$.
Integral $`F_7𝑑\mathrm{\Omega }_n/(p_2k)^2`$:
This is of the type $`\widehat{I}_n^{(2,0)}`$ of $`[`$12$`]`$. The result is
$$F_7=\frac{\pi }{a^2}\frac{1}{1+\epsilon }=\frac{\pi }{a^2}(1\epsilon +𝒪(\epsilon ^2))$$
(C5)
Integral $`F_8𝑑\mathrm{\Omega }_n/(p_2k)^2(p_1p_4)`$:
This is of the type $`\widehat{I}_n^{(2,1)}`$ of $`[`$12$`]`$, and determination of the $`𝒪(\epsilon )`$ term proceeds as follows. First, defining
$$HA+B\mathrm{cos}\theta _1+C\mathrm{sin}\theta _1\mathrm{cos}\theta _2,$$
(C6)
one can show the identity
$$\frac{1}{(1\mathrm{cos}\theta _1)H}=\frac{1}{A+B}\left\{\frac{1}{1\mathrm{cos}\theta _1}+\frac{B}{H}+\frac{C\mathrm{sin}\theta _1\mathrm{cos}\theta _2}{(1\mathrm{cos}\theta _1)H}\right\}$$
(C7)
and by repeated application of it:
$`{\displaystyle \frac{\mathrm{sin}\theta _1}{(1\mathrm{cos}\theta _1)^2H}}=`$ (C8)
$`{\displaystyle \frac{1}{A+B}}\{{\displaystyle \frac{\mathrm{sin}\theta _1}{(1\mathrm{cos}\theta _1)^2}}+{\displaystyle \frac{B}{A+B}}[{\displaystyle \frac{\mathrm{sin}\theta _1}{(1\mathrm{cos}\theta _1)}}+{\displaystyle \frac{B\mathrm{sin}\theta _1}{H}}{\displaystyle \frac{C(1+\mathrm{cos}\theta _1)\mathrm{cos}\theta _2}{H}}]`$
$``$ $`{\displaystyle \frac{C\mathrm{cos}\theta _2}{A+B}}[{\displaystyle \frac{(1+\mathrm{cos}\theta _1)}{(1\mathrm{cos}\theta _1)}}+{\displaystyle \frac{B(1+\mathrm{cos}\theta _1)}{H}}]+{\displaystyle \frac{C^2}{(A+B)^2}}[{\displaystyle \frac{\mathrm{sin}^3\theta _1\mathrm{cos}^2\theta _2}{(1\mathrm{cos}\theta _1)^2}}`$
$`+`$ $`{\displaystyle \frac{B\mathrm{sin}\theta _1(1+\mathrm{cos}\theta _1)\mathrm{cos}^2\theta _2}{H}}{\displaystyle \frac{C(1+\mathrm{cos}\theta _1)^2\mathrm{cos}^3\theta _2}{H}}]\}`$
Note that the fifth term vanishes due to the integration over $`\theta _2`$. Also, terms with only $`H`$ in the denominators are finite, consequently have no poles and cannot produce finite contributions from their $`𝒪(\epsilon )`$ terms. Thus, we are left with the terms 1,2 and 7. After integrating them in n-dimensions, summing and keeping the relevant order $`\epsilon `$ terms, we arrive to the following result:
$$F_8^\epsilon =\frac{\pi }{a^2(A+B)}\epsilon $$
(C9)
Integral $`F_{11}𝑑\mathrm{\Omega }_n/(p_2k)^2(p_1p_4)^2`$:
Proceeding as before, the term of $`𝒪(\epsilon )`$ is provided by
$$F_{11}^\epsilon =\frac{\pi }{a^2(A+B)^2}\epsilon $$
(C10)
Finally, we give the Brems integrals, $`F_i`$, $`i=12,\mathrm{},21`$, appearing in Eqs. (5.4) and (5.6). With
$$P_3=2m^2s+u(ss_2),Z=4m^2s(s+t)s_2^2t$$
(C11)
we have:
$`F_{12}`$ $`=`$ $`{\displaystyle \frac{1}{\overline{y}^2}}[P_1(tu)+{\displaystyle \frac{s}{2}}P_3F_{14}]`$
$`F_{13}`$ $`=`$ $`{\displaystyle \frac{1}{\overline{y}^2}}[{\displaystyle \frac{2S_2}{sm^2}}P_3+P_1(tu)F_{14}]`$
$`F_{14}`$ $`=`$ $`{\displaystyle \frac{2S_2}{s_2\overline{y}}}\mathrm{ln}\left({\displaystyle \frac{s_2(ss_2)+2m^2s+s_2\overline{y}}{s_2(ss_2)+2m^2ss_2\overline{y}}}\right)`$
$`F_{15}^s`$ $`=`$ $`{\displaystyle \frac{2S_2}{s_2(s+t)}},F_{15}^c=0`$
$`F_{16}^s`$ $`=`$ $`{\displaystyle \frac{4S_2}{s_2su}},F_{16}^c={\displaystyle \frac{4S_2}{s_2su}}[\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}+\mathrm{ln}{\displaystyle \frac{u^2}{(s+t)^2}}]`$
$`F_{17}^s`$ $`=`$ $`{\displaystyle \frac{8S_2(s+t)}{s_2s^2u^2}},F_{17}^c={\displaystyle \frac{8S_2(s+t)}{s_2s^2u^2}}[\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}+\mathrm{ln}{\displaystyle \frac{u^2}{(s+t)^2}}{\displaystyle \frac{s_2P_3}{m^2s(s+t)}}]`$
$`F_{18}`$ $`=`$ $`{\displaystyle \frac{4S_2}{s_2\sqrt{tz}}}\mathrm{ln}\left({\displaystyle \frac{Z2m^2s(s+t)+s_2\sqrt{tz}}{Z2m^2s(s+t)s_2\sqrt{tz}}}\right)`$
$`F_{19}^s`$ $`=`$ $`{\displaystyle \frac{8S_2}{s_2^2s}},F_{19}^c={\displaystyle \frac{8S_2}{s_2^2s}}\mathrm{ln}{\displaystyle \frac{sS_2}{ss_2+ut}}`$
$`F_{20}^s`$ $`=`$ $`{\displaystyle \frac{S_2t}{(s+t)^2}},F_{20}^c={\displaystyle \frac{P_1(tu)}{(s+t)^2}}`$
$`F_{21}^s`$ $`=`$ $`{\displaystyle \frac{4S_2}{s_2st}},F_{21}^c={\displaystyle \frac{4S_2}{s_2st}}[\mathrm{ln}{\displaystyle \frac{S_2}{m^2}}+\mathrm{ln}{\displaystyle \frac{t^2}{(s+u)^2}}]`$ (C12)
REFERENCES
* G. Baum et al, COMPASS Collaboration: CERN/SPLC 96-14 and 96-30.
* (a) R. Arnold, P. Bosted et al, SLAC-PROPOSAL-E156, 1997; (b) W.-D. Nowak, DESY 96-095; (c) A. de Roeck and T. Gehrmann, DESY-Proceedings-1998-1.
* M. Gluck and E. Reya, Z. Phys. C39, 569 (1988); B. Lampe and E. Reya, Phys. reports (in press). The latter is also a comprehensive review of polarized particle processes.
* M. Stratmann and W. Vogelsang, Z. Phys. C74, 641 (1997); A. Watson, ibid C12, 123 (1982).
* I. Bojak and M. Stratmanm: (a) Phys. Lett. B433, 411 (1998); (b) Nucl. Phys. B540, 345 (1999); Erratum: ibid B569, 694 (2000).
* B. Kamal, Z. Merebashvili and A.P. Contogouris, Phys. Rev. D51, 4808 (1995); Erratum: ibid D55, 3229 (1997).
* G. Jikia and A. Tkabladze, ibid D54, 2030 (1996).
* J. Smith and W.L. van Neerven, Nucl. Phys. B374, 36 (1992).
* P. Nason et al, ibid B327, 49 (1989).
* T. Muta, ”Foundations of Quantum Chromodynamics” (Word Scientific, 1987)
* G. Passarino and M. Veltman, Nucl. Phys. B160, 151 (1979)
* W. Beenakker et al, Phys. Rev. D40, 54 (1989).
* R. Mertig and W. van Neerven, Z. Phys. C70, 637 (1996).
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* G.’t Hooft and M. Veltman, Nucl. Phys. B44, 189 (1972).
* B. Kamal, Phys. Rev. D53, 1142 (1996); see also A.P. Contogouris and Z. Merebashvili, ibid D55, 2718 (1997).
* A.P. Contogouris, S. Papadopoulos and B. Kamal, Phys. Lett. B246, 523, (1990).
* A. Hearn, REDUCE User’ s Manual Version 3.6 (Rand Corporation, Santa Monica, CA, 1995).
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* M. Gluck and W. Vogelsang, Z. Phys. C55, 353 (1992) and C57, 309 (1993); M. Gluck, M. Stratmann and W. Vogelsang, Phys. Lett. B187, 373 (1994).
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* H. Lai et al, Eur. Phys. J. C12 (2000) 375.
* (a) Talk given by Z. Merebashvili at the International Workshop Spin 99, Prague, Czech Republic, 5-11 Sept. 1999, published in Czech. J. Phys, Vol. 50, No. S1, 153 (2000); and hep-ph/9911506; (b) A.P. Contogouris, Z. Merebashvili and G. Grispos, Phys. Lett. B482, 93 (2000).
* B. Kamal and Z. Merebashvili, Phys. Rev. D58, 074005 (1998).
FIGURE CAPTIONS
* LO and loop graphs. In the loop graphs $`p_1p_2`$ crossed ones are not shown. Note that graph $`(i)`$, representing gluon, quark and ghost loop, does not contribute here.
* A) Gluon Brems graphs; $`p_1p_2`$ crossed ones are not shown. B) Graphs of the subprocess $`\gamma qQ\overline{Q}q`$.
* Quantities related with the $`p_T`$ distributions versus $`x_T=2p_T/\sqrt{S}`$: Parts (a): Polarized differential cross sections; the LO (Born) ones are indicated by $``$. Parts (b): Asymmetries for sets A, B and C. 3I: $`Q=c`$, $`\sqrt{S}=10`$ GeV. 3II: $`Q=c`$, $`\sqrt{S}=100`$ GeV. 3III: $`Q=b`$, $`\sqrt{S}=100`$ GeV.
* Quantities related with the rapidity $`Y`$ distributions: Parts (a) and (b), as well as 4I, 4II and 4III as in Fig. 3.
* Quantities related with the integrated cross sections for $`\stackrel{}{\gamma }\stackrel{}{p}Q+X`$: (a) Factors $`K=\mathrm{\Delta }\sigma /\mathrm{\Delta }\sigma _{LO}`$ (b) Asymmetries.
* At c.m. energies $`\sqrt{S}=10`$ and 100 GeV, for integrated cross sections and with solid lines for $`\mathrm{\Delta }\sigma `$, dashed for $`\sigma `$ and dotted for the asymmetry $`A_{LL}`$: a) The ratio $`R_m`$ (see end of Sect. VI) with $`m=m_c`$; b) The fractional variation $`R_{sc}`$ with the scale $`\mu =M_f`$ and with respect to $`\mu =M_f=3`$ GeV. For both a) and b) the lines specified by + refer to the corresponding quantities for $`\sqrt{S}=100`$ GeV.
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# Search for periodic gravitational radiation with the ALLEGRO gravitational wave detector
## I Introduction
The majority of experimental searches for gravitational radiation have focused on the detection of burst signals, such as those emitted from the collapse of a massive star. There are compelling arguments that nearby millisecond pulsars can provide a detectable source of continuous (CW) gravity waves . We will use the term “pulsar” in the following text to refer to a rotating neutron star, regardless of whether it is emitting detectable electromagnetic radiation.
It is well known that a symmetric object rotating about its symmetry axis does not emit gravitational radiation. Therefore, if a pulsar is to radiate gravity waves, it must either be asymmetric or precessing. For the purposes of this search, we have considered only asymmetric pulsars. The asymmetry can be rotationally induced or due to a “star quake” deforming the neutron star crust or due to some other process. By any mechanism, the amount of ellipticity produced in the neutron star is expected to be small, from a maximum of $`10^5`$, which, for some pulsar models is the maximum supportable asymmetry of the neutron star crust to less than $`10^9`$. The latter number is obtained by assuming the measured spin-down of certain millisecond pulsars is due entirely to the emission of gravitational radiation.
In this article we report the results of our first search for a CW signal in data taken by the ALLEGRO resonant gravitational wave detector at Louisiana State University during the first three months of 1994 . Another search has been performed by the EXPLORER detector team from the University of Rome, which has a similar detector and is pursuing a different type of analysis .
If there were known pulsars with the right spin period, they would be the natural target for any CW search. But there are only a small number of known pulsars with spin periods small enough (2.2 milliseconds) to match our antenna’s 900 Hz frequency and none of those listed in the 1995 Taylor pulsar catalog has a spin period that matches our two narrow reception bands. We would ideally like to make an all sky search, but this is computationally very intensive. Given limited computing resources it nearly impossible to exhaust all of the obvious possibilities. All-sky search strategies have been proposed elsewhere .
Therefore we have chosen another strategy and directed our first search towards two regions where one might reasonably expect to find a high density of CW sources. The first is the globular cluster 47 Tucanae (Tuc), and the second is the center of the galaxy. The celestial coordinates for both are listed in Table I. We chose 47 Tuc because it is relatively close, and because of the large number of fast millisecond pulsars known to be concentrated in this globular cluster , which suggests the possibility of more undetected neutron stars with short spin periods.
The analysis assumes the simplest possible CW source, one that has a stable emission frequency in its rest frame. In other words it is not spinning down and has no companions to cause orbital doppler shifts. In this case only knowledge of the earth’s motion and the antenna’s reception pattern are needed to determine the modulation of the frequency and amplitude of the received signal. These modulations are highly dependent on the source location. At the frequency resolution of our search, if the actual source location is offset from the supposed location by as little as 1.00 in right ascension and 0.25 in declination, the gravitational wave frequency will be mis-identified by 1 bin . However, the range of possible locations for detection of the signal (albeit with the frequency mis-identified) is significantly larger .
## II ALLEGRO
The ALLEGRO gravitational wave detector is located in the Physics Building at Louisiana State University in Baton Rouge, Louisiana ($`30^{}.413`$ N lat, $`91^{}.179`$ W long). ALLEGRO consists of a resonant aluminum bar equipped with a resonant inductive transducer and a dc SQUID amplifier all cooled to 4.2 K. The bar is 60 cm in diameter and 300 cm in length, with a physical mass of 2296 kg. The bar is oriented perpendicular to the plane of the great circle on the Earth that passes through Geneva, the location of the Rome Explorer antenna, and midway between Baton Rouge, LA and Stanford, CA. This orientation results in the axis of ALLEGRO being directed along a line $`40^{}.4`$ west of North.
Vibrations of our resonant antenna produce a voltage out of the SQUID electronics which is proportional to the relative displacement of the bar and transducer. The coupled system of bar and transducer has two normal modes at roughly 896.8 Hz and 920.3 Hz. These are referred to as the “minus” and “plus” resonant modes respectively. ALLEGRO is most sensitive in a small bandwidth around each of these modes. The voltage out of the SQUID electronics is sent to a single lock-in detector which demodulates and low pass filters the signal. The reference frequency of the lock-in is set halfway between the normal mode frequencies of the antenna, thus shifting the frequency of the signal from near 1 kHz to near 10 Hz. The output consists of in-phase ($`x`$) and quadrature ($`y`$) components which contain the full amplitude and phase information. The demodulated data is then sampled 125 times a second and written to disk.
## III target source
We describe the gravity wave in a reference frame which has its $`\widehat{e}_z^{^{}}`$ axis aligned with the direction of propagation of the wave. This frame is referred to as the “wave frame” and is denoted by primed coordinates ($`x^{^{}},y^{^{}},z^{^{}}`$). The most general form of the gravity wave in the wave frame is
$$𝐡_{wave}=\left[\begin{array}{ccc}h_+(t^{^{}})& h_\times (t^{^{}})& \hfill 0\\ h_\times (t^{^{}})& h_+(t^{^{}})& \hfill 0\\ 0& 0& \hfill 0\end{array}\right]$$
(1)
where $`h_+(t^{^{}})`$ and $`h_\times (t^{^{}})`$ are the amplitudes of the two allowed states of linear polarization, referred to as the plus and cross amplitudes respectively.
In General Relativity one can compute the polarization amplitudes for a rotating, asymmetric neutron star using the quadrupole approximation . We choose for our model a neutron star with the three principle moments of inertia about the three principle axes fixed in the body frame, denoted by $`I_1,I_2\text{and}I_3`$. The rotation axis is chosen to be along $`I_3`$ and the neutron star is assumed to be deformed such that $`I_1I_2`$ . The ellipticity of the pulsar is then defined to be
$$ϵ=\frac{I_1I_2}{I_3}.$$
(2)
The resulting gravitational radiation is emitted at twice the pulsar rotation frequency and the two polarization amplitudes are given by
$`h_+(t^{^{}})`$ $`=`$ $`h_c(1+\mathrm{cos}^2i)\mathrm{cos}(2\pi f_0t^{^{}})`$ (3)
$`h_\times (t^{^{}})`$ $`=`$ $`2h_c\mathrm{cos}i\mathrm{sin}(2\pi f_0t^{^{}})`$ (4)
where the angle $`i`$ is measured between the pulsar spin axis and the line of sight to the detector, $`h_c`$ is the ”characteristic amplitude” of the incident wave, and $`f_0=2f_{rot}`$ is the frequency of the gravitational wave, equal to twice the rotation frequency of the pulsar. This is not the most general pulsar model one could construct. If the neutron star deformation is not along a principle axis, then emission occurs at the rotation frequency and twice the rotation frequency . If the star is precessing, emission occurs at the rotation frequency, and the rotation frequency plus or minus the precession frequency . For this work, we consider only emission at twice the rotation frequency.
The characteristic amplitudeNew’s characteristic strain is greater by a factor of 2 as they assumed that the rotation axis of the pulsar was along the line of sight to the Earth. is given by
$$h_c=\frac{2G}{c^4r}ϵI_3(\pi f_0)^2$$
(5)
where $`r`$ is the distance to the source, $`G`$ is the gravitational constant, and $`c`$ the speed of light. Adopting a value of $`I_3=10^{45}\mathrm{g}\mathrm{cm}^2`$ (a reasonable estimate for a $`1.4M_{}`$ pulsar with a radius of 10 km), this can be written as
$$h_c=5.2\times 10^{28}(\frac{I_3}{10^{45}\mathrm{g}\mathrm{cm}^2})(\frac{ϵ}{10^8})(\frac{10\mathrm{kpc}}{r})(\frac{f_0}{1\mathrm{kHz}})^2.$$
(6)
For galactic sources ($`<`$ 10 kpc), and assuming the maximum allowed pulsar ellipticity ($`ϵ10^5`$) we see that the characteristic amplitude is of order $`10^{24}`$. As will be shown, this is of the same level as the measurements.
## IV phase considerations
Reliable detection of a periodic gravitational wave signal depends on tracking the phase of the signal through many cycles. As given in Eq. 3, the polarization amplitudes of the gravity wave are pure sinusoids at an unknown frequency. Detection of such a signal would involve a single Fourier transform of the data, the resulting spectrum then being scanned for anomalous peaks. However, the phase of the gravity wave as observed at the detector at a particular observation time, relative to the emitted phase, depends on a number of factors: 1) intrinsic pulsar spin-down, 2) motion of the Earth within the solar system, 3) if it is part of a binary system, orbital motion of the pulsar, and 4) proper motion of the pulsar. This results in the initially narrow-band signal being spread out over many frequency bins, decreasing the signal to noise ratio of the single Fourier transform. As the expected signals are already weak, this seriously compromises one’s detection capabilities. For the purposes of this work, we assume that the pulsar is solitary with negligible spin-down, and has no other accelerations with respect to the solar system barycenter (SSB). We also assume the proper motion of the pulsar is small such that it does not move significantly during the observation time.
It should be noted that the spindown due to the emission of gravitational radiation at the level of sensitivity of this experiment would be quite significant. The lack of an actual target source has led us to choose the simplest case for the purposes of demonstration of the potential sensitvity of a real detector. With these assumptions, the phase of the carrier frequency is modulated only by motion of the Earth in the solar system, with the modulation being given by (see Fig. 1)
$$\mathrm{\Phi }(t)=\omega _0(\frac{1}{c}𝐫(t)\widehat{𝐧}+1.658^{ms}\mathrm{sin}g(t))$$
(7)
where $`𝐫(t)`$ is the vector from the SSB to the center of mass of the detector at observation time $`t`$, $`\widehat{𝐧}`$ is a unit vector from the SSB directed towards the pulsar, $`c`$ is the speed of light, and $`\omega _0=2\pi f_0`$. The second term is due to the combined effect of time dilation and gravitational redshift due to the solar system bodies. It is periodic over a year with maximum delay of 1.7 ms. We note that the angle $`g(t)`$ varies slightly from year to year. In 1994 its value was
$`g(t)/2\pi =356.60+0.98560\times D`$
where $`D`$ is the day of the year. Including this information into Eq. 3, we write the two polarizations of the gravity wave as (expressed in a frame parallel to the wave frame)
$`h_+(t)`$ $`=`$ $`h_c(1+\mathrm{cos}^2i)\mathrm{cos}(\omega _0t+\mathrm{\Phi }(t)+\mathrm{\Phi }_0)`$ (8)
$`h_\times (t)`$ $`=`$ $`h_c(2\mathrm{cos}i)\mathrm{sin}(\omega _0t+\mathrm{\Phi }(t)+\mathrm{\Phi }_0)`$ (9)
where $`\mathrm{\Phi }_0`$ is an unknown, constant offset.
Astronomical locations are usually given as right ascention and declination coordinates, denoted $`\alpha ,\delta `$ respectively. These are defined in the the “celestial” coordinate frame (CC frame with coordinates $`(X,Y,Z)`$). This frame is centered at the solar system barycenter (SSB) with $`\widehat{e}_Z`$ along the Earth’s rotation axis. $`\widehat{e}_X`$ and $`\widehat{e}_Y`$ are in the Earth’s equatorial plane with $`\widehat{e}_X`$ directed towards the intersection of the equatorial plane with the Earth’s orbital plane (the ecliptic) at the vernal equinox. Right ascension is measured in hours of angle (12 hrs = $`\pi `$ radians) from the vernal equinox eastward along the celestial equator to the celestial object and declination is measured in degrees north (+) or south (-) of the equatorial plane.
To calculate the dot product in the phase modulation term, we also expressed the detector position in the CC frame. This was done in two steps. First, the position of the center of mass of the Earth relative to the SSB was obtained using a commercially available software package from the U.S. Naval Observatory (MICAMultiyear Interactive Computer Almanac, U.S. Naval Observatory, 3450 Mass. Ave., N.W. Washington, DC 20392). MICA uses the Jet Propulsion Laboratory DE200/LE200 ephemeris. It reports positions of the Earth’s center of mass in a cartesian coordinate frame centered at the SSB. Coordinates are given to the nearest $`10^9`$ astronomical units (AU), which is of order $`10^2`$ meters. This corresponds to 1/1000 of a wavelength for the signals of interest here, enabling an accurate tracking of the phase of the gravitational wave.
Second, a GPS receiver was used to obtain the latitude and longitude coordinates of ALLEGRO, again with sufficient accuracy to track the signal. These were converted to cartesian coordinates in a frame centered at the Earth’s center of mass and vectorally added to the center of mass positions to provide ALLEGRO’s position with respect to the SSB. The cartesian coordinates for ALLEGRO were converted to the spherical coordinates ($`r,t_s,\delta _A`$): $`r`$ being the distance from the SSB to the detector, $`t_s`$ the local sidereal time, and $`\delta _A`$ the declination of the detector. Using the spherical trigonometric formula for the cosine of the angle between two vectors, we write the phase delay as
$$\mathrm{\Phi }(t)=2\pi f_0\left(\frac{r(t)}{c}[\mathrm{sin}\delta _A\mathrm{sin}\delta +\mathrm{cos}\delta _A\mathrm{cos}\delta \mathrm{cos}\gamma (t)]+1.658^{ms}\mathrm{sin}g(t)\right)$$
(10)
where $`\gamma =t_s\alpha `$ is the local hour angle.
## V signal considerations
A passing gravity wave provides the largest fractional change in the length of a bar, and therefore the largest signal in a bar detector, when its direction of propagation is perpendicular to the bar axis and one of the polarizations of the gravity wave lies along the bar axis. In general, the gravity wave is incident to the bar with some arbitrary angle and polarization. We define the polarization angle as the angle between the bar axis (west of North) and the direction of the plus polarization of the gravity wave.
We describe the detector in a coordinate frame whose origin is at the center of mass of the antenna, with the $`\widehat{e}_z`$ axis aligned with the local vertical and the $`\widehat{e}_x`$ axis aligned with the bar axis. This frame is referred to as the “lab frame” and is denoted denoted by unprimed coordinates ($`x,y,z`$).
For a gravity wave incident with an arbitrary orientation to the bar axis, it is only the component of the strain tensor along the bar axis which produces a detectable driving force on the bar. This force is commonly written as
$$(t)=\frac{1}{2}\mu L_e\ddot{h}_{xx}(t)$$
(11)
where $`\ddot{h}_{xx}(t)`$ is the second time derivative of the strain component along the bar axis. The quantities $`\mu `$ and $`L_e`$ are the effective mass and length of the first longitudinal eigenmode of the bar, obtained by solving the elastic equations of motion for a long, thin cylinder. The effective mass is equal to half the physical mass of the bar ($`\mu =1148`$ kg) and $`L_e=(4/\pi ^2)`$L where L=3 m is the actual length of the bar. The force acting on the bar dominates the output as the force on the transducer itself produces a much smaller motion.
To calculate the component of the gravitational wave, given by Eq. 1 in the wave frame, along the bar axis requires knowledge of the rotation matrix between the wave frame and the lab frame. As both source direction and detector locations are known, the rotation matrix is completely specified (up to an angle related to the unknown polarization of the gravity wave). We calculate the rotation matrix by first rotating the wave frame to the CC frame, and then the CC frame to the lab frame. The full rotation matrix from the wave frame to the lab frame is given by
$$𝐑_{wavelab}=𝐑_{CClab}𝐑_{waveCC}.$$
(12)
The CC frame is related to the wave frame by the angles $`(\varphi =\alpha \pi /2,\theta =\delta +\pi /2,\psi =\psi _0)`$, where we have used the Euler x-convention to define the axes of rotation. Using Eq. 4-46 of Goldstein , with the stated angular substitutions, the rotation matrix from the wave frame to the CC frame is given by
$`𝐑_{waveCC}=\left(\begin{array}{ccc}(\mathrm{sin}\alpha \mathrm{cos}\psi _0\mathrm{cos}\alpha \mathrm{sin}\delta \mathrm{sin}\psi _0)& (\mathrm{sin}\psi _0\mathrm{sin}\alpha \mathrm{cos}\psi _0\mathrm{cos}\alpha \mathrm{sin}\delta )& (\mathrm{cos}\alpha \mathrm{cos}\delta )\\ (\mathrm{cos}\psi _0\mathrm{cos}\alpha \mathrm{sin}\psi _0\mathrm{sin}\alpha \mathrm{sin}\delta )& (\mathrm{sin}\psi _0\mathrm{cos}\alpha \mathrm{cos}\psi _0\mathrm{sin}\alpha \mathrm{sin}\delta )& (\mathrm{sin}\alpha \mathrm{cos}\delta )\\ (\mathrm{sin}\psi _0\mathrm{cos}\delta )& (\mathrm{cos}\psi _0\mathrm{cos}\delta )& (\mathrm{sin}\delta )\end{array}\right).`$
Again using the x-convention to define the rotation axes, the CC frame is related to the laboratory frame by the angles $`(\varphi =t_s+\pi /2,\theta =\pi /2l,\psi =\pi /2+\eta )`$. $`t_s`$ is the detector local sidereal time as before and $`l`$ is the detector latitude. The particular value for the final rotation angle comes from the choice to define the $`\widehat{e}_x`$ axis pointing west of North along the bar axis. The rotation matrix from CC frame to lab frame is then
$`𝐑_{CClab}=\left(\begin{array}{ccc}(\mathrm{sin}\eta \mathrm{sin}t_s\mathrm{sin}l\mathrm{cos}t_s\mathrm{cos}\eta )& (\mathrm{sin}\eta \mathrm{cos}t_s\mathrm{sin}l\mathrm{sin}t_s\mathrm{cos}\eta )& (\mathrm{cos}\eta \mathrm{cos}l)\\ (\mathrm{cos}\eta \mathrm{sin}t_s+\mathrm{sin}l\mathrm{cos}t_s\mathrm{sin}\eta )& (\mathrm{cos}\eta \mathrm{cos}t_s+\mathrm{sin}l\mathrm{sin}t_s\mathrm{sin}\eta )& (\mathrm{sin}\eta \mathrm{cos}l)\\ (\mathrm{cos}l\mathrm{cos}t_s)& (\mathrm{cos}l\mathrm{sin}t_s)& (\mathrm{sin}l)\end{array}\right).`$
Using Eq. 12, the full rotation matrix may be computed and $`h_{xx}`$ is the (11) component of
$$𝐡_{lab}=𝐑_{wavelab}𝐡_{wave}𝐑_{wavelab}^T.$$
(13)
Considering only the component of the gravity wave along the bar axis,
$`h_{bar}(t)=(R_{11}^2R_{12}^2)h_+(t)+2R_{11}R_{12}h_\times (t)`$
where the relevant components of the rotation matrix are
$`\begin{array}{ccc}R_{11}\hfill & =\mathrm{cos}\psi _0(\mathrm{sin}\eta \mathrm{cos}\gamma +\mathrm{sin}l\mathrm{cos}\eta \mathrm{sin}\gamma )\hfill & \\ & +\mathrm{sin}\psi _0(\mathrm{sin}\eta \mathrm{sin}\delta \mathrm{sin}\gamma +\mathrm{sin}l\mathrm{cos}\eta \mathrm{sin}\delta \mathrm{cos}\gamma +\mathrm{cos}\eta \mathrm{cos}l\mathrm{cos}\delta )\hfill & \\ R_{12}\hfill & =\mathrm{cos}\psi _0(\mathrm{sin}\eta \mathrm{sin}\delta \mathrm{sin}\gamma +\mathrm{sin}l\mathrm{cos}\eta \mathrm{sin}\delta \mathrm{cos}\gamma +\mathrm{cos}\eta \mathrm{cos}l\mathrm{cos}\delta )\hfill & \\ & \mathrm{sin}\psi _0(\mathrm{sin}\eta \mathrm{cos}\gamma +\mathrm{sin}l\mathrm{cos}\eta \mathrm{sin}\gamma ).\hfill & \end{array}`$
Defining
$`f_+(t,\eta ,\gamma ,\delta ,l)`$ $`=`$ $`(\mathrm{sin}\eta \mathrm{cos}\gamma (t)+\mathrm{sin}l\mathrm{cos}\eta \mathrm{sin}\gamma (t))^2`$ (15)
$`(\mathrm{sin}\eta \mathrm{sin}\delta \mathrm{sin}\gamma (t)+\mathrm{sin}l\mathrm{cos}\eta \mathrm{sin}\delta \mathrm{cos}\gamma (t)+\mathrm{cos}\eta \mathrm{cos}l\mathrm{cos}\delta )^2`$
and
$`f_\times (t,\eta ,\gamma ,\delta ,l)`$ $`=`$ $`2(\mathrm{sin}\eta \mathrm{cos}\gamma (t)+\mathrm{sin}l\mathrm{cos}\eta \mathrm{sin}\gamma (t))`$ (17)
$`\times (\mathrm{sin}\eta \mathrm{sin}\delta \mathrm{sin}\gamma (t)+\mathrm{sin}l\mathrm{cos}\eta \mathrm{sin}\delta \mathrm{cos}\gamma (t)+\mathrm{cos}\eta \mathrm{cos}l\mathrm{cos}\delta )`$
we have
$`R_{11}^2R_{12}^2`$ $`=`$ $`\mathrm{cos}2\psi _0f_+(t)+\mathrm{sin}2\psi _0f_\times (t)`$ (18)
$`2R_{11}R_{12}`$ $`=`$ $`\mathrm{cos}2\psi _0f_\times (t)\mathrm{sin}2\psi _0f_+(t)`$ (19)
and
$$h_{bar}(t)=[\mathrm{cos}2\psi _0f_+(t)+\mathrm{sin}2\psi _0f_\times (t)]h_+(t)+[\mathrm{cos}2\psi _0f_\times (t)\mathrm{sin}2\psi _0f_+(t)]h_\times (t).$$
(20)
Using Eq. 8 and Eq. 20, the waveform is given by
$`h_{bar}(t)`$ $`=`$ $`h_c(1+\mathrm{cos}^2i)[\mathrm{cos}2\psi _0f_+(t)+\mathrm{sin}2\psi _0f_\times (t)]\mathrm{cos}(\omega _0t+\mathrm{\Phi }(t)+\mathrm{\Phi }_0)`$ (21)
$`+`$ $`h_c(2\mathrm{cos}i)[\mathrm{cos}2\psi _0f_\times (t)\mathrm{sin}2\psi _0f_+(t)]\mathrm{sin}(\omega _0t+\mathrm{\Phi }(t)+\mathrm{\Phi }_0).`$ (22)
The reception patterns $`f_+(t)`$ and $`f_\times (t)`$ are shown in Fig. 2 for a signal from 47 Tuc. Writing the time-dependent sine and cosine terms in Eq. 21 as exponentials and defining
$`F_+(t)=\mathrm{exp}[j\mathrm{\Phi }(t)]f_+(t)`$ (23)
$`F_\times (t)=\mathrm{exp}[j\mathrm{\Phi }(t)]f_\times (t)`$ (24)
the expression for the gravity wave strain along the bar axis is
$`h_{bar}(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}h_c(1+\mathrm{cos}^2i)\mathrm{exp}(j\mathrm{\Phi }_0)[\mathrm{cos}2\psi _0F_+(t)+\mathrm{sin}2\psi _0F_\times (t)]\mathrm{exp}(j\omega _0t)`$ (25)
$`+`$ $`{\displaystyle \frac{1}{2}}h_c(1+\mathrm{cos}^2i)\mathrm{exp}(j\mathrm{\Phi }_0)[\mathrm{cos}2\psi _0F_+^{}(t)+\mathrm{sin}2\psi _0F_\times ^{}(t)]\mathrm{exp}(j\omega _0t)`$ (26)
$``$ $`{\displaystyle \frac{j}{2}}h_c(2\mathrm{cos}i)\mathrm{exp}(j\mathrm{\Phi }_0)[\mathrm{cos}2\psi _0F_\times (t)\mathrm{sin}2\psi _0F_+(t)]\mathrm{exp}(j\omega _0t)`$ (27)
$`+`$ $`{\displaystyle \frac{j}{2}}h_c(2\mathrm{cos}i)\mathrm{exp}(j\mathrm{\Phi }_0)[\mathrm{cos}2\psi _0F_\times ^{}(t)\mathrm{sin}2\psi _0F_+^{}(t)]\mathrm{exp}(j\omega _0t).`$ (28)
We are now in position to calculate the anticipated signal as it would appear in the ALLEGRO data stream. This is most clearly presented in the frequency domain. Fourier transforming Eq. 11, the force produced on the bar is
$$(\omega )=\frac{1}{2}\mu L_e\omega ^2h_{bar}(\omega ).$$
(29)
This driving force produces motion of the transducer
$`H(\omega )=G(\omega )(\omega )={\displaystyle \frac{1}{2}}\mu L_e\omega ^2G(\omega )h_{bar}(\omega )`$
where $`G(\omega )`$ is the transfer function which relates transducer motion (or equivalently voltage) to the applied force . We note that the overall calibration is contained in the transfer function. Defining a new transfer function (strain to transducer motion)
$`G_f(\omega )={\displaystyle \frac{1}{2}}\mu L_e\omega ^2G(\omega )`$
we write the signal as it appears in the data as simply
$$H(\omega )=G_f(\omega )h_{bar}(\omega )$$
(30)
with $`h_{bar}(\omega )`$ given by the Fourier transform of Eq. 25.
At this point in the signal chain, in the interest of limiting the bandwidth of sample data required, the signal is mixed with a reference (whose frequency is chosen to be between the two detection mode frequencies) and low-pass filtered using a commercial lock-in amplifier . The lock-in provides both in-phase $`(x)`$ and quadrature $`(y)`$ outputs, so both amplitude and phase information is available on a bandwidth of 125 Hz containing the resonant modes of the detector. In software, the in-phase and quadrature components were combined to form a complex data stream
$`z(t)=x(t)+jy(t)`$
which was then demodulted to a 1 Hz bandwidth around each of the resonant modes . We describe the complete demodulation from 1 kHz to 1 Hz, including both the hardware and software lockins, as mixing the signal with a single reference $`\mathrm{exp}[j(\omega _rt+\varphi _r)]`$, where $`\omega _r`$ is the reference frequency and $`\varphi _r`$ is an unknown reference phase (this phase can in fact be measured, but for the data set in question was not). Returning to the time domain for clarity the mixed signal is
$$H(t)=[G_f(t)h_{bar}(t)]\mathrm{exp}[j(\omega _rt+\varphi _r)]$$
(31)
where ‘$``$’ indicates a convolution. Using the expression from Eq. 25 we get
$`H(t)`$ $`=`$ $`G_f(t)({\displaystyle \frac{1}{2}}h_c(1+\mathrm{cos}^2i)[\mathrm{cos}2\psi _0F_+^{}(t)+\mathrm{sin}2\psi _0F_\times ^{}(t)]\mathrm{exp}[j(\mathrm{\Phi }_0\varphi _r+(\omega _0\omega _r)t)]`$ (32)
$`+`$ $`{\displaystyle \frac{j}{2}}h_c(2\mathrm{cos}i)[\mathrm{cos}2\psi _0F_\times ^{}(t)\mathrm{sin}2\psi _0F_+^{}(t)]\mathrm{exp}[j(\mathrm{\Phi }_0\varphi _r+(\omega _0\omega _r)t)]).`$ (33)
Since $`\omega _r`$ is on the order of 1 kHz, the effect of the low-pass filtering is to remove terms from $`H(t)`$ which contain the sum frequencies $`(\omega _r+\omega _0)`$. Now returning to the frequency domain, the signal, after demodulation, can be written
$`H(\omega )`$ $`=`$ $`{\displaystyle \frac{1}{2}}h_c(1+\mathrm{cos}^2i)\mathrm{exp}[j(\mathrm{\Phi }_0\varphi _r)]G_f(\omega )[\mathrm{cos}2\psi _0F_+(\omega ^{^{}})+\mathrm{sin}2\psi _0F_\times (\omega ^{^{}})]`$ (34)
$``$ $`{\displaystyle \frac{j}{2}}h_c(2\mathrm{cos}i)\mathrm{exp}[j(\mathrm{\Phi }_0\varphi _r)]G_f(\omega )[\mathrm{cos}2\psi _0F_\times (\omega ^{^{}})\mathrm{sin}2\psi _0F_+(\omega ^{^{}})]`$ (35)
where $`\omega ^{^{}}=\omega \pm (\omega _0\omega _r)`$ is the (positive or negative) downconverted signal frequency.
## VI detection considerations
Equation 34 gives the form of the CW signal as it would appear in the Allegro data. We now ask: “Is this signal in the data?” Since the strength of the signal is small compared to the detector noise (otherwise we would see it on a spectrum analyzer!), some work needs to be done to answer this question. We will use the standard maximum likelihood method to guide the analysis, but in the end the question of whether a signal was detected or not will be answered experimentally.
The likelihood function can be written
$$\mathrm{\Lambda }=\frac{P(z|H)}{P(z|0)}$$
(36)
with $`P(H|z)`$ the probability that, given the observed data $`z(t)`$, the signal $`H(t)`$ is present, $`P(z|H)`$ the probability that, given the signal $`H(t)`$ is present, we observed the data $`z(t)`$, and $`P(z|0)`$ the probability that, given no signal present, we observe the data $`z(t)`$.
From , if
$`z=z_1,z_2,\mathrm{},z_N`$
is the time-ordered sequence of detector output (in our case imaginary) and
$`H=H_1,H_2,\mathrm{},H_N`$
is the signal waveform at the same sample times, then
$$P(z|0)=\frac{1}{[(4\pi )^NdetR_{gh}]}\mathrm{exp}(\frac{1}{2}\underset{g,h=0}{\overset{N1}{}}R_{gh}^1z_gz_h^{})$$
(37)
and
$$P(z|H)=\frac{1}{[(4\pi )^NdetR_{gh}]}\mathrm{exp}\left(\frac{1}{2}\underset{g,h=0}{\overset{N1}{}}R_{gh}^1(z_gH_g)(z_hH_h)^{}\right)$$
(38)
where the $`R_{gh}^1`$ is the inverse of the autocorrelation matrix.
Substituting from above, the likelihood ratio is
$$\mathrm{\Lambda }=\mathrm{exp}(\frac{1}{2}\underset{g,h=0}{\overset{N1}{}}R_{gh}^1(z_gH_h^{}H_gz_h^{}+H_gH_h))$$
(39)
which, given the properties of the complex autocorrelation matrix , can be written
$$\mathrm{\Lambda }=\mathrm{exp}\left(\mathrm{}\left\{\underset{g,h=0}{\overset{N1}{}}R_{gh}^1z_gH_h^{}\right\}\frac{1}{2}\underset{g,h=0}{\overset{N1}{}}R_{gh}^1H_gH_h^{}\right)$$
(40)
where $`\mathrm{}`$ means “take the real part”.
We define the following notation: $`x_k`$ for the discrete Fourier transform of $`x(t)`$, evaluated at $`\omega _k=2\pi k/N\mathrm{\Delta }t`$, and $`Sn_k^{(1)}`$ is the one-sided PSD. Using the relation
$$\underset{g,h=0}{\overset{N1}{}}R_{gh}^1x_gx_h^{}=2\mathrm{\Delta }f\underset{k=0}{\overset{N1}{}}\frac{x_kx_k^{}}{Sn_k^{(1)}}$$
(41)
and taking the natural logarithm of the the likelihood function (to remove the exponential), we have
$$\mathrm{ln}\mathrm{\Lambda }=2\mathrm{\Delta }f\mathrm{}\left\{\underset{k=0}{\overset{N1}{}}\frac{z_kH_k^{}}{Sn_k^{(1)}}\right\}\mathrm{\Delta }f\underset{k=0}{\overset{N1}{}}\frac{|H_k|^2}{Sn_k^{(1)}}.$$
(42)
Substituting for $`H`$ gives,
$`\mathrm{ln}\mathrm{\Lambda }`$ $`=`$ $`2\mathrm{\Delta }f\mathrm{}(\mathrm{exp}[j(\mathrm{\Phi }_0+\varphi _r)]{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle \frac{z_kG_{f_k}^{}}{Sn_k^{(1)}}}`$ (43)
$`\times `$ $`[{\displaystyle \frac{1}{2}}h_c(1+\mathrm{cos}^2i)\mathrm{cos}2\psi _0F_{+_k}^{}+{\displaystyle \frac{1}{2}}h_c(1+\mathrm{cos}^2i)\mathrm{sin}2\psi _0F_{\times _k}^{}`$ (44)
$``$ $`{\displaystyle \frac{j}{2}}h_c(2\mathrm{cos}i)\mathrm{cos}2\psi _0F_{\times _k}^{}+{\displaystyle \frac{j}{2}}h_c(2\mathrm{cos}i)\mathrm{sin}2\psi _0F_{+_k}^{}])`$ (45)
$``$ $`\mathrm{\Delta }f{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle \frac{1}{Sn_k^{(1)}}}`$ (46)
$`\times `$ $`[{\displaystyle \frac{1}{4}}h_c^2(1+\mathrm{cos}^2i)^2\mathrm{cos}^22\psi _0|G_{f_k}F_{+_k}|^2+{\displaystyle \frac{1}{4}}h_c^2(1+\mathrm{cos}^2i)^2\mathrm{sin}^22\psi _0|G_{f_k}F_{\times _k}|^2`$ (47)
$`+`$ $`{\displaystyle \frac{1}{4}}h_c^2(2\mathrm{cos}i)^2\mathrm{cos}^22\psi _0|G_{f_k}F_{\times _k}|^2+{\displaystyle \frac{1}{4}}h_c^2(2\mathrm{cos}i)^2\mathrm{sin}^22\psi _0|G_{f_k}F_{+_k}|^2]`$ (48)
where terms of the form
$`{\displaystyle \underset{k=0}{\overset{N1}{}}}F_{+_k}F_{\times _k}^{}`$
sum to zero.
To simplify the notation of Eq. 43, we make the following definitions:
$`q_+`$ $`=`$ $`\mathrm{\Delta }f{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle \frac{z_kG_{f_k}^{}F_{+_k}^{}}{Sn_k^{(1)}}}`$ (49)
$`q_\times `$ $`=`$ $`\mathrm{\Delta }f{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle \frac{z_kG_{f_k}^{}F_{\times _k}^{}}{Sn_k^{(1)}}}`$ (50)
$`\rho _+`$ $`=`$ $`\mathrm{\Delta }f{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle \frac{|G_{f_k}F_{+_k}|^2}{Sn_k^{(1)}}}`$ (51)
$`\rho _\times `$ $`=`$ $`\mathrm{\Delta }f{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle \frac{|G_{f_k}F_{\times _k}|^2}{Sn_k^{(1)}}}`$ (52)
$`a_1`$ $`=`$ $`h_c(1+\mathrm{cos}^2i)\mathrm{cos}2\psi _0`$ (53)
$`a_2`$ $`=`$ $`h_c(1+\mathrm{cos}^2i)\mathrm{sin}2\psi _0`$ (54)
$`a_3`$ $`=`$ $`h_c(2\mathrm{cos}i)\mathrm{cos}2\psi _0`$ (55)
$`a_4`$ $`=`$ $`h_c(2\mathrm{cos}i)\mathrm{sin}2\psi _0.`$ (56)
The four quantities $`\{q_+,q_\times ,\rho _+,\rho _\times \}`$ are familiar from signal processing theory. The first two, $`\{q_+,q_\times \}`$, are the outputs from applying independent optimal filters in the frequency domain for the two polarizations of the gravity wave. The second pair, $`\{\rho _+,\rho _\times \}`$ are the signal to noise ratios (energy) for the two polarizations. These four terms are completely specified, up to the unknown signal frequency. The four $`a_i`$’s are the “amplitudes” containing the desired astrophysical information.
Substituting these into the likelihood function,
$`\mathrm{ln}\mathrm{\Lambda }`$ $`=`$ $`a_1\mathrm{}\left\{\mathrm{exp}[j(\mathrm{\Phi }_0\varphi _r)]q_+\right\}+a_2\mathrm{}\left\{\mathrm{exp}[j(\mathrm{\Phi }_0\varphi _r)]q_\times \right\}`$ (57)
$``$ $`a_3\mathrm{}\left\{j\mathrm{exp}[j(\mathrm{\Phi }_0\varphi _r)]q_\times \right\}+a_4\mathrm{}\left\{j\mathrm{exp}[j(\mathrm{\Phi }_0\varphi _r)]q_+\right\}`$ (58)
$``$ $`{\displaystyle \frac{1}{4}}(a_1^2\rho _++a_2^2\rho _\times +a_3^2\rho _\times +a_4^2\rho _+).`$ (59)
We maximize the likelihood function over the $`a_i`$’s which yields the following four expressions:
$`a_1`$ $`=`$ $`2{\displaystyle \frac{|q_+|}{\rho _+}}\mathrm{cos}(\mathrm{\Phi }_0\varphi _r+\varphi _+)`$ (60)
$`a_2`$ $`=`$ $`2{\displaystyle \frac{|q_\times |}{\rho _\times }}\mathrm{cos}(\mathrm{\Phi }_0\varphi _r+\varphi _\times )`$ (61)
$`a_3`$ $`=`$ $`2{\displaystyle \frac{|q_\times |}{\rho _\times }}\mathrm{sin}(\mathrm{\Phi }_0\varphi _r+\varphi _\times )`$ (62)
$`a_4`$ $`=`$ $`2{\displaystyle \frac{|q_+|}{\rho _+}}\mathrm{sin}(\mathrm{\Phi }_0+\varphi _r\varphi _+)`$ (63)
where the filtered ouputs have been expressed as
$`q_{(+,\times )}=|q_{(+,\times )}|\mathrm{exp}(j\varphi _{(+,\times )})`$
with $`\varphi _{(+,\times )}=\mathrm{arg}[q_{(+,\times )}]`$. Finally,
$$\underset{i=1}{\overset{4}{}}a_i^2=4(\frac{|q_+|^2}{\rho _+^2}+\frac{|q_\times |^2}{\rho _\times ^2})=h_c^2[(1+\mathrm{cos}^2i)^2+4\mathrm{cos}^2i]=h_+^2+h_\times ^2$$
(64)
which, to factors, is the energy of the gravitational wave. Note that this depends on the unknown signal frequency $`f_0`$. As described below, we will assume values for the gravitational wave frequency, and use that to calculate Eq. 64. The end result of the analysis is a “spectrum” of energy at a given signal frequency. We shall report this as a “strain amplitude”, given by the square root of Eq. 64
$`h_s(f_0)=\sqrt{h_+^2(f_0)+h_\times ^2(f_0)}.`$
## VII experimental considerations
Calculation of Eq. 64 involved a number of steps. First, the archived data was read off tape and narrowbanded to a 1 Hz bandwidth around each of the resonant modes where ALLEGRO is most sensitive. This was done to reduce the computational requirements. Next, the mode amplitudes were “cleaned” of large, transient events (section VII A). Finally, the cleaned data was discrete Fourier transformed and optimal filters were applied (section VII B).
### A data selection
Even with ALLEGRO’s high duty cycle, there were periods of missing or unusable data. Data losses came in basically three flavors. (1) Transient electronic effects which lasted on the order of a second. (2) Longer periods when the detector was undergoing some form of maintenance. (3) The clock losing phase lock to WWVB. The transient disturbances were the most frequent, occurring at a rate of one or two per day. They usually involved a sudden change in the flux threading the SQUID loop (hence the name “flux jumps”). The majority of the flux jumps occurred when the dc level of the SQUID reached a pre-determined maximum (5 volts). The electronics controlling the SQUID then reset the dc voltage to zero, causing a short and violent jump in the in-phase and quadrature channels of the data, as shown in Fig. 3.
Frequently electronic interference reaching the SQUID caused flux jumps. In the past when data tapes were erased (the degaussing takes place in a separate room from the main experiment), the end of the degaussing cycle produced a noise spike which traveled through the wiring in the wall, through the computers, and from there to the SQUID. Once recognized, an isolating transformer was placed between the degausser and the wall socket, fixing the problem. Another common type of transient signal is a “spike”. These look similar to flux jumps in the data and there is some suspicion that they are in fact flux jumps, but essentially they are of unknown origin. All of these noisy periods were short enough so that the affected data could be removed and the resulting gap interpolated across. This was done using the MATLAB interp1 routine. Interpolation was performed on the resampled data. It was usual to interpolate across 1-5 seconds of data, using the 10-20 seconds of data before and after the gap for the interpolation template.
For longer sections of unusable data or for periods of missing data, the analysis was stopped and restarted after the disturbance. By “restarted” we mean that the accumulation of data was stopped, the accumulated data purged, and the analysis started up again on the data immediately following the disturbance. The most common cause of data loss was transferring liquid helium into the dewar which removed a couple of hours of data every week. Another cause of long stretches of unusable data were large excitations of the resonant modes due to earthquakes around the globe. Earthquakes were identified by a unique signature in the low frequency housekeeping channel. It was usual for an earthquake to produce multiple large excitations over a few tens of minutes, often resulting in saturation of the A/D’s. Computer down time, calibration of the detector and other maintenance all caused gaps in the data, although infrequently.
The final type of data loss was associated with the WWVB clock. Frequently when the weather between Baton Rouge and NIST at Boulder, Colorado was bad, the clock we used to control the sampling of the data lost phase lock to the WWVB radio signal. When this happened, the clock’s internal oscillator “freewheeled” with the result that the time between samples was no longer consistently 8 ms. Deviations in the sampling rate from 8 ms were called “timing jumps”. A jump in the time between when samples were taken produced a corresponding jump in the phase of modes and the calibrator signal as shown in Fig. 4. The most frequent jumps were on the order of 1-2 ms, producing a phase jump in a sinusoidal signal at the mode frequencies of approximately 1/100 of a cycle. This was considered an acceptable level of uncertainty in the ability to track the phase of a potential gravity wave signal. These small jumps were noted but ignored. Larger jumps produced correspondingly larger jumps in phase and were considered unacceptable. When they occurred, the analysis was stopped at the timing jump and restarted again after the glitch. The frequency and size of the timing jumps were highly variable. During the winter of 1994 the smaller jumps occurred almost once per day while the larger jumps rarely happened. By the spring of that year, the trend was reversed and much data was lost due to the inconsistency of the clock.
### B the filters
The ability of our analysis filters to match the phase of the gravity wave determined the amount of continuous data which could be analyzed at one time. We refer to this as a “record”. Ideally, the length of a record would be set by the available computational facilities. This was not the case. Instead, the record length was limited by our electronics. The function generator which supplied the reference frequency to the lock-in detector had some drift associated with it. This drift limits the ability of the filter to match the phase of the gravity wave signal in the data. Careful measurements of the drift led to the conservative conclusion that roughly 28 hours ($`10^5`$ sec) was the longest period of time for which we could expect the filter to remain in phase with the signal . The less conservative conclusion was roughly three times longer, but we prefer to be cautious with the analysis. The reference signal to the lock-in has since been phase locked to GPS, greatly improving the phase stability. Given the above considerations, there were 34 records available from days 1-94 of 1994, for a total of 944.44 hours of data.
Performing a discrete Fourier transform on $`10^5`$ sec worth of data results in a frequency resolution of the search of $`10^5`$ Hz. The optimal filters also depend on the signal frequency. With no known source available, we assume a there is a potential signal every 10 $`\mu `$Hz in the ranges 896.30-897.30 Hz and 919.76-920.76 Hz. Each assumed signal frequency was used in turn to create the $`F_+`$ and $`F_\times `$ components of the signal template.
The other two components of the filters, the bar transfer function and the power spectral density, depend on the resonant frequency and damping time of the coupled bar-transducer system. As both quantities experienced slow drifts due to, for example, temperature changes in the dewar, it was necessary to measure them for each record (see Fig. 5 and Fig. 6 respectively). For each mode, a low variance power spectral density was formed from the cleaned data. A Lorenzian curve was then fit to each PSD using the MATLAB curvefit routine, which finds the best fit to a function in the least-squares sense. The Lorentzian was characterized by four parameters: the white-noise level ($`S_0`$), the peak height ($`S_1`$), the frequency of the resonance ($`\omega _\pm `$), and the width (or the damping time - $`\tau _\pm `$). The resulting fitted parameters for each record were then stored on disk to be retrieved as necessary. The power spectral density component of the optimal filter was calculated for a record by
$$Sn(\omega _k)=S_0+\frac{(S_1S_0)}{\tau _\pm ^2}[(\omega _k^2\omega _\pm ^2)^2+\omega _k^2/\tau _\pm ^2]^1.$$
(65)
Similarly, the bar transfer function was calculated from the functional
$$G(\omega _k)=k_\pm \frac{\omega _\pm }{\tau _\pm }[\omega _\pm ^2+j\omega _k/\tau _\pm \omega _k^2]^1$$
(66)
where the parameters are those given above. The constant $`k_\pm `$ sets the calibration for each mode .
### C computational reduction
Even though this search was directed towards only a small section of the sky, it was still cpu intensive on the available computing facilities. Applying $`2\times 10^5`$ different filters, where each new set of filter coefficients required a Fourier transform on a time sequence of $`10^5`$ elements, was prohibitively time consuming.
The signal frequency enters into calculation of the filter coefficients in two places; through the carrier wave and in the Doppler shift. However, as the phase modulation is slowly varying with respect to the carrier wave, it was not necessary to re-calculate the phase correction due to the Doppler shift for each assumed signal frequency. Instead, the Doppler shift was calculated at specific choices of the signal frequency, which were then used to approximate signal templates over a range of frequencies.
The maximum fractional frequency shift of a signal from the two sources over the course of a year is
$`{\displaystyle \frac{|\delta f(f_0)|}{f_0}}\{\begin{array}{cc}5\times 10^5\hfill & \text{ 47 Tuc}\hfill \\ 10^4\hfill & \text{ galactic center}\hfill \end{array}`$
at the signal frequencies of interest. To approximate all phase shifts in a range of signal frequencies $`f_0\pm \mathrm{\Delta }f_0`$ with a calculation at a single signal frequency, then the error made must be less than the frequency resolution of the search. Defining the error as
$`\delta f(f_0\pm \mathrm{\Delta }f_0)\delta f(f_0)=5\times 10^5\mathrm{\Delta }f_0`$
we have the maximum range of signal frequencies over which the approximation is valid is given by
$`|\mathrm{\Delta }f_0|\{\begin{array}{cc}0.2\hfill & \text{Hz, 47 Tuc}\hfill \\ 0.1\hfill & \text{Hz, galactic center}\hfill \end{array}`$
Using this information, we calculated independent $`F_+,F_\times `$ components of the signal template at the frequencies 896.45 Hz, 896.80 Hz, 897.15 Hz, 919.91 Hz, 920.26 Hz and 920.61 Hz for the 47 Tuc analysis. For the galactic center analysis, signal templates were calculated every 0.2 Hz in the range 896.40 Hz - 897.20 Hz for the minus mode, and 919.86 Hz - 920.66 Hz for the plus mode.
## VIII results
The real and imaginary components of the data stream, in the absence of a signal, are zero mean Gaussian distributed. Neither the Fourier transform or the filtering changes the underlying distribution and therefore the real and imaginary parts of the filtered outputs $`q_{(+,\times )}`$ are also individually zero mean Gaussian distributed. Given this, calculation of Eq. 64 at a particular choice of the signal frequency, taken from a particular data record, results in a sample drawn from a chi-squared distribution with four degrees of freedom . If the signal is absent, this is given by
$$p(E)=\frac{1}{(2\sigma ^2)^2}E\mathrm{exp}(E/2\sigma ^2).$$
(67)
where $`E=h_s^2(f_0)`$.
The total number of available data records provide an ensemble of experiments each making $`2\times 10^5`$ measures of $`E`$ at the assumed values for $`f_0`$. Since the signal is expected to be long-lived, much longer than the operational time of the detector, we improve our estimate of the energy by taking the ensemble average of the individual measures:
$$\overline{E}=\frac{1}{N}\underset{i=1}{\overset{N}{}}E_i$$
(68)
where $`N=34`$ is the number of available data records. We calculate the strain amplitude $`h_s(f_0)`$ by taking the square root of $`\overline{E}`$. As each ensemble is generated by the sum of four terms (Eq. 64) and there are 34 data records averaged together, from the central limit theorm we can accurately describe the distribution of $`h_s(f_0)`$ as gaussian.
The resulting “spectra” of $`h_s(f_0)`$ vs. $`f_0`$ are shown in Fig. 7 for 47 Tuc and Fig. 8 for the galactic center analysis. It is important to remember that the abscissa is not a Fourier frequency but rather the assumed frequency of the gravitational wave. That the plus mode shows a lower strain value than the minus mode reflects the fact that the plus mode has a higher quality factor and is therefore more sensitive. That the results of the galactic center analysis seem somewhat “shifted” with respect to the 47 Tuc results is due to the effects of the Doppler shifts. For a signal from 47 Tuc during the period over which data was analyzed, a signal at a particular frequency, for example 920 Hz, experienced periods of both red shift and blue shift. On average, therefore, the signal was experiencing detector noise at roughly 920 Hz. For the galactic center analysis, during this period an incoming signal was experiencing maximum blue-shift, and was never red-shifted. Therefore, a signal at 920 Hz would be experiencing detector noise at a slightly higher frequency, resulting in a graph that seems “a bit skewed to the right” with respect to the 47 Tuc analysis.
The question still remains if there is a real signal at a particular frequency. It is unlikely that the detector is dominated by CW radiation. We expect the number of real signals which would be extracted by our analysis to be small with respect to the number of possible signals (i.e. the number of frequencies at which we assume there to be a signal). We can therefore answer the question experimentally. We normalize $`h_s(f_0)`$ at each $`f_0`$ to its mean value. This new variable has a distribution independent of the detector noise at each signal frequency.
The histograms of the resulting normalized spectra are shown in Fig. 9 for 47 Tuc and Fig. 10 for the galactic center analysis. Given our reasonable assumption that the detector is not wave-dominated, the majority of the normalized measures form an experimental estimation of the parent distribution for the detector noise from which each individual measure is drawn. If the strain amplitude at a particular signal frequency lies outside this distribution, it is identified as a candidate for a real signal.
## IX discussion
Neither Fig. 9 or Fig. 10 shows any evidence for “outliers” which would suggest the possibility of having observed gravitational radiation from pulsar spin-down. The existence of burst-like non-gaussian noise in resonant bar data is well known, and this noise must be dealt with in any search for burst-like signals of gravitational waves. The encouraging result of this analysis is that for CW signals there is not a serious problem with non-gaussian noise in the resonant bar data. Thus, at least at the level of the current work, all the noise sources are well understood and the overall sensitivity may be improved in the future simply by inceasing the length of the data record. We note that care was taken to use only data when the detector was operating well and procedures were implemeted to smooth out other non-stationary effects. Astrophysically, at this level of sensitivity, the actual spin-down rate of a pulsar due to the energy lost as gravitationl radiation would be large enough that the signal and its filter template would go out of phase by half a cycle after of order 5000 s, much less than the length of one record. This effect can be taken into account by including one or more spin-down parameters in the signal template, at the cost of significantly increasing the computational time of the search. Since there was no known source for ALLEGRO and somewhat limited computing resources available, we have essentially made a demonstration of the capabilities of resonant detectors for this type of search.
However, it is important to note that even with these restrictions, the analysis has reached a level of sensitivity which is astrophysically interesting. If a source were to become known with the correct frequency, and, as our sensitivity was limited by hardware which has been substantially improved, it seems likely that such a signal could be detected and new astrophysics learned.
## X conclusion
We have searched for sources of continuous gravitational radiation from the globular cluster 47 Tuc and the galactic center at signal frequencies near 1 kHz using data taken by the ALLEGRO gravitational radiation detector in the first three months of 1994. No candidate signals were found, and a constraint of $`8\times 10^{24}`$ was put on gravitational radiation emitted from pulsar spin-down at both locations.
###### Acknowledgements.
The authors thank Sam Finn and Joel Tohline for many helpful discussions. Much of the analysis and the final writing of this paper was done at the INFN Laboratori Nazionali di Frascati. E.M. thanks G. Pizzella, E. Coccia, and the ROG collaboration for their support. The LSU group was supported by the National Science Foundation under Grant No. PHY-9311731.
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# Quantum mechanics with non-unitary symmetries
## 1 Overview of Dirac equation and second quantization
Dirac originally proposed the relativistic equation of the form
$$i\frac{\psi }{t}=\left(\frac{1}{i}\stackrel{}{\alpha }\stackrel{}{}+\beta m\right)\psi =H\psi $$
(10)
where $`\alpha _i`$ and $`\beta `$ are anticommuting matrices satisfying:
$`\alpha _i^2=\beta ^2`$ $`=`$ $`1`$ (11)
$`\{\alpha _i,\alpha _j\}`$ $`=`$ $`0`$ (12)
$`\{\alpha _i,\beta \}`$ $`=`$ $`0.`$ (13)
One would then multiply the Dirac equation (10) by $`\beta `$ and rewrite it in a covariant notation
$$\left(i\gamma ^\mu _\mu m\right)\psi =0$$
(14)
where matrices $`\gamma ^\mu `$ are defined as
$$\gamma ^0=\beta ,\gamma ^i=\beta \alpha ^i,\{\gamma ^\mu ,\gamma ^\nu \}=2g^{\mu \nu }.$$
(15)
By taking the Hermitean conjugate of the equation (14) we get the another one for the conjugate of $`\psi `$:
$$_\mu \overline{\psi }i\gamma ^\mu +\overline{\psi }m=0$$
(16)
with $`\overline{\psi }\psi ^{}\gamma _0`$.
The next step would usually be to prove that this equation is relativistically invariant and then find the Lagrangian which reproduces this equation and use it to find the energy-momentum tensor, conserved current and angular momentum tensor. This Lagrangian is found to be<sup>1</sup><sup>1</sup>1 derivative operator is defined with the appropriate sign and factor to give (after partial integration) ψ
ϕ=ψ
ϕ=ψ
ϕ𝜓
italic-ϕ𝜓
italic-ϕ𝜓
italic-ϕ\int\psi\kern 1.66672pt\vbox{\hbox{\kern-1.66672pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}\hbox to-1.00006pt{}\phi=\int\psi\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}\hbox to0.0pt{}\phi=\int\psi\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}\hbox to0.0pt{}\phi, or explicitly f(x)
f(x)/x𝑓𝑥
𝑓𝑥𝑥f(x)\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}\hbox to0.0pt{}\equiv-{\partial f(x)}/{\partial x} and f(x)
g(x){f(x)g(x)/xf(x)/xg(x)}/2𝑓𝑥
𝑔𝑥𝑓𝑥𝑔𝑥𝑥𝑓𝑥𝑥𝑔𝑥2f(x)\kern 1.66672pt\vbox{\hbox{\kern-1.66672pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}\hbox to-1.00006pt{}g(x)\equiv\{f(x)\partial g(x)/\partial x-\partial f(x)/\partial x\>g(x)\}/2. Conjugation properties of these operators are (
)=
superscript
(\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}\hbox to0.0pt{})^{\dagger}=-\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}\hbox to0.0pt{}, (
)=
superscript
(\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}\hbox to0.0pt{})^{\dagger}=-\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}\hbox to0.0pt{}, (
)=
superscript
(\kern 1.66672pt\vbox{\hbox{\kern-1.00006pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}\hbox to-1.00006pt{})^{\dagger}=-\kern 1.66672pt\vbox{\hbox{\kern-1.00006pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}\hbox to-1.00006pt{}.
D=ψ¯(x)(i
/
m)ψ(x)subscript𝐷¯𝜓𝑥𝑖
/
𝑚𝜓𝑥\mathcal{L}_{D}=\bar{\psi}(x)\left(i\kern 1.66672pt\vbox{\hbox{\kern-1.00006pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial\hskip-5.50003pt{/}$}}\hbox to-1.00006pt{}-m\right)\psi(x) (17)
with $`p/\gamma p`$. Going to the momentum space one finds the solutions of equations (14) and (16) to be
$`\psi (x)`$ $`=`$ $`{\displaystyle \frac{d^3q}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_q}\left(e^{iqx}u_{qr}b_{qr}+e^{iqx}v_{qr}d_{qr}^{}\right)}`$
$`\overline{\psi }(x)`$ $`=`$ $`{\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_p}\left(e^{ipx}\overline{u}_{ps}b_{ps}^{}+e^{ipx}\overline{v}_{ps}d_{ps}\right)}`$ (18)
Although at this point $`b_{ps}`$ and $`d_{ps}`$ are complex numbers, they will be treated as noncommuting quantities so that all results derived would still be valid after the second quantization. If we require the solution (1) to be invariant to the gauge transformation
$$\psi (x)\psi ^{}(x)=e^{i\alpha }\psi \overline{\psi }(x)\overline{\psi }^{}(x)=\overline{\psi }e^{i\alpha }$$
(19)
or for infinitesimal $`\alpha `$
$$\psi (x)(1i\alpha )\psi \overline{\psi }(x)\overline{\psi }(1+i\alpha )$$
(20)
we get the conserved current
$$j^\mu (x)=\overline{\psi }(x)\gamma ^\mu \psi (x).$$
(21)
Requiring the translational invariance
$$\psi (x)\psi ^{}(x^{})=\psi (x+a)\overline{\psi }(x)\overline{\psi }^{}(x^{})=\overline{\psi }(x+a)$$
(22)
or for infinitesimal $`a`$
$$\psi (x)(1+a_\mu ^\mu )\psi (x)\overline{\psi }(x)(1+a_\mu ^\mu )\overline{\psi }(x)$$
(23)
one gets the energy momentum tensor
$`\theta ^{\mu \nu }(x)`$ $`=`$ ψ¯(x)γμi
ψν(x)¯𝜓𝑥superscript𝛾𝜇𝑖
superscript𝜓𝜈𝑥\displaystyle\bar{\psi}(x)\gamma^{\mu}i\kern 1.66672pt\vbox{\hbox{\kern-1.66672pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}{}^{\nu}\psi(x) (24)
Following the same procedure for rotations and Lorentz boosts
$$\psi (x)\psi ^{}(x^{})=S(\omega )\psi (x)\overline{\psi }(x)\overline{\psi }^{}(x^{})=\overline{\psi }(x)\gamma ^0S^{}(\omega )\gamma ^0=\overline{\psi }(x)S^1(\omega )$$
(25)
or for infinitesimal $`\omega `$
$`\psi (x)`$ $``$ $`\left\{1{\displaystyle \frac{i}{2}}\omega _{\mu \nu }\left({\displaystyle \frac{\sigma ^{\mu \nu }}{2}}+x^\mu i^\nu x^\nu i^\mu \right)\right\}\psi (x)`$
$`\overline{\psi }(x)`$ $``$ ψ¯(x){1+i2ωμν(σμν2+xμi
νxνi
)μ}\displaystyle\bar{\psi}(x)\left\{1+\frac{i}{2}\omega_{\mu\nu}\left(\frac{\sigma^{\mu\nu}}{2}+x^{\mu}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\nu}-x^{\nu}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\mu}\right)\right\} (26)
we get the generalized angular momentum density tensor
$`J^{\mu ,\alpha \beta }`$ $`=`$ ψ¯(x)[γμ(xαi
βxβi
)α+12{γμ,σαβ2}]ψ(x)\displaystyle\bar{\psi}(x)\left[\gamma^{\mu}\left(x^{\alpha}i\kern 1.66672pt\vbox{\hbox{\kern-1.66672pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}{}^{\beta}-x^{\beta}i\kern 1.66672pt\vbox{\hbox{\kern-1.66672pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}{}^{\alpha}\right)+\frac{1}{2}\left\{\gamma^{\mu},\frac{\sigma^{\alpha\beta}}{2}\right\}\right]\psi(x) (27)
$`=`$ $`x^\alpha \theta ^{\mu \beta }x^\beta \theta ^{\mu \alpha }+{\displaystyle \frac{1}{2}}\overline{\psi }(x)\{\gamma ^\mu ,{\displaystyle \frac{\sigma ^{\alpha \beta }}{2}}\}\psi (x)`$
While all of the above looks nice, there are a few problems with this formulation. It has been believed some of them are solved by second quantization, but some remain even after second quantization.
While the positive energy and negative energy momentum eigenstates of different spin
$`\psi _{ps}^+(x)`$ $`=`$ $`{\displaystyle \frac{e^{ipx}}{(2\pi )^{3/2}}}u_{ps}`$
$`\psi _{ps}^{}(x)`$ $`=`$ $`{\displaystyle \frac{e^{ipx}}{(2\pi )^{3/2}}}v_{ps}`$ (28)
*are* separately orthogonal<sup>2</sup><sup>2</sup>2 Note that the scalar product isn’t defined as $`\psi ^{}\psi d^3x`$ but as $`\psi ^{}\gamma _0\psi d^3x`$. Reasons for this will be explained in section 2.1 <sup>3</sup><sup>3</sup>3there is a minus sign for the norm of $`\psi ^{}`$ solution that seems to contradict positiveness of the norm. That will also be discussed in section 2.1
$`{\displaystyle \overline{\psi }_{ps}^+(x)\psi _{qr}^+(x)d^3x}`$ $`=`$ $`\delta ^3(pq)\delta _{r,s}`$
$`{\displaystyle \overline{\psi }_{ps}^{}(x)\psi _{qr}^{}(x)d^3x}`$ $`=`$ $`\delta ^3(pq)\delta _{r,s}`$ (29)
scalar products of positive and negative energy momentum eigenstates *aren’t* vanishing, making the solutions non-orthogonal
$`{\displaystyle \overline{\psi }_{ps}^+(x)\psi _{qr}^{}(x)d^3x}`$ $`=`$ $`\delta ^3(p+q)e^{2iE_pt}\overline{u}_{ps}v_{\stackrel{~}{p},r}`$
$`{\displaystyle \overline{\psi }_{ps}^{}(x)\psi _{qr}^+(x)d^3x}`$ $`=`$ $`\delta ^3(p+q)e^{2iE_pt}\overline{v}_{ps}u_{\stackrel{~}{p},r}`$ (30)
with $`\stackrel{~}{p}^\mu =(p^0,\stackrel{}{p})`$. The complete set of the solutions to the Dirac equation isn’t orthogonal so the decomposition
$$\psi (x)=\frac{d^3p}{(2\pi )^{3/2}2E_p}\left(e^{ipx}u_{ps}b_{ps}+e^{+ipx}v_{ps}d_{ps}^{}\right)$$
(31)
or
$$\psi (x)=\frac{d^3p}{2E_p}\left(\psi _{ps}^+(x)b_{ps}+\psi _{ps}^{}(x)d_{ps}^{}\right)$$
(32)
isn’t a valid decomposition in a complete orthonormal set of functions. One deals with this by saying that we work with fields, not wave functions, and therefore there’s no need for orthogonality anyway.
Another “problem” is a consequence of the definition of eigenstates (1). Applying the momentum operator $`\stackrel{}{P}i`$ to positive states produces the proper eigenvalue for the positive energy states
$$i\psi _{ps}^+=\stackrel{}{p}\psi _{ps}^+$$
(33)
but it gives us the wrong sign for the negative energy states
$$i\psi _{ps}^{}=\stackrel{}{p}\psi _{ps}^{}.$$
(34)
We can think of momentum eigenstates as states obtained by applying the Lorentz boost in the direction $`\stackrel{}{p}`$ and with the boost parameter $`\left|\stackrel{}{p}\right|`$ to the particle in it’s rest frame ($`\stackrel{}{p}=0`$). After we boost the particle in one direction, negative energy particles appear to move in the opposite direction. To solve this problem Dirac proposed the hole interpretation saying that in physical vacuum all the negative energy states were filled and when an electron from this sea gets excited to positive energy it leaves a hole that we observe as anti-particle carrying the opposite charge. Stückelberg and a little later Feynman proposed the interpretation that “negative energy” solution of energy $`E`$ and momentum $`\stackrel{}{p}`$ represents a particle moving backwards in time which we observe as the particle of the opposite charge, energy and momentum moving forward in time.
Further problems appear when one considers conserved currents of the theory. The current
$`j^\mu (x)`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_p}}{\displaystyle \frac{d^3q}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_q}}(e^{i(pq)x}b_{ps}^{}b_{qr}\overline{u}_{ps}\gamma ^\mu u_{qr}+e^{i(pq)x}d_{ps}d_{qr}^{}\overline{v}_{ps}\gamma ^\mu v_{qr}`$ (35)
$`+e^{i(p+q)x}b_{ps}^{}d_{qr}^{}\overline{u}_{ps}\gamma ^\mu v_{qr}+e^{i(p+q)x}d_{ps}b_{qr}\overline{v}_{ps}\gamma ^\mu u_{qr})`$
is formally conserved $`_\mu j^\mu (x)=0`$ which allow us to construct conserved charge by integrating over the whole volume $`𝐑^3`$
$`Q={\displaystyle j^0(x)d^3x}={\displaystyle \frac{d^3p}{2E_p}\left(b_{ps}^{}b_{ps}+d_{ps}d_{ps}^{}\right)}.`$ (36)
However, space part (with the help of Gordon identities from appendix C.3) of the total current
$`J^i`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d^3p}{2E_p}}{\displaystyle \frac{1}{E_p}}(p^i{\displaystyle \underset{s}{}}[b_{ps}^{}b_{ps}+d_{ps}d_{ps}^{}]`$ (37)
$`+{\displaystyle \underset{s,r}{}}[{\displaystyle \frac{ie^{2iE_pt}}{2m}}\overline{u}_{ps}\sigma ^{i0}v_{\stackrel{~}{p}r}b_{ps}^{}d_{\stackrel{~}{p}r}^{}{\displaystyle \frac{ie^{2iE_pt}}{2m}}\overline{v}_{ps}\sigma ^{i0}u_{\stackrel{~}{p}r}d_{ps}b_{\stackrel{~}{p}r}])`$
(as well as current density) besides the group velocity term has a real oscillating term. This term of the order of magnitude $`10^{21}s`$ traditionally called *zitterbewegung* has been without proper physical interpretation and is another reason physics community decided that single particle theories don’t work and should be discarded/reinterpreted, although it is present even after second quantization.
When integrating the zeroth component over the infinite volume, *zitterbewegung* terms are still present but don’t contribute since all the fields vanish at infinity. If one integrates over *finite* volume $`V`$,
$$\frac{d}{dt}\underset{V}{}j^0(x)d^3x=\underset{V}{}\stackrel{}{j}(x)d^3x=\underset{S}{}\stackrel{}{n}\stackrel{}{j}(x)𝑑S$$
(38)
where $`\stackrel{}{n}`$ is the unit vector normal to surface $`S`$, *zitterbewegung* terms on the right-hand side make it hard to interpret change of the charge contained in volume $`V`$ as the divergence of the flux of the current over the edge of the volume $`S`$.
After second quantization charge $`Q`$ doesn’t annihilate the vacuum, so one has to remove the divergent part “by hand” by introducing so called *normal ordering*. Even after normal ordering, zitterbewegung is still present; creation and annihilation operators from *zitterbewegung* part of the current come in combinations
$$J^i(x)bd+b^{}d^{}$$
(39)
which doesn’t annihilate the vacuum
$$J^i(x)|00$$
(40)
and mixes states with $`n`$ particles with states with $`n+e^+e^{}\text{ pair}`$ states, for example vacuum and electron-positron pair
$$0\left|J^i\right|e^+e^{}0.$$
(41)
After one introduces electromagnetic interactions through
$$_i=j^\mu (x)A_\mu (x),$$
(42)
those terms give infinite contributions to both higher order vacuum to vacuum transition matrix elements as well as infinite contributions to higher order (loop) diagrams.
Next set of conserved currents is energy-momentum tensor. First of all, energy momentum tensor isn’t symmetric. This is in itself enough of a problem and requires procedure for symmetrization which is (again) introduced “by hand” and isn’t a consequence of any deeper principle. That set aside for a moment, again, energy-momentum tensor
$`\theta ^{\mu \nu }(x)`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_p}}{\displaystyle \frac{d^3q}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_q}}(e^{i(pq)x}b_{ps}^{}b_{qr}\overline{u}_{ps}\gamma ^\mu u_{qr}{\displaystyle \frac{(p+q)^\nu }{2}}`$ (43)
$`e^{i(pq)x}d_{ps}d_{qr}^{}\overline{v}_{ps}\gamma ^\mu v_{qr}{\displaystyle \frac{(p+q)^\nu }{2}}`$
$`+e^{i(p+q)x}b_{ps}^{}d_{qr}^{}\overline{u}_{ps}\gamma ^\mu v_{qr}{\displaystyle \frac{(pq)^\nu }{2}}+e^{i(p+q)x}d_{ps}b_{qr}\overline{v}_{ps}\gamma ^\mu u_{qr}{\displaystyle \frac{(pq)^\nu }{2}}).`$
is formally conserved $`_\mu \theta ^{\mu \nu }=0`$ which allows us to formally define the four-momentum vector
$$P^\mu =\theta ^{0\mu }d^3x=\frac{d^3p}{2E_p}p^\mu \left(b_{ps}^{}b_{ps}d_{ps}d_{ps}^{}\right).$$
(44)
Again, tensor components $`\theta ^{i\mu }`$ mix positive and negative energy solutions giving the same *zitterbewegung*-like behavior to energy-momentum tensor as well. Current conservation gives us the time change of the energy
$$\frac{d}{dt}\underset{V}{}\theta ^{00}(x)d^3x=\underset{V}{}\underset{j=1}{\overset{3}{}}\frac{d}{dx_j}\theta ^{j,0}(x)d^3x=\underset{S}{}\underset{j=1}{\overset{3}{}}n^j\theta ^{j,0}(x)dS$$
(45)
which equals zero for infinite volume $`V`$ and field vanishing at infinity. However, if the volume is finite, then the surface integral has the *zitterbewegung*-like behavior having parts that oscillate both like $`e^{\pm iEt}`$ and $`e^{\pm i2Et}`$, while the left hand side has only the $`e^{\pm iEt}`$ oscillating part. The same holds for momentum components as well
$$\frac{d}{dt}\underset{V}{}\theta ^{0i}(x)d^3x=\underset{V}{}\underset{j=1}{\overset{3}{}}\frac{d}{dx_j}\theta ^{ji}(x)d^3x=\underset{S}{}\underset{j=1}{\overset{3}{}}n^j\theta ^{ji}(x)dS$$
(46)
which makes it hard to interpret the $`\theta ^{j0}`$ components as the components of Poynting vector or $`\theta ^{ji}`$ components as the components of stress tensor.
Note that the conserved current and energy-momentum four vector aren’t proportional (both before and after the second quantization) so it’s impossible to interpret the current as probability density-flux current.
Things get even worst with angular momentum density tensor; again, tensor is formally conserved $`_\mu J^{\mu ,\alpha \beta }=0`$, which again allows us to construct the Lorentz transformation generators
Jμν=J0,μνd3x=12ψ¯[(xμi
νxνi
)μγ0+{γ0,σμν2}]ψd3x.J^{\mu\nu}=\int J^{0,\mu\nu}\,d^{3}x=\int\frac{1}{2}\bar{\psi}\left[\left(x^{\mu}i\kern 1.66672pt\vbox{\hbox{\kern-1.66672pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}{}^{\nu}-x^{\nu}i\kern 1.66672pt\vbox{\hbox{\kern-1.66672pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}{}^{\mu}\right)\gamma^{0}+\left\{\gamma^{0},\frac{\sigma^{\mu\nu}}{2}\right\}\right]\psi\,d^{3}x\;. (47)
Coordinate part of Lorentz generators (also sometimes called "orbital" since it gives us the orbital angular momentum operator)
$`J_{\mathrm{𝑐𝑜𝑜𝑟𝑑}}^{\mu \nu }`$ $`=`$ 12ψ(xμi
νxνi
)μψd3x\displaystyle\int\frac{1}{2}\psi^{\dagger}\left(x^{\mu}i\kern 1.66672pt\vbox{\hbox{\kern-1.66672pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}{}^{\nu}-x^{\nu}i\kern 1.66672pt\vbox{\hbox{\kern-1.66672pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}{}^{\mu}\right)\psi\,d^{3}x (48)
$`=`$ $`{\displaystyle \frac{E_p}{m}\left(b_{ps}^{}\left(p^\mu i_p^\nu p^\nu i_p^\mu \right)b_{ps}d_{ps}\left(p^\mu i_p^\nu p^\nu i_p^\mu \right)d_{ps}^{}\right)\frac{d^3p}{2E_p}}.`$
transforms every spinor component independently corresponding to the transformation $`\psi (x)\psi (\mathrm{\Lambda }x)`$. It’s behavior is more or less reasonable (aside from the problem of normal ordering) although derivative operators acting on $`b_{ps}`$ and $`d_{ps}^{}`$ raise additional questions after second quantization. Spin part of generators
$`J_{\mathrm{𝑠𝑝𝑖𝑛}}^{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \overline{\psi }\{\gamma ^0,\frac{\sigma ^{\mu \nu }}{2}\}\psi d^3x}={\displaystyle \frac{1}{2}}{\displaystyle \psi ^{}\left(\frac{\sigma ^{\mu \nu }}{2}+\gamma ^0\frac{\sigma ^{\mu \nu }}{2}\gamma ^0\right)\psi }`$ (49)
doesn’t do so well. For space part we have $`\gamma ^0\sigma ^{ij}\gamma ^0=\sigma ^{ij}`$ so we get rotation generators
$`J_k`$ $``$ $`ϵ_{ijk}J^{ij}=ϵ_{ijk}{\displaystyle \psi ^{}\frac{\sigma ^{ij}}{2}\psi d^3x}`$ (50)
$`=`$ $`{\displaystyle }{\displaystyle \frac{d^3p}{(2\pi )^32E_p}}{\displaystyle \frac{1}{E_p}}{\displaystyle \underset{s,r}{}}(u_{ps}^{}{\displaystyle \frac{\sigma ^k}{2}}u_{pr}b_{ps}^{}b_{pr}+v_{ps}^{}{\displaystyle \frac{\sigma ^k}{2}}v_{pr}d_{ps}d_{pr}^{}`$
$`+e^{2iE_pt}u_{ps}^{}{\displaystyle \frac{\sigma ^k}{2}}v_{\stackrel{~}{p}r}b_{ps}^{}d_{\stackrel{~}{p}r}^{}+e^{2iE_pt}v_{ps}^{}{\displaystyle \frac{\sigma ^k}{2}}u_{\stackrel{~}{p}r}d_{ps}b_{\stackrel{~}{p}r})`$
which *again* shows zitterbewegung-like behavior not only in the space parts of the currents but in zeroth components as well. Those zitterbewegung-like terms
$$e^{2iE_pt}u_{ps}^{}\frac{\sigma ^k}{2}v_{\stackrel{~}{p}r}b_{ps}^{}d_{\stackrel{~}{p}r}^{}+e^{2iE_pt}v_{ps}^{}\frac{\sigma ^k}{2}u_{\stackrel{~}{p}r}d_{ps}b_{\stackrel{~}{p}r}$$
(51)
mix positive and negative energy states (aside from making generators time dependent which is clearly a contradiction to the idea of a time conserved quantity). After second quantization these terms (and therefore the whole generators) no longer annihilate the vacuum so the finite rotations
$`e^{i\stackrel{}{\omega }\stackrel{}{J}}`$ $`=`$ $`1i\stackrel{}{\omega }\stackrel{}{J}+(i\stackrel{}{\omega }\stackrel{}{J})^2+\mathrm{}`$ (52)
$``$ $`1+(b^{}b+d^{}d+d^{}b^{}+db)+(b^{}b+d^{}d+d^{}b^{}+db)^2+\mathrm{}`$
would either lead to
$$0\left|e^{i\stackrel{}{\omega }\stackrel{}{J}}\right|01$$
(53)
or if one insists $`0\left|e^{i\stackrel{}{\omega }\stackrel{}{J}}\right|0=1`$ it would lead to an infinite series of integral constraints on operators $`b`$ and $`d`$ with the only solution being the trivial one. Spontaneous symmetry breaking doesn’t help since it would lead to the vacuum with the preferred direction which would seem to contradict the experiment.
Another consequence of these terms is that they mix states with $`n`$ and $`n\pm 2`$ particles, for example vacuum and electron positron state
$$0\left|\stackrel{}{J}\right|e^+e^{}00\left|e^{i\stackrel{}{\omega }\stackrel{}{J}}\right|e^+e^{}0$$
(54)
which would imply that what looks like electron-positron pair from one angle looks like vacuum from the other.
For mixed space-time part of the tensor we have $`\gamma ^0\sigma ^{0i}\gamma ^0=\sigma ^{0i}`$, so the spin part of boost generator vanishes exactly:
$`K_k`$ $``$ $`J^{0k}={\displaystyle \frac{1}{2}}{\displaystyle \overline{\psi }\{\gamma ^0,\frac{\sigma ^{0k}}{2}\}\psi d^3x}={\displaystyle \frac{1}{4}}{\displaystyle \psi ^{}(\sigma ^{0k}+\gamma ^0\sigma ^{0k}\gamma ^0)\psi }=0`$ (55)
This would imply that every component of Dirac field transforms as scalar under boosts and as spinor under rotations which is a contradiction in definition.
Space part of the angular momentum density tensor $`J^{i,\alpha \beta }`$ again shows all the above problems and some more. Writing them down in detail wouldn’t be particularly illuminating and so it will be skipped.
These problems are present in other representations as well. While a second quantization, some reinterpretation of symbols and some renormalization a bit later do offer possible solutions for some of these problems, they don’t solve all of them. Although the approach used in this article does call for reinterpretation as well, it doesn’t merely offer solutions for these problems/peculiarities; none of them are *present* in the theory in the first place.
## 2 Relativity and quantum mechanics
Non-relativistic quantum mechanics is based on a few simple postulates. First one says that physical states are represented by a complete set of normalized, complex vectors $`\psi `$ in Hilbert space $``$:
$$\stackrel{}{\psi }{}_{}{}^{}\stackrel{}{\psi }1$$
(56)
where $`\stackrel{}{\psi }^{}`$ is defined to be complex conjugate of the vector $`\stackrel{}{\psi }`$. We can write the vector $`\stackrel{}{\psi }`$ (in some orthogonal basis) as a column matrix traditionally denoted as $`|\psi `$. Equation (56) then becomes
$$\psi |\psi 1$$
(57)
where $`\psi |`$ is now defined to be Hermitean conjugate of the matrix $`|\psi `$: $`\psi |\left(|\psi \right)^{}`$. If we are to interpret the square of the wave function as the probability, equation $`\psi (t)|\psi (t)1`$ should hold at all times. to ensure that, we require $`\left|\psi (\stackrel{}{x},t)\right|^2=\rho `$ to be the zeroth component of a conserved currend $`j^\mu =(\rho ,\stackrel{}{j})`$ with j=(i/m)ψ
ψ𝑗𝑖Planck-constant-over-2-pi𝑚superscript𝜓
𝜓\vec{j}=(-i\hbar/m)\psi^{*}\kern 1.66672pt\vbox{\hbox{\kern-1.00006pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\nabla$}}\hbox to-1.00006pt{}\psi
$$\frac{\rho (\stackrel{}{x},t)}{t}+\stackrel{}{j}(\stackrel{}{x},t)=_\mu j^\mu =0$$
(58)
If the system is in a state $`|\psi `$, then the probability of finding it in another state $`|\varphi `$ is $`\left|\varphi |\psi \right|^2`$. Physical observables are described by Hermitean operators $`A=A^{}`$ on the space $``$; the expectation value of operator $`A`$ is defined to be $`\psi \left|A\right|\psi `$. If the system is invariant under certain symmetries, a theorem of Wigner states that such symmetries are represented by unitary (or antiunitary) operators $`U^1=U^{}`$. If that wasn’t the case, then the first postulate *wouldn’t* be valid since if $`|\varphi =U|\psi `$, then
$$\varphi |\varphi =\psi \left|U^{}U\right|\psi $$
(59)
is no longer invariant under the symmetry. Wave function $`\psi (x)`$ can be represented as a superposition of orthogonal states $`\{\varphi _n\}`$
$$\psi (\stackrel{}{x},t)=\underset{n}{}c_n(t)\varphi _n(\stackrel{}{x},t)$$
(60)
where the set of functions $`\{\varphi _n\}`$ is a complete set
$$\underset{n}{}\varphi _n(\stackrel{}{x},t)\varphi _n^{}(\stackrel{}{y},t)\delta (\stackrel{}{x}\stackrel{}{y}).$$
(61)
In fact, the wave function $`\psi (x,t)`$ can be viewed as the factor in the decomposition of the abstract Hilbert space vector $`|\psi (t)`$ in the orthonormal basis of vectors $`|\stackrel{}{x}`$
$$|\psi (t)=|\stackrel{}{x}\underset{\psi (x,t)}{\underset{}{\stackrel{}{x}|\psi (t)}}d^3x$$
(62)
where the vectors $`|\stackrel{}{x}`$ are eigenvectors of the position operator $`\stackrel{}{𝐱}`$
$$\stackrel{}{𝐱}|\stackrel{}{x}=\stackrel{}{x}|\stackrel{}{x}.$$
(63)
This is the basis for the interpretation that the value of $`\psi (x,t)`$ is the amplitude (and therefore $`\left|\psi (x,t)\right|^2`$ the probability) for particle in the state $`|\psi (t)`$ to be found at the position $`\stackrel{}{x}`$. Position eigenstates can be eliminated completely from the equation (60) so it becomes
$$|\psi (t)=\underset{n}{}c_n(t)|\varphi _n(t)$$
(64)
where the coefficients $`c_n`$ are given by $`c_n=\varphi _n|\psi `$. This principle of superposition is the most fundamental principle in quantum mechanics. It is only natural to keep it through the rest of the article.
On the other hand, special theory of relativity requires that the speed of light $`c`$ is constant in all inertial frames, or mathematically the distance between two points in space-time
$$c^2(t_2^2t_1^2)(\stackrel{}{x}_2\stackrel{}{x}_1)^2$$
(65)
should be the same in all inertial frames. The group of transformations that obey this rule (Poincare group) consists of all the translations in space-time, rotations in space as well as of Lorentz transformations (or boosts). Here we run into trouble; while translations in space-time *are* unitary, general Lorentz transformations *aren’t*. The group of proper Lorentz transformations (usually called simply Lorentz group) is $`SO(1,3)`$ which is known to be non-compact and therefore has no finite-dimensional unitary representations. Using infinite-dimensional representations would imply that (free) particle of given energy and momentum has infinite number of (degenerate) spin states. This doesn’t seem to be the case in nature. So we’re stuck with non-unitary representations of Lorentz group.
This brings us to the fundamental incompatibility; since the representations are non-unitary, expressions like $`\psi ^{}|\phi ^{}\left(|\psi ^{}\right)^{}|\phi ^{}`$ won’t be invariant under Lorentz boosts to another frame
$$|\phi |\phi ^{}=e^{i\stackrel{}{\omega }\stackrel{}{K}}|\phi =U|\phi ,\psi |\psi ^{}|=\psi |\left(e^{i\stackrel{}{\omega }\stackrel{}{K}}\right)^{}=\psi |e^{+i\stackrel{}{\omega }(\stackrel{}{K})^{}}=\psi |U$$
(66)
$$\psi ^{}|\phi ^{}=\psi \left|U^{}U\right|\phi =\psi \left|UU\right|\phi \psi |\phi .$$
(67)
The same holds for expectation values of operators $`\psi \left|A\right|\psi `$. On the other hand, rotation as well as translation generators *are* Hermitean which makes rotations and translations unitary and scalar product invariant; this must not be changed whatever we do. QFT deals with this problem by keeping the definition of scalar product, reinterpreting equations as equations for fields, not wave functions, and postulating that these fields act on Fock space and create or destroy states in that space. Commutation or anticommutation relations are postulated in such a way that will keep the scalar product positive definite. However, it turns out that merely assigning the name "scalar product" to a different, relativistically invariant quantity will allow us to keep the idea of wave function and construct consistent single-particle relativistically invariant theory, without negative energies, without zitterbewegung, with nice conserved current proportional with energy and momentum density proportional to the zeroth component of this current, as we had in the non-relativistic case which enabled the famous probabilistic interpretation.
### 2.1 Lorentz invariant scalar product
The “new” definition of scalar product turns out to be quite familiar. Lets take another look at equation (65); in general, for complex four-vectors $`a`$, we have the invariant (and real) quantity
$$a_\mu ^{}a^\mu =a_0^{}a_0a_1^{}a_1a_2^{}a_2a_3^{}a_3.$$
(68)
If we use the Dirac’s bra-ket notation, four vector $`a^\mu `$ can be written as
$$a^\mu |a=\left(\begin{array}{c}a_0\\ a_1\\ a_2\\ a_3\end{array}\right).$$
(69)
Equation (68) can be written in matrix form as
$$a\left|P\right|a=a_0^{}a_0a_1^{}a_1a_2^{}a_2a_3^{}a_3.$$
(70)
where $`P`$ is the parity matrix
$$P=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right)$$
(71)
Since this is the invariant quantity in the defining representation of Lorentz group, it should be no surprise that it works in all another representations as well. In general, we can redefine the scalar product by inserting another operator between bra and ket states; to make scalar product invariant, one has to have
$$\psi ^{}\left|O\right|\phi ^{}=\psi \left|U^{}OU\right|\phi =\psi \left|O\right|\phi O^1U^{}O=U^1.$$
(72)
or in terms of generators, this operator must satisfy
$$O^1i\stackrel{}{J}O=i\stackrel{}{J}O^1i\stackrel{}{K}O=i\stackrel{}{K}.$$
(73)
If we take a look at the commutators of parity and Lorentz generators (appendix A)
$$P\stackrel{}{J}P=\stackrel{}{J}P\stackrel{}{K}P=\stackrel{}{K}.$$
(74)
and equation (70), it’s clear parity satisfies this condition. So far any coordinate dependence of states $`|\psi `$, $`|\phi `$ or $`|a`$ wasn’t mentioned. Operators above are in fact only spin parts of full generators so the operator $`O`$ which makes the scalar product invariant is in fact only the spin part of full parity operator. “Orbital” parts of Lorentz generators as well as translation generators act on the coordinate part of vectors
$$e^{\frac{i}{2}J_{coord}^{\mu \nu }\omega _{\mu \nu }iP^\rho a_\rho }\varphi _{\mathrm{}}(x)=\varphi _{\mathrm{}}(\mathrm{\Lambda }(\omega )x+a)$$
(75)
independently and don’t mix different spin components. Their products
$$\varphi ^{}(x)P_{spin}\varphi (x)$$
(76)
are function of position $`x`$ and as such in general are not Lorentz invariant for any representation. Their integrals, however, *will* be invariant if they are integrated with the proper Lorentz-invariant measure. Lorentz invariant scalar product will then be
$$\psi \left|P\right|\phi \overline{\psi }|\phi $$
(77)
where we define the “new” conjugate vector $`\overline{\psi }|\psi |P`$. This is in fact nothing new. We have already shown that contractions of covariant and contravariant vectors can be interpreted as parity operator sandwiched between two states. Another widely present example comes from Dirac representation; there Lorentz invariant product isn’t $`\psi ^{}(x)\psi (x)`$ but $`\overline{\psi }(x)\psi (x)`$, where bar on $`\psi `$ means $`\psi ^{}\gamma _0`$. Again, $`\gamma _0`$ is nothing but parity operator for $`(1/2,0)(0,1/2)`$ representation (see appendix B.1).
Note that equation (77) has one fundamental consequence: if we interpret the scalar product of the state in Hilbert space with itself $`\overline{\psi }|\psi `$ as a norm of that state, there will have to be some states with negative norm in every nontrivial representation. From mathematical viewpoint, such a definition isn’t really a norm; however, since in physics one uses the phrases like “negative metric” to describe Minkowski metric tensor, a phrase “norm” will also be misused here in the same tradition to describe the relativistically invariant scalar product of a vector with itself. Since the “norm” has a parity operator in it’s definition, it *need not* be positive. Therefore one has to think twice before calling a zeroth component of the continuity current probability density.
THis however isn’t inconsistent with the definition of the scalar product in nonrelativistic quantum mechanics. Since the parity operator isn’t uniquely defined by commutation relations (74), neither is the norm. If $`P`$ satisfies relations (74) so will $`P`$; therefore, if we choose parity operator to be $`P`$ instead of $`P`$, we have effectivlly multiplied the norm of all states with $`1`$. In another words, for given state $`|\psi `$ we can always choose parity operator in such a way that the norm of that particlar state is positive. As a consequence, in the non-relativistic limit where boosts no longer mix states of different spin
$$\psi (\stackrel{}{x},t)\psi ^{}(\stackrel{}{x}^{},t^{})=\psi (\stackrel{}{x}+\stackrel{}{v}t,t)$$
(78)
particle and antiparticle states transform separately under boosts and form a group of Galilean transformations. In this limit one can always choose the parity operator for different representations of Galilean group of transformations that will give positive definite norm for all spin states. Thus we can recover the proper non-relativistic scalar product as well as probabilistic interpretation.
What will finally fix parity operator is the requirement of it’s actions on particle states. For four-vector representations $`(1/2,1/2)`$ we have “physical” requirement that spin 1 (vector states) part should have negative parity and spin 0 part (scalar state) positive. Same logic works for all $`(j,j)`$ representations which have states with spin $`\mathrm{}(0,1,\mathrm{},2j)`$ having parity $`(1)^{\mathrm{}}`$.
## 3 Representations of Lorentz Group and parity equations
To have Lorentz invariant theory, besides the relativistically invariant scalar product, wave functions have to belong to various representations of Lorentz group. As it was mentioned before, any wave functions in a Hilbert space can always be decomposed in a complete set of functions. This fundamental property combined with the relativistic invariance will determine the behavior of wave functions for any given representation.
### 3.1 Construction of the spinor representation in RWFM
Lorentz boost and rotation generators for spinor representation are nothing else but Dirac $`\sigma `$ matrices. Since their derivation is straightforward only the results in chiral representation are quoted here (derivation can be found in appendix B.1)
$$J_k=\frac{1}{2}\left(\begin{array}{cc}\sigma _k& 0\\ 0& \sigma _k\end{array}\right)=\frac{1}{4}\underset{ij}{}ϵ_{ijk}\sigma ^{ij}K_k=\frac{i}{2}\left(\begin{array}{cc}\sigma _k& 0\\ 0& \sigma _k\end{array}\right)=\frac{1}{2}\sigma ^{0k}$$
(79)
Finite transformations are then generated by
$$S(\omega )=\mathrm{exp}\left(\frac{i}{2}J^{\mu \nu }\omega _{\mu \nu }\right)=\mathrm{exp}\left(\frac{i}{4}\sigma ^{\mu \nu }\omega _{\mu \nu }\right).$$
(80)
Now, let’s construct Lorentz-invariant momentum and spin eigenfunctions; wave function will be a direct sum of $`(1/2,0)`$ and $`(0,1/2)`$ terms
$$\psi (x)=\psi _R^{}(x)\psi _L^{}(x)=\left(\begin{array}{c}\psi _R^{}(x)\\ \psi _L^{}(x)\end{array}\right).$$
(81)
It is more convenient to add the null matrix to $`\psi _R^{}`$ and $`\psi _L^{}`$ matrices and make the sum in (81) normal instead of direct
$$\psi _R^{}\psi _R(x)=\left(\begin{array}{c}\psi _R^{}(x)\\ 0\end{array}\right),\psi _L^{}\psi _L(x)=\left(\begin{array}{c}0\\ \psi _L^{}(x)\end{array}\right),\psi (x)=\psi _R(x)+\psi _L(x)$$
(82)
Invariance conditions for translations and pure Lorentz transformations are
$$\psi (x)\psi ^{}(x^{})=\psi (x+a)\psi (x)\psi ^{}(x^{})=S(\omega )\psi (x).$$
(83)
Note that there is no $`S(a)`$ term for translations; it implies that each spin component transforms separately under translations. This will make it possible to factor out the same translation-generator-eigenfunctions (momentum eigenfunctions) from the whole wave function $`\psi (x)`$. At this point it’s convenient to (again) introduce Dirac’s braket notation. Wave function can then be written as
$$\psi (x)=\stackrel{}{x}|\psi (t)=\stackrel{}{x}|\psi _R(t)+\stackrel{}{x}|\psi _L(t)=\left(\begin{array}{c}\stackrel{}{x}|\psi _R^{}(t)\\ \stackrel{}{x}|\psi _L^{}(t)\end{array}\right).$$
(84)
We can insert a complete set of momentum and (time dependent) orthogonal spin eigenstates
$$d^3q\underset{r,𝒫}{}\frac{|\psi _{q,r,𝒫}\psi _{q,r,𝒫}|}{\psi _{q,r,𝒫}|\psi _{q,r,𝒫}}=d^3q\underset{r}{}\left(|\psi _{q,r,+}\psi _{q,r,+}||\psi _{q,r,}\psi _{q,r,}|\right)1$$
(85)
where we have acknowledged the fact that negative parity states have negative norm as well. Wave function then becomes
$`\stackrel{}{x}|\psi (t)`$ $`=`$ $`{\displaystyle d^3q\underset{r,𝒫}{}\frac{\stackrel{}{x}|\psi _{q,r,𝒫}\psi _{q,r,𝒫}|\psi }{\psi _{q,r,𝒫}|\psi _{q,r,𝒫}}}`$ (86)
$`=`$ $`{\displaystyle d^3q\underset{r,𝒫}{}\left(\stackrel{}{x}|\psi _{q,r,+}\psi _{q,r,+}|\psi \stackrel{}{x}|\psi _{q,r,}\psi _{q,r,}|\psi \right)}`$
Translational invariance tells us that spin and coordinate dependence can be factored
$$\stackrel{}{x}|\psi _{q,r,𝒫}(t)=\frac{e^{i\stackrel{}{q}\stackrel{}{x}}}{(2\pi )^{3/2}}w_{q,r,𝒫}(t)$$
(87)
where $`w_{q,r,𝒫}`$ is some matrix in spin space which in general can depend on momentum and parity of the state. Sign in the exponential is chosen to make the eigenvalue of momentum operator “positive” $`\stackrel{}{q}`$
$$i\stackrel{}{x}|\psi _{q,r,𝒫}(t)=\stackrel{}{q}\stackrel{}{x}|\psi _{q,r,𝒫}(t).$$
(88)
For momentum and spin eigenstate $`|\psi (t)=|\psi _{p,s,𝒫}`$ we have
$$\stackrel{}{x}|\psi _{p,s,𝒫}(t)=d^3q\underset{r,𝒫^{}}{}\frac{e^{i\stackrel{}{q}\stackrel{}{x}}}{(2\pi )^{3/2}}w_{q,r,𝒫}\psi _{q,r,𝒫^{}}|\psi _{p,s,𝒫}=\frac{e^{i\stackrel{}{p}\stackrel{}{x}}}{(2\pi )^{3/2}}w_{q,r,𝒫}(t)$$
(89)
with the states normalized (at equal time) as
$$\psi _{q,r,𝒫^{}}|\psi _{p,s,𝒫}=𝒫\delta _{\stackrel{}{p},\stackrel{}{q}}\delta _{r,s}\delta _{𝒫^{},𝒫}.$$
(90)
Since translations in space (and time) don’t mix different spin components, we can factor out the common $`x`$-dependent exponential. In $`(1/2,0)(0,1/2)`$ basis we have
$$\stackrel{}{x}|\psi _{p,s}(t)=\stackrel{}{x}|\psi _{p,s}^R(t)+\stackrel{}{x}|\psi _{p,s}^L(t).$$
(91)
Now, parity will transform states with momentum $`\stackrel{}{p}`$ and spin $`s`$ to the state with momentum $`\stackrel{}{p}`$ and spin $`s`$; we can divide the coordinate and spin part $`P=P_{coord}P_{spin}`$ so that the coordinate part acts on $`|\stackrel{}{x}`$ while spin part acts on $`|\psi _{p,s}`$
$$P\stackrel{}{x}|\psi _{p,s,𝒫}(t)=P_{coord}\frac{e^{i\stackrel{}{p}\stackrel{}{x}}}{(2\pi )^{3/2}}P_{spin}w_{q,r,𝒫}=\frac{e^{i\stackrel{}{p}\stackrel{}{x}}}{(2\pi )^{3/2}}P_{spin}w_{q,r,𝒫}(t)$$
(92)
For massive particles there is always a nontrivial momentum eigenstate with $`\stackrel{}{p}=0`$; this eigenstate has to be also parity eigenstate
$$P\stackrel{}{x}|\psi _{0,s}(t)=\pm \stackrel{}{x}|\psi _{0,s}(t)$$
(93)
On the other hand, parity just exchanges $`\psi _R`$ and $`\psi _L`$ so obviously, parity eigenstates in the rest frame will be states with equal $`\psi _R`$ and $`\psi _L`$ with either the same or different relative sign
$`\stackrel{}{x}|\psi _{0,s}^+(t)`$ $`=`$ $`\stackrel{}{x}|\psi _{0,s}^R(t)+\stackrel{}{x}|\psi _{0,s}^L(t)={\displaystyle \frac{1}{(2\pi )^{3/2}}}\underset{1}{\underset{}{e^{\stackrel{}{x}\stackrel{}{0}}}}u_{0,s}(t)`$
$`\stackrel{}{x}|\psi _{0,s}^{}(t)`$ $`=`$ $`\stackrel{}{x}|\psi _{0,s}^R(t)\stackrel{}{x}|\psi _{0,s}^L(t)={\displaystyle \frac{1}{(2\pi )^{3/2}}}\underset{1}{\underset{}{e^{\stackrel{}{x}\stackrel{}{0}}}}v_{0,s}(t)`$ (94)
where matrices $`u`$ and $`v`$ are defined to be
$$u_{0,s}(t)=\left(\begin{array}{c}\chi _s(t)\\ \chi _s(t)\end{array}\right)v_{0,s}(t)=\left(\begin{array}{c}\chi _s(t)\\ \chi _s(t)\end{array}\right)$$
(95)
Translation invariance in time requires that the time dependence of matrices $`\chi _s(t)`$ can be factored to
$$\chi _s(t)=e^{\pm i\kappa t}\chi _s$$
(96)
where $`\kappa `$ is a real positive number, and $`\chi _s`$ time independent bispinor. Parity eigenstates then become
$$u_{0,s}(t)=e^{\pm i\kappa t}\left(\begin{array}{c}\chi _s\\ \chi _s\end{array}\right)=e^{\pm i\kappa t}u_{0,s}v_{0,s}(t)=e^{\pm i\kappa t}\left(\begin{array}{c}\chi _s\\ \chi _s\end{array}\right)=e^{\pm i\kappa t}v_{0,s}.$$
(97)
Spinors for finite momentum can be obtained from (97) by applying Lorentz boost in direction $`\stackrel{}{\theta }`$. Boost operator can again be decomposed into parts acting only on coordinates only and the part acting only on spin degrees of freedom.
$$\stackrel{}{x}^{}|\psi _{p,s}^+=S(\stackrel{}{\theta })\stackrel{}{x}|\psi _{0,s}^+=S(\stackrel{}{\theta })\left[\frac{e^{\pm i\kappa t}}{(2\pi )^{3/2}}u_{0,s}\right]=S_{coord}(\stackrel{}{\theta })\frac{e^{\pm i\kappa t}}{(2\pi )^{3/2}}S_{spin}(\stackrel{}{\theta })u_{0,s}=S_{coord}(\stackrel{}{\theta })\frac{e^{\pm i\kappa t}u_{ps}}{(2\pi )^{3/2}}$$
(98)
where we define spinors $`u_{ps}`$ and $`v_{ps}`$ to be the spinor obtained by boosting the rest frame spinors $`u_{0,s}`$ and $`v_{0,s}`$ to a frame where it has the momentum $`\stackrel{}{p}`$
$$u_{ps}S(\stackrel{}{\theta })u_{0,s},v_{ps}S(\stackrel{}{\theta })v_{0,s}.$$
(99)
If we *require* that the resulting state has the momentum $`\stackrel{}{p}`$ parallel to $`\stackrel{}{\theta }`$, we can again on the grounds of translational invariance in space-time conclude
$$\stackrel{}{x}^{}|\psi _{ps}^+(t^{})=e^{\pm i\kappa ^{}t^{}}\stackrel{}{x}^{}|\psi _{ps}^+=\frac{e^{i\stackrel{}{p}\stackrel{}{x}^{}}}{(2\pi )^{3/2}}e^{\pm i\kappa ^{}t^{}}u_{ps}$$
(100)
Comparing equations (98) and (100) we get
$$S_{coord}(\stackrel{}{p})e^{\pm i\kappa t}=e^{i\stackrel{}{p}\stackrel{}{x}^{}}e^{\pm i\kappa ^{}t^{}}.$$
(101)
To get the factor $`e^{i\stackrel{}{p}\stackrel{}{x}^{}}`$ on the right hand side, sign in front of $`\kappa `$ must be positive and equal to particle energy in the rest frame, or in another hands $`\kappa =m`$. Then on the other side we have $`\kappa ^{}=E_\stackrel{}{p}=+\sqrt{\stackrel{}{p}^2+m^2}`$.
$$\stackrel{}{x}|\psi _{ps}^+(t)=\frac{e^{ipx}}{(2\pi )^{3/2}}u_{ps}$$
(102)
This however determines the behavior of $`\stackrel{}{x}|\psi _{ps}^{}(t)`$ completely as well since they are both superpositions of *same* chiral wave functions
$$\stackrel{}{x}|\psi _{ps}^\pm (t)=\stackrel{}{x}|\psi _{ps}^R(t)\pm \stackrel{}{x}|\psi _{ps}^L(t)$$
(103)
and therefore must have the same space-time dependent exponential
$$\stackrel{}{x}|\psi _{ps}^{}(t)=\frac{e^{ipx}}{(2\pi )^{3/2}}v_{ps}.$$
(104)
The conclusion that both energies are positive is based on the requirement that the state boosted by $`\stackrel{}{\theta }`$ has the momentum parallel to $`\stackrel{}{\theta }`$. While this seams reasonable requirement, there’s nothing preventing us to require it to be antiparallel. That would lead to negative energies, and negative for all four solutions. But there’s no consistent way to have solutions with opposite energies as solutions of Dirac equation (1) do.
### 3.2 Parity equations for spinor representation
Zero momentum eigenfunctions satisfy
$`P\stackrel{}{x}|\psi _{0,s}^+(t)=\stackrel{}{x}|\psi _{0,s}^+(t)P\stackrel{}{x}|\psi _{0,s}^{}(t)=\stackrel{}{x}|\psi _{0,s}^{}(t)`$ (105)
or for spinors
$$P_{spin}u_{0,s}=\gamma ^0u_{0,s}=u_{0,s},P_{spin}v_{0,s}=\gamma ^0v_{0,s}=v_{0,s}.$$
(106)
In analogy with vector properties under parity, we’ll call solutions $`u_{ps}`$ axial spinors or pseudospinors, and solutions $`v_{ps}`$ polar spinors.
Now we may ask the question: in our frame these are zero momentum spinors; mirroring them to a fixed point reproduces themselves multiplied with $`\pm 1`$. In our frame this is parity operation; however, observer in a frame moving with some velocity will neither see zero momentum particle eigenstate nor will mirroring that state to a point moving with the frame velocity look like parity to him. So how will that symmetry operation look to him?
At this point it’s convenient to apply unitary transformation to all matrices in spinor space which will make the parity operator diagonal and leave spin generators $`\stackrel{}{J}`$ unchanged
$$M_{ch}M_D=U^{}M_{ch}U\psi _D=U^{}\psi _{ch}U=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)$$
(107)
This gives us the original set of matrices used by Dirac to describe spinors. Spin eigenstates in rest frame are then
$$u_{0,s}=\left(\begin{array}{c}\chi _s\\ 0\end{array}\right),v_{0,s}=\left(\begin{array}{c}0\\ \chi _s\end{array}\right).$$
(108)
Boosting equation (106) to frame where particles have momentum $`\stackrel{}{p}`$ we have
$$S(\stackrel{}{p})\gamma ^0S^1(\stackrel{}{p})S(\stackrel{}{p})u_{0,s}=S(\stackrel{}{p})u_{0,s}S(\stackrel{}{p})\gamma ^0S^1(\stackrel{}{p})S(\stackrel{}{p})v_{0,s}=S(\stackrel{}{p})v_{0,s}.$$
(109)
After evaluating
$$S(\stackrel{}{p})PS^1(\stackrel{}{p})=S(\stackrel{}{p})\gamma ^0S^1(\stackrel{}{p})=\left(\begin{array}{cc}\frac{E}{m}& \frac{\stackrel{}{p}\stackrel{}{\sigma }}{m}\\ \frac{\stackrel{}{p}\stackrel{}{\sigma }}{m}& \frac{E}{m}\end{array}\right)$$
(110)
In Dirac representation one can introduce $`\gamma `$ matrices and express this in closed form as
$$S(\stackrel{}{p})PS^1(\stackrel{}{p})=\frac{p/}{m}$$
(111)
so equation (109) becomes
$$\frac{p/}{m}u_{ps}=u_{ps}\frac{p/}{m}v_{ps}=v_{ps}$$
(112)
or
$$(p/m)u_{ps}=0(p/+m)v_{ps}=0.$$
(113)
Making a superposition of these momentum eigenstates
$`\psi ^+(x)`$ $``$ $`{\displaystyle d^3p\underset{s}{}\stackrel{}{x}|u_{ps}u_{ps}|\psi ^+}={\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_p}\underset{s}{}e^{ipx}u_{ps}b_{ps}}=\left(\begin{array}{c}\phi ^+(x)\\ \chi ^+(x)\end{array}\right)`$ (116)
$`\psi ^{}(x)`$ $``$ $`{\displaystyle d^3p\underset{s}{}\stackrel{}{x}|v_{ps}v_{ps}|\psi ^+}={\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_p}\underset{s}{}e^{ipx}v_{ps}d_{ps}}=\left(\begin{array}{c}\phi ^{}(x)\\ \chi ^{}(x)\end{array}\right)`$ (119)
one gets parity equations for spin 1/2
$$(/m)\psi ^+(x)=0,(/+m)\psi ^{}(x)=0.$$
(120)
Quantity $`b_{ps}/2E_p=u_{ps}|\psi ^+`$ can be interpretedas amplitude for finding the state $`|\psi ^+`$ in momentum-spin eigenstate $`|u_{ps}`$. Same holds for $`d_{ps}/2E_p=v_{ps}|\psi ^{}`$.
First equation in (120) is what Dirac proposed as the relativistic equation of the first order which he hoped wouldn’t have negative energies. Since the spinor representation has 4 independent solutions, if polar spinors were to satisfy the Dirac equation, one artificially had to multiply the spin part of solution for polar spinors with coordinate factor $`\mathrm{exp}(+ipx)`$ instead of $`\mathrm{exp}(ipx)`$. In QFT this lead to solutions with negative energies which are clearly not present here.
In terms of $`x`$-space bispinors $`\phi ^\pm (x)`$ and $`\chi ^\pm (x)`$ parity equations give us a set of coupled first order equations
$`i{\displaystyle \frac{\phi ^\pm (x)}{t}}`$ $`=`$ $`\pm m\phi ^\pm (x)i\stackrel{}{\sigma }\chi ^\pm (x)`$ (121)
$`i{\displaystyle \frac{\chi ^\pm (x)}{t}}`$ $`=`$ $`m\chi ^\pm (x)i\stackrel{}{\sigma }\phi ^\pm (x).`$ (122)
Complete wave function for $`(1/2,0)(0,1/2)`$ will be the sum of polar and axial part
$`\psi (x)`$ $``$ $`\stackrel{}{x}|\psi =\stackrel{}{x}\left|\left({\displaystyle d^3p\underset{s,𝒫}{}\frac{|w_{p,s,𝒫}w_{p,s,𝒫}|}{w_{p,s,𝒫}|w_{p,s,𝒫}}}\right)\right|\psi `$ (127)
$`=`$ $`{\displaystyle d^3p\underset{s}{}\left(\stackrel{}{x}|u_{p,s}u_{ps}|\psi \stackrel{}{x}|v_{p,s}v_{ps}|\psi \right)}`$
$`=`$ $`{\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_p}\underset{s}{}e^{ipx}\left(b_{ps}u_{ps}+d_{ps}v_{ps}\right)}=\left(\begin{array}{c}\phi ^+(x)+\phi ^{}(x)\\ \chi ^+(x)+\chi ^{}(x)\end{array}\right)=\left(\begin{array}{c}\phi (x)\\ \chi (x)\end{array}\right)`$
in the massless limit $`m0`$ satisfies the set of equations
$`{\displaystyle \frac{\phi (x)}{t}}=\stackrel{}{\sigma }\chi (x){\displaystyle \frac{\chi (x)}{t}}=\stackrel{}{\sigma }\phi (x).`$ (128)
Independent solutions to this system of equations satisfy $`\phi (x)=\pm \chi (x)`$ and correspond to separate $`(1/2,0)`$ and $`(0,1/2)`$ transformations. Interpretation of this result is given in section 5.
This same argument about parity in different frames applies to all representations of Lorentz group. The fact that equations can be expressed in relativistically-covariant way and that the equations are linear in both energy and momentum (or alternatively both time and space derivatives) are unique propery of Dirac representation since only the anticommutators of generators in Dirac representation satisfy Clifford algebra as well as usual $`SU(2)`$-generators algebra.
Another thing worth mentioning here is the fact that these differential equations are a consequence of the relativistic mixing of space and time. Applying the same arguments to representations of Galilean group will not have any time derivatives since the time is the same in all Galilean frames.
### 3.3 Parity equations for spin 1 representations
Now lets see what parity symmetry yields for higher spin representations. There are two representations that have eigenstates with spin 1, $`(1,0)(0,1)`$ and $`(1/2,1/2)`$ . For $`(1,0)(0,1)`$ representation transformation matrix $`S(\theta )`$ is given by
$$S(\theta )=e^{i\stackrel{}{\theta }\stackrel{}{K}}=\mathrm{exp}\left(\begin{array}{cc}\stackrel{}{\theta }\stackrel{}{S}& 0\\ 0& \stackrel{}{\theta }\stackrel{}{S}\end{array}\right)$$
(129)
where matrices $`S`$ are given explicitly in (391) or (396). The exponential of $`\pm \stackrel{}{\theta }\stackrel{}{S}`$ in spin representation
$$\pm \stackrel{}{\theta }\stackrel{}{S}=\pm \left(\begin{array}{ccc}\theta _0& \theta _{}& 0\\ \theta _+& 0& \theta _{}\\ 0& \theta _+& \theta _0\end{array}\right)\theta _\pm =\frac{\theta _1\pm \theta _2}{\sqrt{2}},\theta _0=\theta _3$$
(130)
or in coordinate representation
$$\pm \stackrel{}{\theta }\stackrel{}{S}=\pm \left(\begin{array}{ccc}0& \theta _3& \theta _2\\ \theta _3& 0& \theta _1\\ \theta _2& \theta _3& 0\end{array}\right)$$
(131)
calculated explicitly yields
$$e^{\pm \stackrel{}{\theta }\stackrel{}{S}}=\mathrm{𝟏}\pm \frac{\mathrm{sinh}\theta }{\theta }\stackrel{}{\theta }\stackrel{}{S}+\frac{\mathrm{cosh}\theta 1}{\theta ^2}\left(\stackrel{}{\theta }\stackrel{}{S}\right)^2$$
(132)
or
$$e^{i\stackrel{}{\theta }\stackrel{}{K}}=\mathrm{𝟏}+\frac{\mathrm{sinh}\theta }{\theta }\left(i\stackrel{}{\theta }\stackrel{}{K}\right)+\frac{\mathrm{cosh}\theta 1}{\theta ^2}\left(i\stackrel{}{\theta }\stackrel{}{K}\right)^2$$
(133)
where we have used the fact that $`\left(\stackrel{}{\theta }\stackrel{}{S}\right)^3=\theta ^2\left(\stackrel{}{\theta }\stackrel{}{S}\right)`$. Repeating the same procedure for $`(1/2,1/2)`$ representation yields formally the same result
$$S(\theta )=e^{i\stackrel{}{\theta }\stackrel{}{K}}=\mathrm{𝟏}+\frac{\mathrm{sinh}\theta }{\theta }\left(i\stackrel{}{\theta }\stackrel{}{K}\right)+\frac{\mathrm{cosh}\theta 1}{\theta ^2}\left(i\stackrel{}{\theta }\stackrel{}{K}\right)^2$$
(134)
but with different set of generators $`\stackrel{}{K}`$. In spin representation one has
$$i\stackrel{}{\theta }\stackrel{}{K}=\left(\begin{array}{cccc}0& \theta _+& \theta _0& \theta _{}\\ \theta _{}& & & \\ \theta _0& & \mathrm{𝟎}& \\ \theta _+& & & \end{array}\right)$$
(135)
while in coordinate representation this is
$$i\stackrel{}{\theta }\stackrel{}{K}=\left(\begin{array}{cccc}0& \theta _1& \theta _2& \theta _3\\ \theta _1& & & \\ \theta _2& & \mathrm{𝟎}& \\ \theta _3& & & \end{array}\right).$$
(136)
Since
$`\mathrm{cosh}\theta `$ $`=`$ $`{\displaystyle \frac{E}{m}}\mathrm{sinh}\theta ={\displaystyle \frac{\left|\stackrel{}{p}\right|}{m}}{\displaystyle \frac{\stackrel{}{\theta }\stackrel{}{S}}{\theta }}={\displaystyle \frac{\stackrel{}{p}\stackrel{}{S}}{\left|\stackrel{}{p}\right|}}`$
$`\mathrm{cosh}{\displaystyle \frac{\theta }{2}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{E+m}{2m}}}\mathrm{sinh}{\displaystyle \frac{\theta }{2}}=\sqrt{{\displaystyle \frac{Em}{2m}}}`$ (137)
we can express both equations as
$$S(\theta )=e^{i\stackrel{}{\theta }\stackrel{}{K}}=\mathrm{𝟏}+\frac{1}{m}\left(i\stackrel{}{p}\stackrel{}{K}\right)+\frac{1}{m(E+m)}\left(i\stackrel{}{p}\stackrel{}{K}\right)^2$$
(138)
with proper interpretation of generators $`\stackrel{}{K}`$. Using this to transform parity we get
$$S(\theta )PS(\theta )=e^{i\stackrel{}{\theta }\stackrel{}{K}}Pe^{i\stackrel{}{\theta }\stackrel{}{K}}=P\left(1\frac{2E}{m^2}\left(i\stackrel{}{p}\stackrel{}{K}\right)+\frac{2E}{m^2}\left(i\stackrel{}{p}\stackrel{}{K}\right)^2\right)$$
(139)
where we have used the fact that $`P\stackrel{}{K}=\stackrel{}{K}P`$ and $`\left(i\stackrel{}{p}\stackrel{}{K}\right)^3=p^2\left(i\stackrel{}{p}\stackrel{}{K}\right)`$. Switching to parity basis for $`(1,0)(0,1)`$ representation, equation (139) in block form becomes
$$S(\theta )PS(\theta )=\left(\begin{array}{cc}1+\frac{2}{m^2}(\stackrel{}{p}\stackrel{}{S})^2& \frac{2E}{m^2}(\stackrel{}{p}\stackrel{}{S})\\ \frac{2E}{m^2}(\stackrel{}{p}\stackrel{}{S})& \left(1+\frac{2}{m^2}(\stackrel{}{p}\stackrel{}{S})^2\right)\end{array}\right)$$
(140)
Using the same notation for parity eigenstates as for $`(1/2,0)(0,1/2)`$ representation
$$u_{ps}\psi _{ps}^+=\left(\begin{array}{cc}\phi _{ps}^+& \\ \chi _{ps}^+& \end{array}\right)v_{ps}\psi _{ps}^{}=\left(\begin{array}{cc}\phi _{ps}^{}& \\ \chi _{ps}^{}& \end{array}\right)$$
(141)
where $`\phi `$ and $`\chi `$ are now matrices with three rows, parity condition for particle in the rest frame
$$Pu_{0,s}=u_{0,s},Pv_{0,s}=v_{0,s}$$
(142)
boosted to a frame where particle has momentum $`\stackrel{}{p}`$ becomes
$$e^{i\stackrel{}{\theta }\stackrel{}{K}}Pe^{+i\stackrel{}{\theta }\stackrel{}{K}}e^{i\stackrel{}{\theta }\stackrel{}{K}}\psi _{0,s}^\pm =\pm e^{i\stackrel{}{\theta }\stackrel{}{K}}\psi _{0,s}^\pm .$$
(143)
Since by definition
$$\psi _{ps}^+u_{ps}=e^{i\stackrel{}{\theta }\stackrel{}{K}}\psi _{0,s}^+=e^{i\stackrel{}{\theta }\stackrel{}{K}}u_{0,s}\psi _{ps}^{}v_{ps}=e^{i\stackrel{}{\theta }\stackrel{}{K}}\psi _{0,s}^{}=e^{i\stackrel{}{\theta }\stackrel{}{K}}v_{0,s}$$
(144)
we get the set of equations
$`\left[{\displaystyle \frac{11}{2}}m^2+\left(\stackrel{}{p}\stackrel{}{S}\right)^2\right]\phi _{ps}^\pm +\left(\stackrel{}{p}\stackrel{}{S}\right)E\chi _{ps}^\pm `$ $`=`$ $`0`$ (145)
$`\left(\stackrel{}{p}\stackrel{}{S}\right)E\phi _{ps}^\pm \left[{\displaystyle \frac{1\pm 1}{2}}m^2+\left(\stackrel{}{p}\stackrel{}{S}\right)^2\right]\chi _{ps}^\pm `$ $`=`$ $`0`$ (146)
Following the same logic as in the case of Dirac representation, we can construct the wave function $`\psi ^+`$ and $`\psi ^{}`$ of parts with positive and negative parity by making the superposition of a complete set of eigenstates of given parity
$`\psi ^+(x)`$ $``$ $`{\displaystyle d^3p\underset{s}{}\stackrel{}{x}|u_{ps}u_{ps}|\psi ^+}={\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_p}\underset{s}{}e^{ipx}b_{ps}u_{ps}}\left(\begin{array}{c}\phi ^+(x)\\ \chi ^+(x)\end{array}\right)`$ (149)
$`\psi ^{}(x)`$ $``$ $`{\displaystyle d^3p\underset{s}{}\stackrel{}{x}|v_{ps}v_{ps}|\psi ^{}}={\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_p}\underset{s}{}e^{ipx}d_{ps}v_{ps}}\left(\begin{array}{c}\phi ^{}(x)\\ \chi ^{}(x)\end{array}\right).`$ (152)
Since $`\stackrel{}{S}`$ in coordinate basis can be expressed as $`S_{jk}^i=iϵ^{ijk}`$, identifying components of column-matrix $`f`$ with components of vector
$$\stackrel{}{f}=\left(\begin{array}{ccc}f_1& & \\ f_2& & \\ f_3& & \end{array}\right)$$
(153)
product $`(\stackrel{}{p}\stackrel{}{S})_{jk}=ip^iϵ^{ijk}`$ can be expressed as
$$(\stackrel{}{p}\stackrel{}{S})f=ip^iϵ^{ijk}f^k=i(\stackrel{}{p}\times \stackrel{}{f})_j=(\times \stackrel{}{f})_j$$
(154)
which gives us coupled differential equations for matrices $`\phi ^\pm (x)`$ and $`\chi ^\pm (x)`$ representing regular 3-vectors $`\stackrel{}{\phi }^\pm (x)`$ and $`\stackrel{}{\chi }^\pm (x)`$
$`{\displaystyle \frac{11}{2}}m^2\stackrel{}{\phi }^\pm (x)+\times \left(\times \stackrel{}{\phi }^\pm (x)\right)+\times \left({\displaystyle \frac{\stackrel{}{\chi }^\pm (x)}{t}}\right)`$ $`=`$ $`0`$ (155)
$`\times \left({\displaystyle \frac{\stackrel{}{\phi }^\pm (x)}{t}}\right){\displaystyle \frac{1\pm 1}{2}}m^2\stackrel{}{\chi }^\pm (x)\times \left(\times \stackrel{}{\phi }^\pm (x)\right)`$ $`=`$ $`0`$ (156)
In the limit $`m0`$ mass term becomes negligible, and the total wave function
$$\psi (x)=\psi ^+(x)+\psi ^{}(x)=\left(\begin{array}{cc}\stackrel{}{\phi }{}_{}{}^{+}(x)+\stackrel{}{\phi }{}_{}{}^{}(x)& \\ \stackrel{}{\chi }{}_{}{}^{+}(x)+\stackrel{}{\chi }{}_{}{}^{}(x),& \end{array}\right)\left(\begin{array}{cc}i\stackrel{}{b}(x)& \\ \stackrel{}{e}(x),& \end{array}\right)$$
(157)
satisfies equations
$`i\times \left(\times \stackrel{}{b}(x){\displaystyle \frac{\stackrel{}{e}(x)}{t}}\right)`$ $`=`$ $`0`$
$`\times \left(\times \stackrel{}{e}(x)+{\displaystyle \frac{\stackrel{}{b}(x)}{t}}\right)`$ $`=`$ $`0`$ (158)
where we have changed the notation to emphasize the similarity with Maxwell equations in vacuum. For very small but non-vanishing mass $`m`$, right-hand side of equations (3.3) will be proportional to $`m^2`$ and so the corrections to (157) will be of order $`(m/E)^2`$. Taking this mass to be bellow the experimental limit for the mass of the photon gives corrections suppressed by about 40 orders of magnitude, far beyond any experimental detection.
For $`(1/2,1/2)`$ representation, generators $`\stackrel{}{K}`$ to be inserted in (139) in coordinate representation are given by
$$i\stackrel{}{p}\stackrel{}{K}=\left(\begin{array}{cccc}0& p_1& p_2& p_3\\ p_1& & & \\ p_2& & \mathrm{𝟎}& \\ p_3& & & \end{array}\right).$$
(159)
Parity states
$$\psi _{ps}^\pm e^{i\stackrel{}{\theta }\stackrel{}{K}}\psi _{0,s}^\pm $$
(160)
obtained by boosting the rest-frame eigenstates
$$\psi _{0,s}^+=u_{0,s}=\left(\begin{array}{cc}\psi _{0,s}^0& \\ 0& \end{array}\right)\psi _{0,s}^{}=v_{0,s}=\left(\begin{array}{cc}0& \\ \stackrel{}{\psi }_{0,s}& \end{array}\right).$$
(161)
can be identified (again in coordinate representation) with components of four vector $`\psi ^\mu =(\psi ^0,\stackrel{}{\psi })`$
$$\psi _{ps}^\pm =\left(\begin{array}{cc}\psi _{ps}^{\pm \mathrm{\hspace{0.17em}0}}& \\ \stackrel{}{\psi }_{ps}^\pm & \end{array}\right).$$
(162)
Action of generators $`\stackrel{}{K}`$ on these states in terms of “regular” components of four-vector $`\psi ^\mu `$ can be expressed as
$`i\stackrel{}{p}\stackrel{}{K}\psi (x)`$ $`=`$ $`\left(\begin{array}{cc}\stackrel{}{p}\stackrel{}{\psi }(x)& \\ \stackrel{}{p}\psi _0(x)& \end{array}\right)=\left(\begin{array}{cc}i\stackrel{}{\psi }(x)& \\ i\psi _0(x)& \end{array}\right)`$ (167)
$`\left(i\stackrel{}{p}\stackrel{}{K}\right)^2\psi (x)`$ $`=`$ $`\left(\begin{array}{cc}\stackrel{}{p}^2\psi _0(x)& \\ \stackrel{}{p}\left(\stackrel{}{p}\stackrel{}{\psi }(x)\right)& \end{array}\right)=\left(\begin{array}{cc}^2\psi _0(x)& \\ \left(\stackrel{}{\psi }(x)\right)& \end{array}\right).`$ (172)
Wave function constructed as the superposition of eigenstates
$`A^+(x)`$ $``$ $`{\displaystyle d^3p\underset{s}{}\stackrel{}{x}|u_{ps}u_{ps}|\psi ^+}={\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_p}\underset{s}{}e^{ipx}b_{ps}u_{ps}}\left(\begin{array}{c}A_0^+(x)\\ \stackrel{}{A}^+(x)\end{array}\right)`$ (175)
$`A^{}(x)`$ $``$ $`{\displaystyle d^3p\underset{s}{}\stackrel{}{x}|v_{ps}v_{ps}|\psi ^{}}={\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_p}\underset{s}{}e^{ipx}d_{ps}v_{ps}}\left(\begin{array}{c}A_0^{}(x)\\ \stackrel{}{A}^{}(x)\end{array}\right).`$ (178)
again satisfies a set of coupled differential equations
$`\left({\displaystyle \frac{11}{2}}m^2^2\right)A_0^\pm (x){\displaystyle \frac{}{t}}\stackrel{}{A}^\pm (x)`$ $`=`$ $`0`$ (179)
$`{\displaystyle \frac{}{t}}A_0^\pm (x){\displaystyle \frac{1\pm 1}{2}}m^2\stackrel{}{A}^\pm (x)\left(\stackrel{}{A}^\pm (x)\right)`$ $`=`$ $`0`$ (180)
Looking again either at massless limit or ultra-relativistic regime, mass terms can be neglected compared do other terms and we end up with familiar equations for the total wave function $`A^\mu (x)`$
$`\left({\displaystyle \frac{\stackrel{}{A}(x)}{t}}+A_0(x)\right)`$ $`=`$ $`0`$ (181)
$`\left({\displaystyle \frac{A_0(x)}{t}}+\stackrel{}{A}(x)\right)`$ $`=`$ $`\left(_\mu A^\mu (x)\right)=0.`$ (182)
If one tries to construct $`(1,0)(0,1)`$ representation from two $`(1/2,1/2)`$ representations, $`k^\mu `$ and $`A^\mu `$, one gets electric and magnetic fields $`\stackrel{}{E}=\stackrel{}{A}/tA_0`$, $`\stackrel{}{B}=\times \stackrel{}{A}`$. First equation in (181) then becomes
$$\stackrel{}{E}(x)=0$$
(183)
while the last of Maxwell equations (in vacuum) $`\stackrel{}{B}(x)=0`$ is satisfied automatically. So one gets Maxwell equations plus the gauge condition $`_\mu A^\mu (x)=0`$ as a consequence of parity symmetry!
## 4 Lagrangians and conserved currents in RWFM
To find the conserved currents for given representation one needs the Lagrangian. Since parity symmetry already incorporates Dirac’s equation or Maxwell equations in the theory, it seems reasonable not to “force” them on the system as Euler-Lagrange equations. So one would need some other guiding light for finding the proper Lagrangian.
Let $`\phi (x)`$ belong to some representation of Lorentz group. In QFT it represents a quantum field while in RWFM it represents a wave function. All known Lagrangian densities can be written in the form
$$=\phi _A^{}(x)P_{AB}𝒪_{BC}\phi _C(x)$$
(184)
where capital Latin indices represents spin and all other “internal” indices and operators $`𝒪`$ generally have some constants, some derivative operators, some spin matrices, and for theories with internal symmetries also some matrices in internal symmetry spaces. Action can be then written as<sup>4</sup><sup>4</sup>4 From now on it will be understood that $`\phi |`$ and $`\overline{\phi }`$ mean $`(|\phi )^{}P`$ and $`\phi ^{}P`$ for any representation
$$I=d^4x=\underset{0}{\overset{\mathrm{}}{}}\phi \left|𝒪\right|\phi 𝑑t$$
(185)
Symmetries of wave functions will give us conserved currents which after integration over whole space $`𝐑^3`$ give us conserved quantities. Those conserved quantities will again have a form
$$Q^{\mu ,\mathrm{}}=j^{0,\mu ,\mathrm{}}d^3x=\overline{\phi }_A(x)𝒥_{AB}^{0,\mu ,\mathrm{}}\phi _B(x)d^3x=\phi \left|𝒥^{0,\mu ,\mathrm{}}\right|\phi $$
(186)
which in RWFM suggests the interpretation of quantities $`Q^{\mu ,\mathrm{}}`$ as (conserved) expectation values of some operators $`J^{0,\mu ,\mathrm{}}`$. The question arises which operator $`𝒪`$ should we choose for a particular representation of Lorentz group that will give us “good” conserved currents?
There are two fundamental conserved “currents” which every representation must reproduce properly: energy-momentum tensor and angular momentum density tensor. Their conserved “charges” will give matrix elements of generators of generators, energy, momentum, spin, etc. If there are no interactions, total energy should be additive, i.e. a sum of energies of all orthogonal modes, for all momentum, spin and parity eigenstates. General relativity puts even stronger restriction: energy-momentum tensor should be symmetric. Since symmetric energy-momentum tensor also gives us the proper total energy, the symmetry requirement will be adopted as another fundamental requirement of RWFM. This puts enough restrictions on operator $`𝒪`$ to determine it completely. Let’s take a look at the translational invariance requirement:
$$\phi (x)\phi (x^{})=\phi (x+a)(1+a_\alpha ^\alpha )\phi (x)\text{ for infinitesimal }a.$$
(187)
For Lagrangians in the form (184) this means the change in action will be
$$\delta I=\left(\delta \overline{\phi }_A𝒪_{AB}\phi _B+\overline{\phi }_A𝒪_{AB}\delta \phi _B\right).$$
(188)
Substituting the constant infinitesimal parameter $`a_\alpha `$ with function $`a_\alpha (x)`$ (, sec. 7.3) will give us the change in action proportional to $`_\mu a_\alpha `$ which is just the energy momentum tensor
$$\delta I=\theta ^{\mu \alpha }_\mu a_\alpha d^4x$$
(189)
with
$$\theta ^{\mu \nu }=\frac{}{^\mu \phi _A}^\nu \phi _A+^\nu \overline{\phi }_A\frac{}{^\mu \overline{\phi }_A}.$$
(190)
Effectively, what we are doing is substituting $`^\mu \phi _A`$ with $`^\nu \phi _A`$, for every spin component. We already know that one of the indices comes from derivative operator acting on field (or it’s conjugate); if energy-momentum tensor is to be symmetric, the other Lorentz index must also come from derivative operator, or in another words, indices of derivative operators must be coupled together. Contracting derivative operator with lets say Dirac’s gamma matrix wouldn’t give us symmetric energy momentum tensor. This condition narrows the form of Lagrangian to
=φ¯A
OABμ
φBμ+φ¯AOAB′′φBsubscript¯𝜑𝐴
superscriptsuperscriptsubscript𝑂𝐴𝐵𝜇
subscriptsubscript𝜑𝐵𝜇subscript¯𝜑𝐴superscriptsubscript𝑂𝐴𝐵′′subscript𝜑𝐵\mathcal{L}=\bar{\varphi}_{A}\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\mu}O_{AB}^{\prime}\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}_{\mu}\varphi_{B}+\bar{\varphi}_{A}O_{AB}^{\prime\prime}\varphi_{B} (191)
where operators $`O^{}`$ and $`O^{\prime \prime }`$ don’t depend on any derivative operators.
Invariance to Lorentz transformations will put additional restriction on Lagrangian. For infinitesimal transformation
$$\phi _A(x)\phi _A^{}(x^{})=S_{AB}(\omega )\phi _B(x)\left\{1\frac{i}{2}\omega _{\mu \nu }\underset{L^{\mu \nu }}{\underset{}{\left[J_{AB}^{\mu \nu }+(x^\mu i^\nu x^\nu i^\mu )\delta _{AB}\right]}}\right\}\phi _B(x)$$
(192)
we again get the change in action to be proportional to the angular momentum density tensor $`J^{\alpha ,\mu \nu }`$
$$\delta I=J^{\alpha ,\mu \nu }\frac{_\alpha \omega _{\mu \nu }}{2}d^4x$$
(193)
Looking at the expression for the
$`J^{\alpha ,\mu \nu }`$ $`=`$ $`{\displaystyle \frac{}{^\alpha \phi _A}}L_{AB}^{\mu \nu }\phi _B+\overline{\phi }_A\overline{L}_{AB}^{\mu \nu }{\displaystyle \frac{}{^\alpha \overline{\phi }_B}}`$ (194)
$`=`$ φ¯A
OABαLBCμνφC+φ¯AL¯ABμνOBC
φCαsubscript¯𝜑𝐴
superscriptsuperscriptsubscript𝑂𝐴𝐵𝛼subscriptsuperscript𝐿𝜇𝜈𝐵𝐶subscript𝜑𝐶subscript¯𝜑𝐴subscriptsuperscript¯𝐿𝜇𝜈𝐴𝐵superscriptsubscript𝑂𝐵𝐶
superscriptsubscript𝜑𝐶𝛼\displaystyle\bar{\varphi}_{A}\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\alpha}O_{AB}^{\prime}L^{\mu\nu}_{BC}\varphi_{C}+\bar{\varphi}_{A}\bar{L}^{\mu\nu}_{AB}O_{BC}^{\prime}\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{\alpha}\varphi_{C}
where $`\overline{L}_{AB}^{\mu \nu }=P_{AC}L_{CD}^{\mu \nu }P_{DB}`$. If the integral of zeroth component must give us expectation values of rotation and boost generators, operator $`O^{}`$ must be a constant.
To get the Lorentz invariant Lagrangian, the last factor in (191) must satisfy
$$e^{+\frac{i}{2}\omega _{\mu \nu }J^{\mu \nu }}O^{\prime \prime }e^{\frac{i}{2}\omega _{\mu \nu }J^{\mu \nu }}=O^{\prime \prime }$$
(195)
for all $`\omega _{\mu \nu }`$, or in another words it must be a scalar. Therefore it will be proportional to unit matrix in spin space, so it also has to be just a number. What we’re left with is
=φ¯A(c1
+c2)φA.subscript¯𝜑𝐴subscript𝑐1
subscript𝑐2subscript𝜑𝐴\mathcal{L}=\bar{\varphi}_{A}\left(c_{1}\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}\hbox to0.0pt{}\cdot\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}\hbox to0.0pt{}+c_{2}\right)\varphi_{A}\;. (196)
Values of $`c_1`$ and $`c_2`$ are finally fixed by the requirement of onshellness $`p^2=m^2`$, or in another words, Euler-Lagrange equations should give us Klein-Gordon equation for every component of the wave function. This finally yields the Lagrangian
=φ¯A(i
i
m2)φAsubscript¯𝜑𝐴𝑖
𝑖
superscript𝑚2subscript𝜑𝐴\mathcal{L}=\bar{\varphi}_{A}\left(i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}\hbox to0.0pt{}\cdot i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}\hbox to0.0pt{}-m^{2}\right)\varphi_{A} (197)
which is almost identical to the Klein-Gordon Lagrangian for (complex) scalar field
KG=ϕ(i
i
m2)ϕsubscript𝐾𝐺superscriptitalic-ϕ𝑖
𝑖
superscript𝑚2italic-ϕ\mathcal{L}_{KG}=\phi^{*}\left(i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}\hbox to0.0pt{}\cdot i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}\hbox to0.0pt{}-m^{2}\right)\phi (198)
with the field $`\varphi ^{}`$ replaced with the proper field $`\overline{\phi }_A`$ to make the Lagrangian relativistically invariant. Note that this derivation doesn’t depend on the representation of Lorentz group which gives us unified description of all representations.
### 4.1 Energy-momentum tensor
Let’s now derive all conserved currents for the spinor representation explicitly. Starting from Klein-Gordon-like Lagrangian (197)
=ψ¯(i
i
m2)ψ¯𝜓𝑖
𝑖
superscript𝑚2𝜓\mathcal{L}=\bar{\psi}\left(i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}\hbox to0.0pt{}\cdot i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}\hbox to0.0pt{}-m^{2}\right)\psi (199)
and requiring the translational invariance in space-time we get the (obviously symmetric) energy-momentum tensor
θμν(x)=ψ¯(x)(i
iν
+μi
iμ
)νψ\theta^{\mu\nu}(x)=\bar{\psi}(x)\left(i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\nu}\>i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{\mu}+i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\mu}\>i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{\nu}\right)\psi (200)
leading to the energy-momentum four-vector
$`P^\mu `$ $`=`$ d3xψ|x(i
i0
+μi
iμ
)0x|ψ\displaystyle\int\!\!\,d^{3}x\;\langle\psi|\vec{x}\rangle\left(i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{0}\>i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{\mu}+i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\mu}\>i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{0}\right)\langle\vec{x}|\psi\rangle\; (201)
$`=`$ d3xd3pd3qs,r,𝒫,𝒫ψ|ψp,s,𝒫ψp,s,𝒫|x(i
i0
+μi
iμ
)0x|ψq,r,𝒫ψq,r,𝒫|ψ\displaystyle\int\!\!\,d^{3}x\>\,d^{3}p\>\,d^{3}q\sum\limits_{s,r,\mathcal{P},\mathcal{P}^{\prime}}\langle\psi|\psi_{p,s,\mathcal{P}}\rangle\langle\psi_{p,s,\mathcal{P}}|\vec{x}\rangle\left(i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{0}\>i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{\mu}+i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\mu}\>i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{0}\right)\langle\vec{x}|\psi_{q,r,\mathcal{P}^{\prime}}\rangle\langle\psi_{q,r,\mathcal{P}^{\prime}}|\psi\rangle
$`=`$ $`{\displaystyle d^3pd^3q\underset{s,r,𝒫,𝒫^{}}{}\psi |\psi _{p,s,𝒫}\underset{\delta ^3(pq)\delta _{sr}\delta _{𝒫𝒫^{}},2E_pp^\mu 𝒫}{\underset{}{d^3x\psi _{p,s,𝒫}|\stackrel{}{x}\left(p^0q^\mu +p^\mu q^0\right)\stackrel{}{x}|\psi _{q,r,𝒫^{}}}}\psi _{q,r,𝒫^{}}|\psi }`$
$`=`$ $`{\displaystyle d^3p\underset{s,𝒫}{}\psi |\psi _{p,s,𝒫}\left\{2p^0p^\mu 𝒫\right\}\psi _{p,s,𝒫}|\psi }={\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_p}p^\mu \underset{s}{}\left(b_{ps}^{}b_{ps}d_{ps}^{}d_{ps}\right)}`$
There is a negative sign here which comes from negative norm of negative parity states $`v_{ps}`$ states and has nothing to do with the energies of the solutions which are positive for all solutions.
States $`|\psi _{p,s,𝒫}`$ form a basis of Hilbert space which enables us to express energy-momentum four-vector operator through it’s matrix elements
$$\widehat{P}^\mu =|\psi _{p,s,𝒫}\psi _{p,s,𝒫}\left|\widehat{P}^\mu \right|\psi _{p,s,𝒫}\psi _{p,s,𝒫}|$$
(202)
Since the state $`|\psi `$ is arbitrary, we can read the energy-momentum operator from it’s matrix element (201)
$`\widehat{P}^\mu `$ $`=`$ $`{\displaystyle d^3p\underset{s,𝒫}{}|\psi _{p,s,𝒫}\left\{2p^0p^\mu 𝒫\right\}\psi _{p,s,𝒫}|}`$ (203)
$`=`$ $`{\displaystyle d^3p\underset{s}{}\left\{2p^0p^\mu \right\}\left(|\psi _{p,s,+}\psi _{p,s,+}||\psi _{p,s,}\psi _{p,s,}|\right)}`$ (204)
Note that both equations (204) and (201) have negative sign for both energy operator and expectation value. Never the less, energy-momentum operator always gives *positive* result for energy
$`\widehat{P}^\mu |\psi _{q,r,+}`$ $`=`$ $`{\displaystyle d^3p\underset{s}{}\left\{2p^0p^\mu \right\}\left(|\psi _{p,s,+}\underset{\delta ^3(\stackrel{}{p}\stackrel{}{q})\delta _{sr}}{\underset{}{\psi _{p,s,+}|\psi _{q,r,+}}}|\psi _{p,s,}\underset{0}{\underset{}{\psi _{p,s,}|\psi _{q,r,+}}}\right)}`$ (205)
$`=`$ $`\left\{2q^0q^\mu \right\}|\psi _{q,r,+}`$
$`\widehat{P}^\mu |\psi _{q,r,}`$ $`=`$ $`{\displaystyle d^3p\underset{s}{}\left\{2p^0p^\mu \right\}\left(|\psi _{p,s,+}\underset{0}{\underset{}{\psi _{p,s,+}|\psi _{q,r,}}}|\psi _{p,s,}\underset{\delta ^3(\stackrel{}{p}\stackrel{}{q})\delta _{sr}}{\underset{}{\psi _{p,s,}|\psi _{q,r,}}}\right)}`$ (206)
$`=`$ $`\left\{2q^0q^\mu \right\}|\psi _{q,r,}.`$
Expectation values have negative part since *norms* of those states are negative, not *energies*. Dividing with $`\psi |\psi `$ we get the quantity
$$E=\frac{\psi \left|\widehat{P}^0\right|\psi }{\psi |\psi }>0$$
(207)
which is by definition positive for all states in Hilbert space. Comparing equation (201)
$$P^\mu =\frac{d^3p}{2E_p}p^\mu \underset{s}{}\left(b_{ps}^{}b_{ps}d_{ps}^{}d_{ps}\right)$$
(208)
with the expression (44) from the Dirac Lagrangian
$$P^\mu =\theta ^{0\mu }d^3x=\frac{d^3p}{2E_p}p^\mu \left(b_{ps}^{}b_{ps}d_{ps}d_{ps}^{}\right).$$
(209)
we can see that they are the same aside from the ordering of the $`dd^{}`$ term. Primary reason for introducing anticommutators was to make the expectation value (209) positive definite; there’s no reason to do that here since the negative sign of EV doesn’t imply negative energy.
### 4.2 Angular momentum density tensor
Requirement for infinitesimal rotational and boost invariance for spinor of Lorentz group gives us
ψ(x){1i2ωμν[σμν2+(xμi
νxνi
)μ]Lμν}ψ(x)\psi(x)\to\left\{1-\frac{i}{2}\omega_{\mu\nu}\underbrace{\left[\frac{\sigma^{\mu\nu}}{2}+(x^{\mu}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{\nu}-x^{\nu}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{\mu})\right]}_{L^{\mu\nu}}\right\}\psi(x) (210)
and
$`\overline{\psi }(x)`$ $``$ ψ¯(x){1+i2ωμν[γ0σμν2γ0+(xμi
νxνi
)μ]}\displaystyle\bar{\psi}(x)\left\{1+\frac{i}{2}\omega_{\mu\nu}\left[\gamma^{0}\frac{\sigma^{\mu\nu\,\dagger}}{2}\gamma^{0}+(x^{\mu}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\nu}-x^{\nu}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\mu})\right]\right\} (211)
$``$ ψ¯(x){1+i2ωμν[σμν2+(xμi
νxνi
)μ]L¯μν}\displaystyle\bar{\psi}(x)\left\{1+\frac{i}{2}\omega_{\mu\nu}\underbrace{\left[\frac{\sigma^{\mu\nu}}{2}+(x^{\mu}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\nu}-x^{\nu}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\mu})\right]}_{\bar{L}^{\mu\nu}}\right\} (212)
where matrix $`\gamma ^0`$ is parity matrix $`P`$ in spinor space. This invariance gives us conserved currents
$`J^{\alpha ,\mu \nu }`$ $`=`$ ψ¯(x)i
Lμναψ+ψ¯(x)L¯μνi
ψα¯𝜓𝑥𝑖
superscriptsuperscript𝐿𝜇𝜈𝛼𝜓¯𝜓𝑥superscript¯𝐿𝜇𝜈𝑖
superscript𝜓𝛼\displaystyle\bar{\psi}(x)i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\alpha}L^{\mu\nu}\psi+\bar{\psi}(x)\bar{L}^{\mu\nu}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{\alpha}\psi (213)
$`=`$ xμθανxνθαν+ψ¯(i
σμν2α+σμν2i
)αψ\displaystyle x^{\mu}\theta^{\alpha\nu}-x^{\nu}\theta^{\alpha\nu}+\bar{\psi}\left(i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\alpha}\frac{\sigma^{\mu\nu}}{2}+\frac{\sigma^{\mu\nu}}{2}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{\alpha}\right)\psi (214)
which lead to conserved quantities
$$J^{\mu \nu }=J^{0,\mu \nu }d^3x=J_{coord}^{\mu \nu }+J_{spin}^{\mu \nu }$$
(215)
where coordinate or orbital part is defined to be the part proportional to unit matrix in spin space
Jcoordμν=ψ¯{(xμi
νxνi
)μi
+0i
(xμi
νxνi
)ν0}ψd3xJ_{coord}^{\mu\nu}=\int\bar{\psi}\left\{\left(x^{\mu}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\nu}-x^{\nu}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{\mu}\right)i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{0}+i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{0}\left(x^{\mu}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{\nu}-x^{\nu}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{\nu}\right)\right\}\psi\,d^{3}x (216)
and the spin part the rest
Jspinμν=ψ¯{i
σμν20+σμν2i
}0ψd3x.J_{spin}^{\mu\nu}=\int\bar{\psi}\left\{i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{0}\frac{\sigma^{\mu\nu}}{2}+\frac{\sigma^{\mu\nu}}{2}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{0}\right\}\psi\,d^{3}x\;. (217)
Both parts can be decomposed in momentum eigenfunctions $`|u_{ps}`$ and $`|v_{ps}`$
$`J^{\mu \nu }`$ $`=`$ $`{\displaystyle \left[b_{ps}^{}\left(p^\mu i_p^\nu p^\nu i_p^\mu \right)b_{ps}d_{ps}^{}\left(p^\mu i_p^\nu p^\nu i_p^\mu \right)d_{ps}\right]\frac{d^3p}{2E_p}}`$ (218)
Spin part of generators in the momentum states is
$`J_{spin}^{\mu \nu }`$ $`=`$ ψ¯{i
σμν20+σμν2i
}0ψd3x\displaystyle\int\bar{\psi}\left\{i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\leftarrow$}}\hbox{$\partial$}}{}^{0}\frac{\sigma^{\mu\nu}}{2}+\frac{\sigma^{\mu\nu}}{2}i\kern 1.66672pt\vbox{\hbox{\kern-0.50003pt\lower 3.99994pt\hbox{\scriptsize$\rightarrow$}}\hbox{$\partial$}}{}^{0}\right\}\psi\,d^{3}x (219)
$`=`$ $`{\displaystyle }{\displaystyle \frac{d^3p}{2E_p}}{\displaystyle \underset{s,r}{}}(b_{ps}^{}b_{pr}\overline{u}_{ps}{\displaystyle \frac{\sigma ^{\mu \nu }}{2}}u_{pr}+d_{ps}^{}d_{pr}\overline{v}_{ps}{\displaystyle \frac{\sigma ^{\mu \nu }}{2}}v_{pr}`$
$`+b_{ps}^{}d_{pr}\overline{u}_{ps}{\displaystyle \frac{\sigma ^{\mu \nu }}{2}}v_{pr}+d_{ps}^{}b_{pr}\overline{v}_{ps}{\displaystyle \frac{\sigma ^{\mu \nu }}{2}}u_{pr}).`$
Mixed elements are not going to vanish; for boost part it is what we expect, but for rotational part it means trouble. This makes it hard to interpret states $`u_{ps}`$ as particles and states $`v_{ps}`$ as anti-particles; it would be strange to have rotations mix particles and anti-particles (electrons and positrons for example). States $`u_{ps}`$ and $`v_{ps}`$ are (by definition) obtained from rest frame states by applying the boost operator
$$|u_{ps}=e^{i\stackrel{}{\omega }\stackrel{}{K}}|u_{0,s}|v_{ps}=e^{i\stackrel{}{\omega }\stackrel{}{K}}|v_{0,s}.$$
(220)
Since boosts and rotations don’t commute, so if we started with spin eigenstates $`|u_{0,s}`$ and $`|v_{0,s}`$, final states $`|u_{ps}`$ and $`|v_{ps}`$ won’t be spin eigenstates which is another reason why they shouldn’t be used to describe particles. However, they *do* form a basis for given momentum $`\stackrel{}{p}`$ so as long as we work with unpolarized states it doesn’t really matter which basis we use. The question of spine eigenstates is finally addressed in section 5.
It’s again instructive to compare these results with the generators obtained from Dirac Lagrangian. Coordinate part
$`J_{\mathrm{𝑐𝑜𝑜𝑟𝑑}}^{\mu \nu }`$ $`=`$ 12ψ(xμi
νxνi
)μψd3x\displaystyle\int\frac{1}{2}\psi^{\dagger}\left(x^{\mu}i\kern 1.66672pt\vbox{\hbox{\kern-1.66672pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}{}^{\nu}-x^{\nu}i\kern 1.66672pt\vbox{\hbox{\kern-1.66672pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}{}^{\mu}\right)\psi\,d^{3}x (221)
$`=`$ $`{\displaystyle \frac{E_p}{m}\left(b_{ps}^{}\left(p^\mu i_p^\nu p^\nu i_p^\mu \right)b_{ps}d_{ps}\left(p^\mu i_p^\nu p^\nu i_p^\mu \right)d_{ps}^{}\right)\frac{d^3p}{2E_p}}.`$
is again (almost) the same, differing by a factor $`E/m`$ and the ordering of $`dd^{}`$ terms. Spin part of rotation generators is again quite similar to the Dirac case
$`J_k`$ $``$ $`ϵ_{ijk}J^{ij}=ϵ_{ijk}{\displaystyle \psi ^{}\frac{\sigma ^{ij}}{2}\psi d^3x}`$ (222)
$`=`$ $`{\displaystyle }{\displaystyle \frac{d^3p}{2E_p}}{\displaystyle \frac{1}{E_p}}{\displaystyle \underset{s,r}{}}(u_{ps}^{}{\displaystyle \frac{\sigma ^k}{2}}u_{pr}b_{ps}^{}b_{pr}+v_{ps}^{}{\displaystyle \frac{\sigma ^k}{2}}v_{pr}d_{ps}d_{pr}^{}`$
$`+e^{2iE_pt}u_{ps}^{}{\displaystyle \frac{\sigma ^k}{2}}v_{\stackrel{~}{p}r}b_{ps}^{}d_{\stackrel{~}{p}r}^{}+e^{2iE_pt}v_{ps}^{}{\displaystyle \frac{\sigma ^k}{2}}u_{\stackrel{~}{p}r}d_{ps}b_{\stackrel{~}{p}r})`$
but it doesn’t have time dependent *zitterbewegung* exponentials. If one were to perform a second quantization it would annihilate the vacuum as it should and it wouldn’t mix states with different number of particles under rotation. Furthermore, boost part of generators doesn’t vanish as in Dirac case and it has the proper functional dependence similar to rotation generators.
### 4.3 “Probability” current
Finally, since wave functions for all representations are complex numbers, and since all observables are real numbers, changing the phase of wave function should leave all observables unchanged. In another words, symmetry transformation
$$\psi (x)e^{i\alpha }\psi (x)=\left(1i\alpha \right)\psi (x)$$
(223)
creates the current
$`j^\mu (x)`$ $`=`$ 2ψ¯(x)i
ψμ(x)2¯𝜓𝑥𝑖
superscript𝜓𝜇𝑥\displaystyle 2\bar{\psi}(x)i\kern 1.66672pt\vbox{\hbox{\kern-1.00006pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}{}^{\mu}\psi(x) (224)
$`=`$ $`{\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.25em}2}E_p}\frac{d^3q}{(2\pi )^{3/2}\mathrm{\hspace{0.25em}2}E_q}\underset{s,𝒫}{}\underset{r,𝒫^{}}{}e^{i(pq)x}}`$
$`\times \left(b_{ps}^{}\overline{u}_{ps}+d_{ps}^{}\overline{v}_{ps}\right)(p+q)^\mu \left(b_{qr}u_{qr}+d_{qr}v_{qr}\right).`$
Conserved quantity will be the integral of zeroth component
$`Q`$ $`=`$ j0(x)d3x=φ|2i
|0φ=d3p2Eps(bpsbpsdpsdps)\displaystyle\int j^{0}(x)\,d^{3}x=\left\langle{\varphi}\left|{2i\kern 1.66672pt\vbox{\hbox{\kern-1.00006pt\lower 3.99994pt\hbox{$\leftrightarrow$}}\hbox{$\partial$}}{}^{0}}\right|{\varphi}\right\rangle\;=\int\frac{\,d^{3}p}{2E_{p}}\sum\limits_{s}\left(b^{*}_{ps}b_{ps}-d^{*}_{ps}d_{ps}\right) (225)
Note that there is no zitterbewegung terms due to the orthogonality of states. In non-relativistic quantum mechanics zero-th component of the four-current is positive definite so one can interpret it as the probability density and the current as probability flux. Since the current is no longer neither positive definite nor bound, one must think twice before calling it the probability current.
One can again compare equation (225) with the equation for charge from Dirac’s equation (36) *after* the second integration and normal ordering
$$:Q:=:j^0(x):d^3x=\frac{d^3p}{2E_p}(b_{ps}^{}b_{ps}d_{ps}^{}d_{ps}).$$
(226)
We get the same expression for charge, but the origin of the minus sign is different; here it comes from the negative norm, and there it comes from anticommutators. Space part of the total current
$`J^i`$ $`=`$ $`{\displaystyle \frac{d^3p}{2E_p}\frac{1}{E_p}\left(p^i\underset{s}{}\left[b_{ps}^{}b_{ps}d_{ps}^{}d_{ps}\right]\right)}`$ (227)
can be compared with the equivalent current from Dirac Lagrangian (37)
$`:J^i:`$ $`=`$ $`{\displaystyle }:j^i(x):d^3x={\displaystyle }{\displaystyle \frac{d^3p}{2E_p}}{\displaystyle \frac{1}{E_p}}(p^i{\displaystyle \underset{s}{}}:[b_{ps}^{}b_{ps}+d_{ps}d_{ps}^{}]:`$ (228)
$`+{\displaystyle \underset{s,r}{}}[{\displaystyle \frac{ie^{2iE_pt}}{2m}}\overline{u}_{ps}\sigma ^{i0}v_{\stackrel{~}{p}r}b_{ps}^{}d_{\stackrel{~}{p}r}^{}{\displaystyle \frac{ie^{2iE_pt}}{2m}}\overline{v}_{ps}\sigma ^{i0}u_{\stackrel{~}{p}r}d_{ps}b_{\stackrel{~}{p}r}])`$
Note that aside from zitterbewegung terms, current (227) is the same the normal ordered current from Dirac Lagrangian.
## 5 Matrix elements and particle interpretation
Let’s take another look at mixed matrix elements for rotation generators. It is most convenient to use Dirac’s representation; in this representation rotation generators and spinors are given by
$$J_i=\frac{1}{2}\left(\begin{array}{cc}\sigma _i& 0\\ 0& \sigma _i\end{array}\right)$$
(229)
$$u_{ps}=\frac{1}{\sqrt{2m(E+m)}}\left(\begin{array}{c}(E+m)\chi _s\\ \stackrel{}{p}\stackrel{}{\sigma }\chi _s\end{array}\right),v_{ps}=\frac{1}{\sqrt{2m(E+m)}}\left(\begin{array}{c}\stackrel{}{p}\stackrel{}{\sigma }\chi _s\\ (E+m)\chi _s\end{array}\right).$$
(230)
Mixed matrix elements of rotation generators are then
$`\overline{u}_{ps}{\displaystyle \frac{\sigma _i}{2}}v_{pr}`$ $`=`$ $`{\displaystyle \frac{1}{2m(E+m)}}\left(\begin{array}{cc}(E+m)\chi _s& \chi _s\stackrel{}{p}\stackrel{}{\sigma }\end{array}\right){\displaystyle \frac{1}{2}}\left(\begin{array}{cc}\sigma _i& 0\\ 0& \sigma _i\end{array}\right)\left(\begin{array}{c}\stackrel{}{p}\stackrel{}{\sigma }\chi _r\\ (E+m)\chi _r\end{array}\right)`$ (236)
$`=`$ $`{\displaystyle \frac{1}{2m}}\left(\chi _s\sigma _i\stackrel{}{p}\stackrel{}{\sigma }\chi _r\chi _s\stackrel{}{p}\stackrel{}{\sigma }\sigma _i\chi _r\right)={\displaystyle \frac{2}{m}}\chi _s(\stackrel{}{p}\times \stackrel{}{\sigma })_i\chi _r`$ (237)
$`\overline{v}_{ps}{\displaystyle \frac{\sigma _i}{2}}u_{pr}`$ $`=`$ $`{\displaystyle \frac{2}{m}}\chi _s(\stackrel{}{p}\times \stackrel{}{\sigma })_i\chi _r`$ (238)
which will always have at least one non-vanishing component. This would imply that for example applying rotation to pure electron state would give us mixture of electrons and positrons
$$\text{positron}\left|e^{i\stackrel{}{\omega }\stackrel{}{J}}\right|\text{electron}0$$
(239)
In the QFT this problem didn’t exist since the second quantization would just “sweep it under the carpet” by multiplying the non-vanishing elements with either two creation or two destruction operators (which of course produced another problems). This cannot be done here so the only thing we can do is revising our interpretation of states $`u_{ps}`$ and $`v_{ps}`$. As it was said earlier, since they are defined as states obtained by boosting the rest frame spin eigenstates, obviously they are not spin eigenstates at all. In non-relativistic quantum mechanics spin and momentum operators commute so proper spin eigenstates (in parity representation) are constructed as
$$w_{ps}^+=e^{ipx}\left(\begin{array}{c}\chi _s\\ 0\end{array}\right),w_{ps}^{}=e^{ipx}\left(\begin{array}{c}0\\ \chi _s\end{array}\right).$$
(240)
We can think about them as states obtained by applying *only coordinate* part of the boost operator (which commutes with spin generators) to rest-frame spin eigenstates. Particle wave function can then be decomposed as
$`\psi (x)`$ $``$ $`\stackrel{}{x}|\psi =\stackrel{}{x}\left|\left({\displaystyle d^3p\underset{s,𝒫}{}\frac{|w_{p,s,𝒫}w_{p,s,𝒫}|}{w_{p,s,𝒫}|w_{p,s,𝒫}}}\right)\right|\psi `$ (241)
$`=`$ $`{\displaystyle d^3p\underset{s}{}\left(\stackrel{}{x}|w_{p,s}^+w_{ps}^+|\psi \stackrel{}{x}|w_{p,s}^{}w_{ps}^{}|\psi \right)}`$
$`=`$ $`{\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_p}\underset{s}{}e^{ipx}\left(a_{ps}^+w_{ps}^++a_{ps}^{}w_{ps}^{}\right)}`$
which is the same as the decomposition in $`u`$, $`v`$ states with the substitution $`ba^+`$, $`da^{}`$, $`uw^+`$ and $`vw^{}`$. Their mixed matrix elements of rotation generators all vanish by their construction
$`w_{ps}^+\left|\stackrel{}{J}\right|w_{pr}^{}`$ $``$ $`\left(\begin{array}{cc}\chi _s& 0\end{array}\right)\left(\begin{array}{cc}\stackrel{}{\sigma }& 0\\ 0& \stackrel{}{\sigma }\end{array}\right)\left(\begin{array}{c}0\\ \chi _r\end{array}\right)=0`$ (247)
$`w_{ps}^{}\left|\stackrel{}{J}\right|w_{pr}^+`$ $``$ $`\left(\begin{array}{cc}0& \chi _s\end{array}\right)\left(\begin{array}{cc}\stackrel{}{\sigma }& 0\\ 0& \stackrel{}{\sigma }\end{array}\right)\left(\begin{array}{c}\chi _r\\ 0\end{array}\right)=0.`$ (253)
which implies
$$w_{ps}^\pm \left|e^{i\stackrel{}{\omega }\stackrel{}{J}}\right|w_{pr}^{}=0$$
(254)
as it should. Since boost generators aren’t diagonal in this representation, they *will* mix states of different parity
$`w_{ps}^+\left|\stackrel{}{K}\right|w_{pr}^{}`$ $``$ $`\left(\begin{array}{cc}\chi _s& 0\end{array}\right)\left(\begin{array}{cc}0& \stackrel{}{\sigma }\\ \stackrel{}{\sigma }& 0\end{array}\right)\left(\begin{array}{c}0\\ \chi _r\end{array}\right)0`$ (260)
$`w_{ps}^{}\left|\stackrel{}{K}\right|w_{pr}^+`$ $``$ $`\left(\begin{array}{cc}0& \chi _s\end{array}\right)\left(\begin{array}{cc}0& \stackrel{}{\sigma }\\ \stackrel{}{\sigma }& 0\end{array}\right)\left(\begin{array}{c}\chi _r\\ 0\end{array}\right)0.`$ (266)
This implies that Lorentz boosts mix positive and negative parity and spin eigenstates
$$w_{ps}^\pm \left|e^{i\stackrel{}{\omega }\stackrel{}{K}}\right|w_{pr}^{}0$$
(267)
This phenomenon is already known from $`(1,0)(0,1)`$ representation where Lorentz boosts mix electric and magnetic fields. In fact, it’s a fundamental property of every $`(j,0)(0,j)`$ representation which tells us that axial and polar states will exist and mix for every spin and can’t be separated.
When negative parity states are interpreted as wave functions of antiparticles, it also offers a possible and quite intriguing explanation why anti-particles have opposite quantum numbers from particles: all states (both polar and axial) will have the same eigenvalues for some “internal” symmetry generators like for example electric charge $`\widehat{Q}_i=\widehat{j}^0(x)d^3x`$
$$\widehat{Q}_i|\psi ^\pm =q_i|\psi ^\pm $$
(268)
but where that operator is coupled to whatever field through it’s (conserved) current
$$_I=\psi ^\pm \left|j^{\mu ,\mathrm{}}\right|\psi ^\pm 𝒪_\mu (\text{some fields})q_i\left(\underset{\stackrel{}{p},s}{}\left|a_{ps}^+\right|^2p^\mu 𝒪_\mu \left|a_{ps}^{}\right|^2p^\mu 𝒪_\mu \right)$$
(269)
*expectation values* of states with opposite parities will have relative minus sign which change the sign of interaction Hamiltonian which we used to interpret as opposite quantum numbers. In the classical limit this relative minus sign leads to attractive or repulsive force. This is a property of states themselves, not the operators $`\widehat{Q}_i`$. Since every interaction of fermions in standard model is of the form (269), what we observe is the difference in sign of the energy and interpret it as opposite quantum number; in this model it comes from opposite parity instead.
### 5.1 Massless particles
Finally, a few words about massless (or equivalently ultra-relativistic) limit. In section 3 it was shown that in massless limit whole wave function satisfies a system of coupled differential equations. In the case of spinor representation these equations are
$`i{\displaystyle \frac{}{t}}\left(\begin{array}{c}\phi (x)\\ \chi (x)\end{array}\right)=\left(\begin{array}{cc}\stackrel{}{\sigma }\left(i\right)& 0\\ 0& \stackrel{}{\sigma }\left(i\right)\end{array}\right)\left(\begin{array}{c}\phi (x)\\ \chi (x)\end{array}\right)`$ (276)
with solutions
$`\psi _{R,L}=\left(\begin{array}{c}\phi _{R,L}(x)\\ \pm \phi _{R,L}(x)\end{array}\right)`$ (279)
corresponding to separate $`(1/2,0)`$ and $`(0,1/2)`$ transformations. Now, regardless of which one we choose, equations (276) become
$`i{\displaystyle \frac{}{t}}\phi _{R,L}(x)=\stackrel{}{\sigma }\left(i\right)\phi _{R,L}(x).`$ (280)
With the interpretation of the previous subsection, we can decompose wave functions $`\phi _{R,L}(x)`$ in momentum and spin eigenstates in chiral basis
$`\phi _{R,L}(x)`$ $`=`$ $`{\displaystyle \frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.17em}2}E_p}\underset{s}{}e^{ipx}a_{ps}^{R,L}w_{ps}^{R,L}}`$ (281)
which transforms the equation (276) to
$`w_{ps}^{R,L}(x)={\displaystyle \frac{\stackrel{}{p}\stackrel{}{S}}{E}}w_{ps}^{R,L}(x)=𝚺w_{ps}^{R,L}(x).`$ (283)
where the product of translation and rotation generators $`𝚺\stackrel{}{p}\stackrel{}{S}/\left|\stackrel{}{p}\right|`$ is nothing but helicity operator. Since all basis states are positive helicity eigenstates, the total wave function as a superposition of these states will also be a positive helicity eigenstate. However, just as in the case of the zeroth component of conserved currents, expectation values for particle and antiparticle states will have opposite signs due to the definition of scalar product
$$\psi \left|𝚺\right|\psi \left(\underset{\stackrel{}{p},s}{}\left|a_{ps}^+\right|^2\left|a_{ps}^{}\right|^2\right)$$
(284)
In nature neutrinos have been observed to have only positive helicity while antineutrinos have only negative helicities. This is just the matter of convention since we could have chosen to assign names “particle” and “antiparticle” in the reverse order and the parity operator is defined up to a sign anyway. What is physically important is that in the massless limit particles become helicity eigenstates with opposite expectation values.
Analog thing happens in $`(1,0)(0,1)`$ representation as well; in massless limit electric and magnetic field satisfy equations (3.3). In terms of particle eigenstates this becomes
$`\left(\stackrel{}{p}\stackrel{}{S}\right)\left[\left(\stackrel{}{p}\stackrel{}{S}\right)\phi _{ps}^\pm +E\chi _{ps}^\pm \right]`$ $`=`$ $`0`$ (285)
$`\left(\stackrel{}{p}\stackrel{}{S}\right)\left[\left(\stackrel{}{p}\stackrel{}{S}\right)\chi _{ps}^\pm E\phi _{ps}^\pm \right]`$ $`=`$ $`0`$ (286)
which again has solutions $`w_{ps}^{R,L}`$ corresponding to $`\chi _{ps}^\pm =\pm \phi _{ps}^\pm `$ for separate $`(1,0)`$ and $`(0,1)`$ transformations. These solutions in terms of fields $`\stackrel{}{}`$ and $`\stackrel{}{B}`$ are $`\stackrel{}{}=\pm i\stackrel{}{B}`$, consistent with the result familiar from the classical electrodynamics $`\stackrel{}{}=i\stackrel{}{B}`$. In this framework, one can interpret both the Dirac and Maxwell equations as a statement about parity properties of particles in the massless limit.
## 6 Summary and Conclusions
Symmetry treatment of non-relativistic quantum mechanics is generalized to include non-unitary representations of Lorentz group by redefining the scalar product of states in Hilbert space to make it relativistically invariant. This inevitably leads to states with negative norm. However, with this definition of scalar product, it is shown that superposition principle and orthogonality of quantum mechanics leads to the conclusion that all energies are positive. Furthermore, it is shown that the treatment of parity in different Lorentz frames leads to Dirac equation (for spinor representation) and to (sourceless) Maxwell equations for vector fields. It is demonstrated that if one requires “proper” behavior of Noether currents corresponding to translations, rotations and boosts, form of Lagrangian is determined completely for any representation of Lorentz group. The resulting theory is just the “ordinary” single particle quantum mechanics, but now relativistically invariant. This theory doesn’t have negative energies, doesn’t show zitterbewegung-like effects, has clear transformation properties under Lorentz transformations as well as clear interpretation of physical particle states as momentum, spin and parity eigenstates. Continuity current density is proportional to the energy density and momentum density, just like in non-relativistic quantum mechanics, suggesting the possible probabilistic interpretation. One has to be careful and think twice before doing just that since norms are no longer positive definite which could potentially lead to trouble. While this might look a bit discouraging, one has to remember that states with negative norms are unavoidable in covariant formulation of QFT as well. Although that may seem to be the problem with the theory, it does also offer a nice interpretation why particles and antiparticles have opposite quantum numbers or why neutrinos always appear in nature with positive helicity and antineutrinos with negative. Due to the length of the subject, discussion of multiparticle states and the theory of interactions will be presented in a separate paper.
## Acknowledgments
The author would like to thank Krešimir Kumerički, Predrag Prester, Marko Kolanović and Saša Bistrović for many helpful discussions during the preparation of the manuscript. Many thanks to John D. Trout for many helpful explanations and links about the mathematical background of quantum field theory. All errors and misconceptions are mine exclusively.
## Appendix A General properties of Lorentz group
Special theory of relativity requires that the speed of light $`c`$ remains the same in all inertial frames. Mathematically this means the 4-distance between 2 real 4-vectors
$$s=(t_2t_1)^2(\stackrel{}{x}_2\stackrel{}{x}_1)^2$$
(287)
must remain the same in all inertial frames. In quantum physics we have to deal with complex representations as well. If we write complex 4-vector $`a^\mu `$ in matrix notation as $`|a`$, with
$$|a\left(\begin{array}{c}a_0\\ a_1\\ a_2\\ a_3\end{array}\right)$$
(288)
then the invariant quantity isn’t
$$a|a=a_0^{}a_0+a_1^{}a_1+a_2^{}a_2+a_3^{}a_3$$
(289)
but
$$a_\mu ^{}a^\mu =a_0^{}a_0a_1^{}a_1a_2^{}a_2a_3^{}a_3$$
(290)
where $`a_\mu `$ is defined to be $`a^\mu `$ multiplied by 4-tensor called metric tensor
$$g_{\mu \nu }=\text{diag}\{1,1,1,1\}.$$
(291)
Pure rotations will mix only the space components of 4-vector
$$|a^{}=S(\omega )|a,\text{or}\left(\begin{array}{c}a_0\\ a_1^{}\\ a_2^{}\\ a_3^{}\end{array}\right)=S(\omega )\left(\begin{array}{c}a_0\\ a_1\\ a_2\\ a_3\end{array}\right)$$
(292)
so they can be represented in block form as
$$S(\omega )=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& & & \\ 0& & R(\omega )& \\ 0& & & \end{array}\right)=e^{i\stackrel{}{\omega }\stackrel{}{J}}$$
(293)
with rotation generators
$$J_x=i\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right)J_y=i\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 0& 0\\ 0& 1& 0& 0\end{array}\right)J_z=i\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 0\end{array}\right)$$
(294)
Pure Lorentz boosts $`S(\omega )=e^{i\stackrel{}{\omega }\stackrel{}{K}}`$ mix space and time coordinates
$$|a^{}=S(\omega )|a,\text{or}\left(\begin{array}{c}a_0^{}\\ a_1^{}\\ a_2^{}\\ a_3^{}\end{array}\right)=S(\omega )\left(\begin{array}{c}a_0\\ a_1\\ a_2\\ a_3\end{array}\right)$$
(295)
and have generators
$$K_x=i\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right)K_y=i\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right)K_y=i\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 1& 0& 0& 0\end{array}\right).$$
(296)
They can be put together in 4D notation with $`J^{0i}K_i`$, and $`J^{ij}ϵ_{ijk}J_k`$. Commutation relations
$$[J_i,J_j]=iϵ_{ijk}J_k[K_i,K_j]=iϵ_{ijk}J_k[J_i,K_j]=iϵ_{ijk}K_k$$
(297)
in 4D notation become
$$[J^{\mu \nu },J^{\alpha \beta }]=i\left(g^{\mu \alpha }J^{\nu \beta }+g^{\nu \beta }J^{\mu \alpha }g^{\mu \beta }J^{\nu \alpha }g^{\nu \alpha }J^{\mu \beta }\right)$$
(298)
and the finite Lorentz transformations can be written as
$$S(\omega )=e^{i\omega _{\mu \nu }J^{\mu \nu }}.$$
(299)
Up to this point it wasn’t important if vectors depend on coordinates or not. Generators (294) and (296) just mix different components of vectors. If those vectors *do* depend on coordinates, then those coordinates will have to be transformed as well. For state $`|a(x)`$ we have
$$|a^{}(x^{})=S(\omega )|a(x)$$
(300)
which for infinitesimal $`\omega `$ becomes
$$|a^{}(x^{})=(1i\omega _{\mu \nu }J^{\mu \nu })|a(x^{})=(1i\omega _{\mu \nu }\underset{L^{\mu \nu }}{\underset{}{\{J^{\mu \nu }+x^\mu i^\nu x^\nu i^\mu \}}})|a(x).$$
(301)
The first $`J^{\mu \nu }`$ term will be called spin part of generators and the second coordinate part. It is trivial to see that spin and coordinate parts of generators commute and using that fact to explicitly show that the whole Lorentz transformation operator $`L^{\mu \nu }`$ satisfies the same commutation relations (298).
Equation (287) will also be invariant to the translations in space-time
$$x^\mu x^\mu +b^\mu $$
(302)
How should that affect vector $`|a`$? If it doesn’t depend on the coordinates $`x^\mu `$ it should obviously be left unchanged by the translation of coordinate system. On the other hands, if it does depend on the coordinates, it should be just the same vector in new coordinates
$$|a(x)|a^{}(x^{})=|a(x+b).$$
(303)
This obviously transforms each component of the vector $`|a(x)`$ independently. Expanding each component in Taylor series and keeping only the first term, we get the infinitesimal transformation
$$|a(x)|a^{}(x^{})=e^{ib^\mu P_\mu }|a(x)=(1+b^\mu _\mu )|a(x)$$
(304)
which give us the generator of transformations
$$P^\mu (x)=i^\mu .$$
(305)
It is important to notice that there is no spin part of this generator; it has only coordinate part.
### A.1 Parity
We can easily see that operator with the property
$$P\stackrel{}{a}=\stackrel{}{a}\text{or}P\left(\begin{array}{c}a_0\\ a_1\\ a_2\\ a_3\end{array}\right)=\left(\begin{array}{c}a_0\\ a_1\\ a_2\\ a_3\end{array}\right)$$
(306)
leave (287) invariant as well. In matrix form, it can be written as
$$P=𝒫_{\mu }^{}{}_{}{}^{\nu }=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right)$$
(307)
By direct multiplication we can show that
$$P\stackrel{}{J}P=\stackrel{}{J}P\stackrel{}{K}P=\stackrel{}{K}.$$
(308)
or in 4D notation
$$PJ^{\mu \nu }P=𝒫_{\alpha }^{}{}_{}{}^{\mu }𝒫_{\beta }^{}{}_{}{}^{\nu }J^{\alpha \beta }$$
(309)
In mathematical terms, six generators of Lorentz group can be decomposed as a direct product of two $`SU(2)`$ groups. Vectors will then belong to a direct product space of these two $`SU(2)`$ groups. If we use
$$\stackrel{}{J}=\stackrel{}{A}\mathrm{𝟏}+\mathrm{𝟏}\stackrel{}{B}\stackrel{}{K}=i\left(\stackrel{}{A}\mathrm{𝟏}\mathrm{𝟏}\stackrel{}{B}\right)$$
(310)
we can re-express (297) as
$`[J_i,J_j]`$ $`=`$ $`[A_i\mathrm{𝟏},A_j\mathrm{𝟏}]+[\mathrm{𝟏}B_i,A_j\mathrm{𝟏}]+[A_i\mathrm{𝟏},\mathrm{𝟏}B_j]+[\mathrm{𝟏}B_i,\mathrm{𝟏}B_j]`$ (311)
$`=`$ $`[A_i\mathrm{𝟏},A_j\mathrm{𝟏}]+[\mathrm{𝟏}B_i,\mathrm{𝟏}B_j]`$
$`=`$ $`[A_i,A_j]\mathrm{𝟏}+\mathrm{𝟏}[B_i,B_j]`$
$`[K_i,K_j]`$ $`=`$ $`i^2([A_i\mathrm{𝟏},A_j\mathrm{𝟏}][\mathrm{𝟏}B_i,A_j\mathrm{𝟏}][A_i\mathrm{𝟏},\mathrm{𝟏}B_j]`$ (312)
$`+[\mathrm{𝟏}B_i,\mathrm{𝟏}B_j])=i^2([A_i\mathrm{𝟏},A_j\mathrm{𝟏}]+[\mathrm{𝟏}B_i,\mathrm{𝟏}B_j])`$
$`=`$ $`\left([A_i,A_j]\mathrm{𝟏}+\mathrm{𝟏}[B_i,B_j]\right)`$
$`[J_i,K_j]`$ $`=`$ $`i([A_i\mathrm{𝟏},A_j\mathrm{𝟏}][\mathrm{𝟏}B_i,A_j\mathrm{𝟏}]+[A_i\mathrm{𝟏},\mathrm{𝟏}B_j]`$ (313)
$`[\mathrm{𝟏}B_i,\mathrm{𝟏}B_j])=i([A_i\mathrm{𝟏},A_j\mathrm{𝟏}][\mathrm{𝟏}B_i,\mathrm{𝟏}B_j])`$
$`=`$ $`i\left([A_i,A_j]\mathrm{𝟏}\mathrm{𝟏}[B_i,B_j]\right).`$
We can easily see that equations (297) are satisfied if both $`A`$ and $`B`$ satisfy $`SU(2)`$ algebra
$$[A_i,A_j]=iϵ_{ijk}A_k\text{and}[B_i,B_j]=iϵ_{ijk}B_k.$$
(314)
General representation of Lorentz group will be labeled by two "spin" degrees of freedom $`(j,j^{})`$. Note that the first equation in (310) is identical to rules for addition of spin in non-relativistic quantum mechanics. However, this “direct product” isn’t a real direct product since parity can’t be expressed as a direct product of two $`SU(2)`$ operators. To show that, let’s assume the contrary. Parity would then be
$$P=P_1P_2.$$
(315)
Since $`P^{}=P^1=P`$, we have to have
$$P_1P_1=P_2P_2=1.$$
(316)
Let’s assume operators $`\stackrel{}{J}`$ and $`\stackrel{}{K}`$ are irreducible. Inserting (315) into (308)
$`P\stackrel{}{J}P`$ $`=`$ $`(P_1P_2)(\stackrel{}{A}\mathrm{𝟏}+\mathrm{𝟏}\stackrel{}{B})(P_1P_2)`$ (317)
$`P\stackrel{}{K}P`$ $`=`$ $`(P_1P_2)i(\stackrel{}{A}\mathrm{𝟏}\mathrm{𝟏}\stackrel{}{B})(P_1P_2)`$ (318)
we get the system of equations
$`(P_1\stackrel{}{A}P_1)\mathrm{𝟏}+\mathrm{𝟏}(P_2\stackrel{}{B}P_2)`$ $`=`$ $`\stackrel{}{A}\mathrm{𝟏}+\mathrm{𝟏}\stackrel{}{B}`$ (319)
$`(P_1\stackrel{}{A}P_1)\mathrm{𝟏}\mathrm{𝟏}(P_2\stackrel{}{B}P_2)`$ $`=`$ $`\stackrel{}{A}\mathrm{𝟏}+\mathrm{𝟏}\stackrel{}{B}`$ (320)
which would imply
$$(P_1\stackrel{}{A}P_1)\mathrm{𝟏}=\mathrm{𝟏}\stackrel{}{B}\text{and}\mathrm{𝟏}(P_2\stackrel{}{B}P_2)=\stackrel{}{A}\mathrm{𝟏}$$
(321)
which is clearly impossible. So, inserting decomposition (310) into commutation relations (308) we get
$$P(\stackrel{}{A}\mathrm{𝟏})P=\mathrm{𝟏}\stackrel{}{B}P(\mathrm{𝟏}\stackrel{}{B})P=\stackrel{}{A}\mathrm{𝟏}.$$
(322)
Since parity exchanges $`A`$ and $`B`$ eigenstates in $`(a,b)`$ representations, only combinations $`(j,j)`$ will be invariant under parity. However, even if operators $`\stackrel{}{J}`$ and $`\stackrel{}{K}`$ can be decomposed in a direct sum of operators
$$\stackrel{}{J}=\stackrel{}{J}_1\stackrel{}{J}_2=\left(\begin{array}{cc}\stackrel{}{J}_1& 0\\ 0& \stackrel{}{J}_2\end{array}\right)\stackrel{}{K}=\stackrel{}{K}_1\stackrel{}{K}_2=\left(\begin{array}{cc}\stackrel{}{K}_1& 0\\ 0& \stackrel{}{K}_2\end{array}\right),$$
(323)
representation of Lorentz group will be irreducible as long as at least one operator can’t be decomposed in a direct sum. If operators $`\stackrel{}{J}_1`$ and $`\stackrel{}{K}_1`$ belong to the $`(a_1,b_1)`$ representation and $`\stackrel{}{J}_2`$ and $`\stackrel{}{K}_2`$ belong to the $`(a_2,b_2)`$, then their direct sums will also satisfy (297). In block-diagonal form requirement $`P^{}=P`$ gives us
$$P=\left(\begin{array}{cc}P_{11}& P_{12}\\ P_{12}^{}& P_{22}\end{array}\right)$$
(324)
with matrices $`P_{11}`$ and $`P_{22}`$ Hermitean. Adding the condition $`PP=1`$ gives us
$$PP=\left(\begin{array}{cc}P_{11}^2+P_{12}P_{12}^{}& P_{11}P_{12}+P_{12}^{}P_{22}\\ P_{12}^{}P_{11}+P_{22}^{}P_{11}& P_{12}^{}P_{12}+P_{22}^2\end{array}\right)=1.$$
(325)
This has two obvious solutions; first one is $`P_{12}=0`$, $`P_{11}^2=1`$, $`P_{22}^2=1`$. In this case the parity is a direct sum of parities, so this solution doesn’t yield irreducible representation of $`SO(1,3)`$. The other solution is $`P_{12}P_{12}^{}=P_{12}^{}P_{12}=1`$, $`P_{11}=P_{22}=0`$. After reinserting this back into (308) we have
$`\left(\begin{array}{cc}0& P_{12}\\ P_{12}^1& 0\end{array}\right)\left(\begin{array}{cc}\stackrel{}{J}_1& 0\\ 0& \stackrel{}{J}_2\end{array}\right)\left(\begin{array}{cc}0& P_{12}\\ P_{12}^1& 0\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}P_{12}\stackrel{}{J}_2P_{12}^1& 0\\ 0& P_{12}^1\stackrel{}{J}_1P_{12}\end{array}\right)=\stackrel{}{J}`$ (334)
$`\left(\begin{array}{cc}0& P_{12}\\ P_{12}^1& 0\end{array}\right)\left(\begin{array}{cc}\stackrel{}{K}_1& 0\\ 0& \stackrel{}{K}_2\end{array}\right)\left(\begin{array}{cc}0& P_{12}\\ P_{12}^1& 0\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}P_{12}\stackrel{}{K}_2P_{12}^1& 0\\ 0& P_{12}^1\stackrel{}{K}_1P_{12}\end{array}\right)=\stackrel{}{K}.`$ (343)
From this we can see that representations $`(a_1,b_1)`$ and $`(a_2,b_2)`$ have to be of a same dimension. After decomposing $`\stackrel{}{J}_{1,2}`$ and $`\stackrel{}{K}_{1,2}`$ into $`SU(2)`$ products, we get
$`P_{12}^1(\stackrel{}{A}_1\mathrm{𝟏}_1)P_{12}=\mathrm{𝟏}_2\stackrel{}{B}_2P_{12}(\mathrm{𝟏}_2\stackrel{}{B}_2)P_{12}^1=\stackrel{}{A}_1\mathrm{𝟏}_1`$ (344)
$`P_{12}(\stackrel{}{A}_2\mathrm{𝟏}_2)P_{12}^1=\mathrm{𝟏}_1\stackrel{}{B}_1P_{12}^1(\mathrm{𝟏}_1\stackrel{}{B}_1)P_{12}=\stackrel{}{A}_2\mathrm{𝟏}_2`$ (345)
which will be parity-invariant only if $`a_1=b_2`$ and $`b_1=a_2`$. Direct sums of three or more $`SU(2)SU(2)`$ representations will produce only representations of $`SO(1,3)`$ reducible to direct sum of $`(j,j)`$ and $`(j,j^{})(j^{},j)`$ representations.
## Appendix B Representations of Lorentz group
Since all representations of the same type have some similarities, they will also have some similar properties which can be derived and studied jointly.
### B.1 $`(j,0)(0,j)`$ representations
Starting point for these representations is equation (310)
$$\stackrel{}{J}=\stackrel{}{A}\mathrm{𝟏}+\mathrm{𝟏}\stackrel{}{B}\stackrel{}{K}=i\stackrel{}{A}\mathrm{𝟏}i\mathrm{𝟏}\stackrel{}{B}.$$
(346)
By definition, all finite transformations belonging to spin 0 representation must map to identity operator
$$e^{i\stackrel{}{\omega }\stackrel{}{B}}|\mathrm{0\hspace{0.17em}0}=|\mathrm{0\hspace{0.17em}0}$$
(347)
which implies that generators $`\stackrel{}{B}`$ must annihilate $`|\mathrm{0\hspace{0.17em}0}`$ state.
$$\stackrel{}{B}|\mathrm{0\hspace{0.17em}0}0.$$
(348)
Surviving part of generators will then be
$$\stackrel{}{J}=\stackrel{}{A}\mathrm{𝟏},\stackrel{}{K}=i\stackrel{}{A}\mathrm{𝟏}$$
(349)
acting on states $`|jm|\mathrm{0\hspace{0.17em}0}`$. Vector $`|\mathrm{0\hspace{0.17em}0}`$ belongs to scalar representation of Lorentz group which is 1-dimensional, so the direct product $`|jm|\mathrm{0\hspace{0.17em}0}`$ is essentially yhe same as the vector $`|jm`$ and so we can disregard it
$$|jm|\mathrm{0\hspace{0.17em}0}|jm.$$
(350)
Therefore, we can write the matrices for this representation as
$$\stackrel{}{J}\stackrel{}{S},\stackrel{}{K}i\stackrel{}{S}$$
(351)
where matrices $`\stackrel{}{S}`$ are $`2j+1`$-dimensional generator matrices of Lorentz group from appendix C.1.
Similar holds for $`(0,j)`$ representation only here $`\stackrel{}{A}`$ annihilates the states and $`\stackrel{}{B}`$ gives us the generators. From equation (346) we can see that the only difference will be in the sign of $`\stackrel{}{K}`$ generator
$$\stackrel{}{J}\stackrel{}{S},\stackrel{}{K}i\stackrel{}{S},|\mathrm{0\hspace{0.17em}0}|jm|jm$$
(352)
Complete generators in block-matrix form will then be
$$\stackrel{}{J}=\left(\begin{array}{cc}\stackrel{}{S}& 0\\ 0& \stackrel{}{S}\end{array}\right)\stackrel{}{K}=i\left(\begin{array}{cc}\stackrel{}{S}& 0\\ 0& \stackrel{}{S}\end{array}\right)$$
(353)
From (344) and (345) we can see that the parity operator must then be
$$P=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)$$
(354)
Parity eigenstates (for zero momentum) in spin space are
$$u_{0s}=\left(\begin{array}{c}|jm\\ |jm\end{array}\right),v_{0s}=\left(\begin{array}{c}|jm\\ |jm\end{array}\right).$$
(355)
For finite momentum we get spinors $`u_{ps}`$ and $`v_{ps}`$ by applying Lorentz boost to rest-frame spinors
$$u_{ps}S(\stackrel{}{\omega })u_{0s},v_{ps}S(\stackrel{}{\omega })v_{0s}$$
(356)
where $`\stackrel{}{\omega }`$ points in the direction of $`\stackrel{}{p}`$ and has the appropriate magnitude. Note however that they are *not* spin eigenstates. For zero momentum spinors $`u_{0,s}`$ and $`v_{0,s}`$ *are* spin eigenfunctions
$$u_{0,s}=\left(\begin{array}{c}\chi _s\\ \chi _s\end{array}\right)v_{0,s}=\left(\begin{array}{c}\chi _s\\ \chi _s\end{array}\right)$$
(357)
is $`\chi _s`$ are spin eigenstates for spin $`j`$ representation of $`SU(2)`$ group. Now consider any spinor $`\phi _{ps}`$ which is spin eigenstate in direction $`\widehat{n}`$ with some value $`\lambda `$
$$\left(\widehat{n}\stackrel{}{J}\right)\phi _{ps}=\lambda \phi _{ps}.$$
(358)
Applying boost operator to state $`\phi _{ps}`$ gives us (by definition) state $`\phi _{p^{},s^{}}`$
$$\phi _{p^{},s^{}}=e^{i\stackrel{}{\omega }\stackrel{}{K}}\phi _{ps}.$$
(359)
If that state were also spin eigenstate
$`\left(\widehat{n}\stackrel{}{J}\right)\phi _{p^{},s^{}}`$ $`=`$ $`\left(\widehat{n}\stackrel{}{J}\right)e^{i\stackrel{}{\omega }\stackrel{}{K}}\phi _{ps}=e^{i\stackrel{}{\omega }\stackrel{}{K}}\left(\widehat{n}\stackrel{}{J}\right)\phi _{ps}[\widehat{n}\stackrel{}{J},e^{i\stackrel{}{\omega }\stackrel{}{K}}]\phi _{ps}`$ (360)
$`=`$ $`e^{i\stackrel{}{\omega }\stackrel{}{K}}\lambda \phi _{ps}[\widehat{n}\stackrel{}{J},e^{i\stackrel{}{\omega }\stackrel{}{K}}]\phi _{ps}`$ (361)
$`=`$ $`\lambda \phi _{p^{},s^{}}[\widehat{n}\stackrel{}{J},e^{i\stackrel{}{\omega }\stackrel{}{K}}]\phi _{ps}`$ (362)
it would imply that all components of $`J`$ commute with all components of $`K`$ which we know isn’t true.
Since coordinate (“orbital”) part of boost generators commute with spin part of rotation generators, the proper way to create momentum *and* spin eigenstates is to apply only the coordinate part of boost generators to rest frame momentum and spin eigenstates
$$e^{i\stackrel{}{\omega }\stackrel{}{K}_{coord}}\psi _{0,s}(x)=e^{i\stackrel{}{\omega }\stackrel{}{K}_{coord}}\frac{e^{imt}}{(2\pi )^{3/2}}u_{0,s}=\frac{e^{ipx}}{(2\pi )^{3/2}}u_{0,s}\psi _{p,s}(x)$$
(363)
We will call those states $`w_{ps}^\pm `$ (although they don’t depend on momentum)
$$\psi _{p,s}(x)=e^{i\stackrel{}{\omega }\stackrel{}{K}_{coord}}\frac{e^{imt}}{(2\pi )^{3/2}}w_{0,s}=\frac{e^{ipx}}{(2\pi )^{3/2}}w_{p,s}.$$
(364)
Note that momentum and spin dependence factor as they should since translation generator commutes with spin part of rotation generator. This is by no means specific to $`(j,0)(0,j)`$ representations. Since it’s a consequence of commutation relations for operators, it will be valid for all representations. It is generally possible to make the transformation of operators
$$M_{ch}M_P=U^{}M_{ch}U\phi _P=U^{}\phi _{ch}$$
(365)
in which the parity operator is diagonal. This is achieved with the unitary matrix
$$U=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)$$
(366)
where $`1`$ is unit matrix of appropriate dimension, which diagonalizes parity and leave rotation generators unchanged
$$P=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\stackrel{}{J}=\left(\begin{array}{cc}\stackrel{}{S}& 0\\ 0& \stackrel{}{S}\end{array}\right)$$
(367)
but makes boost operators non-diagonal
$$\stackrel{}{K}=i\left(\begin{array}{cc}0& \stackrel{}{S}\\ \stackrel{}{S}& 0\end{array}\right).$$
(368)
Momentum-parity-spin eigenstates in this representation decompose to a direct sum of positive and negative parity part
$$w_{ps}^+=\left(\begin{array}{c}\chi _s\\ 0\end{array}\right),w_{ps}^{}=\left(\begin{array}{c}0\\ \chi _s\end{array}\right).$$
(369)
### B.2 $`(j,j)`$ representations
Lets take a look at the first equation in (346). A look at any quantum mechanics textbook (i. e. , section 3.7) will show that it’s identical to the equation for addition of spin in non-relativistic quantum mechanics. So the spin eigenstates for $`(j,j)`$ representation will then be a sum of $`2j,2j1,\mathrm{},0`$ irreducible spin representations. They are related to spin $`j`$ states through Clebsch-Gordon coefficients $`j_1j_2;jm|j_1j_2;m_1m_2`$
$$|j^{},m=\underset{m_1,m_2}{}jj;j^{}m|jj;m_1m_2|j,m_1|j,m_2j^{}=0,1,\mathrm{},2j$$
(370)
We can in general choose such a basis where rotation generator is reducible
$$\stackrel{}{J}=\left(\begin{array}{cccc}\stackrel{}{J}^{(0)}& 0& 0& \mathrm{}\\ 0& \stackrel{}{J}^{(1)}& 0\mathrm{}& \\ 0& 0& \stackrel{}{J}^{(2)}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right).$$
(371)
This doesn’t give irreducible representations of whole Lorentz group since boost generators $`K`$ won’t be diagonal. A look at the second equation in (346) shows that operators $`K^\pm =K_1\pm iK_2`$ can change values $`(a,b)`$ only by 1. Operator $`K_3`$ multiplies each $`(a,b)`$ state with some factor and so changes the relative sign for sums of states. $`K`$ operators are the sum of the same operators as $`J`$ operators, so the action of both operators will give the same states with different coefficients
$`J|jm`$ $``$ $`J|ja;jb=|jm^{}c_1|ja^{};jb+c_2|ja;jb^{}`$
$`K|jm`$ $``$ $`K|ja;jb=|j^{}m^{}d_1|ja^{};jb+d_2|ja;jb^{}`$ (372)
Since those states are orthogonal, they cannot have the same $`j`$ value; since $`J`$ doesn’t change the $`j`$ value, $`K`$ must. Since the values of $`a`$ and $`b`$ have been changed by 1 or 0, $`K`$ can only raise or lower the $`j`$ value by 1, or in another words, $`K`$ will connect only states with spin $`j^{}`$ and $`j^{}\pm 1`$. In the basis where $`\stackrel{}{J}`$ is diagonal this can be represented in block-matrix form as
$$\stackrel{}{K}=\left(\begin{array}{cccc}0& \stackrel{}{K}_{01}& 0& \mathrm{}\\ \stackrel{}{K}_{10}& 0& \stackrel{}{K}_{12}& \mathrm{}\\ 0& \stackrel{}{K}_{21}& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right).$$
(373)
Since boost generators anti-commutes with $`K`$, action of $`K`$ will change parity of the state
$$P|jm=\lambda |jm,PK|jm=\underset{K}{\underset{}{PKP}}\underset{\lambda |jm}{\underset{}{P|jm}}=\lambda K|jm$$
(374)
Therefore, states with spin $`j^{}`$ and $`j^{}\pm 1`$ will have opposite parity. This also implies there will never be the same number of states with opposite parities.
Again, we’ll call the states obtained by boosts $`u_{ps}`$ and $`v_{ps}`$
$$u_{ps}=e^{i\stackrel{}{\omega }\stackrel{}{K}}u_{0s},v_{ps}=e^{i\stackrel{}{\omega }\stackrel{}{K}}v_{0s},$$
(375)
and spin eigenstates $`w_{ps}^\pm `$. Wave function will then be a superposition of those eigenfunctions
$$\phi ^{(j,j)}(x)=\frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.25em}2}E_p}e^{ipx}\left(\underset{s}{}b_{ps}u_{ps}+\underset{s^{}}{}d_{ps^{}}v_{ps^{}}\right)=\frac{d^3p}{(2\pi )^{3/2}\mathrm{\hspace{0.25em}2}E_p}e^{ipx}\underset{s,𝒫}{}a_{ps}^𝒫w_{ps}^𝒫$$
(376)
with some coefficients which are in general complex functions which depend on momentum, spin and parity of the state. It is important to emphasize that all $`(j,j)`$ wave functions on the quantum level are *complex* numbers which cannot be restricted to be real functions.
### B.3 $`(j,j^{})(j^{},j)`$ representations
This class of representations mixes the properties of former two. Like in the $`(j,0)(0,j)`$ representations, rotation and boost generators will be the direct sum of two irreducible parts of equal dimensions $`(j,j^{})`$; similar to $`(j,j)`$ representation, rotation generators for each $`(j,j^{})`$ part will be a direct sum of generators for different spin, and again, boost generators will be irreducible. Rotations will act separately on states with different spin and parity while the boost will mix them.
### B.4 Negative definite scalar products and norms
Let’s show that states with negative “norms” *are* unavoidable. Consider the action of parity (322) on states belonging to representation $`(j,j)`$ $`|jm|jm^{}`$:
$$P|jm|jm^{}=|jm^{}|jm$$
(377)
States with $`m=m^{}`$ will obviously be invariant; for states with $`mm^{}`$ there are two linear combinations which are parity eigenstates
$$P\frac{1}{\sqrt{2}}\left(|jm|jm^{}\pm |jm^{}|jm\right)=\pm \frac{1}{\sqrt{2}}\left(|jm|jm^{}\pm |jm^{}|jm\right).$$
(378)
For $`(j,j^{})(j^{},j)`$ representation states $`\left(|jm|j^{}m^{}|j^{}n^{}|jn\right)`$ parity exchanges $`mn`$ and $`m^{}n^{}`$. Again, parity invariant combinations will be
$`P\left(|jm|j^{}m^{}|j^{}n^{}|jn\pm |jn|j^{}n^{}|j^{}m^{}|jm\right)`$
$`=\pm \left(|jm|j^{}m^{}|j^{}n^{}|jn\pm |jn|j^{}n^{}|j^{}m^{}|jm\right).`$ (379)
In both cases, we have explicitly constructed states of both negative and positive parity. This construction shows that it’s impossible to have states with either parity without states with opposite parity. Now take the state with negative parity $`\phi `$ and take it’s norm
$$\overline{\phi }|\phi =\underset{A}{}\phi _A^{}P_{AB}\phi _B=\underset{A}{}\phi _A^{}\phi _A$$
(380)
which is by definition negative number *q.e.d.*
## Appendix C Groups and Dirac matrices
### C.1 Representations of $`SU(2)`$ generators
$`SU(2)`$ generators satisfy commutation relations
$$[S_i,S_j]=iϵ_{ijk}S_k$$
(381)
It is sometimes convenient to define raising and lowering operators $`S_+`$ and $`S_{}`$ as
$$S_\pm =S_1\pm iS_2$$
(382)
We choose the basis in which $`S_3`$ is diagonal matrix
$$(S_3)_{kl}=\delta _{k,l}(jk+1)\text{or}S_3=\left(\begin{array}{cccc}j& 0& \mathrm{}& 0\\ 0& j1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& j\end{array}\right).$$
(383)
From commutation relations for $`S_\pm `$ we get
$$S_+|j,m=\sqrt{j(j+1)m(m+1)}|j,m+1$$
(384)
from which we can construct raising operator with only off-diagonal elements
$$(S^+)_{kl}=\delta _{k,l1}\sqrt{j(j+1)l(l1)}\text{or}S_3=\left(\begin{array}{cccc}0& \sqrt{2j}& 0& \mathrm{}\\ 0& 0& \sqrt{2(j1)}& \mathrm{}\\ 0& 0& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right).$$
(385)
$`S^{}`$ will be Hermitean conjugate of $`S^+`$, from which we can calculate $`S_1`$ and $`S_2`$
$$S_1=\frac{S_++S_{}}{2}S_2=\frac{S_+S_{}}{2i}$$
(386)
Lowest-dimensional representation of $`SU(2)`$ generators is spin $`1/2`$ representation where generators are Pauli matrices multiplied by factor $`1/2`$
$$S_i=\frac{\sigma _i}{2}$$
(387)
where Pauli matrices are given by
$$\sigma ^1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\sigma ^2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)\sigma ^3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(388)
$$S_+=\frac{\sigma _1+i\sigma _2}{2}=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)S_{}=\frac{\sigma _1i\sigma _2}{2}=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right).$$
(389)
Spin eigenstates $`|jm`$ in this form are
$$|\frac{1}{2}\frac{1}{2}=\left(\begin{array}{c}1\\ 0\end{array}\right)|\frac{1}{2}\frac{1}{2}=\left(\begin{array}{c}0\\ 1\end{array}\right)$$
(390)
Next representation is spin $`1`$ representation; generators for it are given by
$$S_1=\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 1\\ 0& 1& 0\end{array}\right)S_2=\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}0& i& 0\\ i& 0& i\\ 0& i& 0\end{array}\right)S_3=\left(\begin{array}{ccc}1& 0& 0\\ 0& 0& 0\\ 0& 0& 1\end{array}\right).$$
(391)
Raising and lowering operators are
$$S_+=\left(\begin{array}{ccc}0& \sqrt{2}& 0\\ 0& 0& \sqrt{2}\\ 0& 0& 0\end{array}\right)S_{}=\left(\begin{array}{ccc}0& 0& 0\\ \sqrt{2}& 0& 0\\ 0& \sqrt{2}& 0\end{array}\right),$$
(392)
while spin eigenstates are
$$|\mathrm{1\hspace{0.25em}1}=\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right)|\mathrm{1\hspace{0.25em}0}=\left(\begin{array}{c}0\\ 1\\ 0\end{array}\right)|11=\left(\begin{array}{c}0\\ 0\\ 1\end{array}\right).$$
(393)
Spin $`1`$ eigenstates are can be also thought of as components of vector in spherical coordinates. Generators in spherical coordinates are related to the Cartesian coordinates through
$$S_i^{spher}=VS_i^{cart}V^{}$$
(394)
where $`V`$ is given by
$$V=\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}& \frac{i}{\sqrt{2}}& 0\\ 0& 0& 1\\ \frac{1}{\sqrt{2}}& \frac{i}{\sqrt{2}}& 0\end{array}\right)$$
(395)
where rotation generators in Cartesian coordinates are given by
$$S_1=i\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right)S_2=i\left(\begin{array}{ccc}0& 0& 1\\ 0& 0& 0\\ 1& 0& 0\end{array}\right)S_3=i\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right).$$
(396)
Vectors in spherical coordinates are related to Cartesian coordinates through
$$\left(\begin{array}{c}r_1\\ r_0\\ r_{+1}\end{array}\right)=V^{}\left(\begin{array}{c}r_x\\ r_y\\ r_z\end{array}\right)=\left(\begin{array}{ccc}(r_x+ir_y)/\sqrt{2}& & \\ r_z& & \\ (r_xir_y)/\sqrt{2}& & \end{array}\right)$$
(397)
For spin $`3/2`$ generators are given by
$$S_1=\frac{1}{2}\left(\begin{array}{cccc}0& \sqrt{3}& 0& 0\\ \sqrt{3}& 0& 2& 0\\ 0& 2& 0& \sqrt{3}\\ 0& 0& \sqrt{3}& 0\end{array}\right)S_2=\frac{1}{2}\left(\begin{array}{cccc}0& i\sqrt{3}& 0& 0\\ i\sqrt{3}& 0& 2i& 0\\ 0& 2i& 0& i\sqrt{3}\\ 0& 0& i\sqrt{3}& 0\end{array}\right)$$
$$S_3=\frac{1}{2}\left(\begin{array}{cccc}3& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 3\end{array}\right).$$
(398)
$$S_+=\left(\begin{array}{cccc}0& \sqrt{3}& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& \sqrt{3}\\ 0& 0& 0& 0\end{array}\right)S_{}=\left(\begin{array}{cccc}0& 0& 0& 0\\ \sqrt{3}& 0& 0& \\ 0& 1& 0& 0\\ 0& 0& \sqrt{3}& 0\end{array}\right),$$
(399)
and spin eigenstates are
$$|\frac{3}{2}\frac{3}{2}=\left(\begin{array}{c}1\\ 0\\ 0\\ 0\end{array}\right)|\frac{3}{2}\frac{1}{2}=\left(\begin{array}{c}0\\ 1\\ 0\\ 0\end{array}\right)\text{etc.}$$
(400)
Finally, for spin 2 we have
$$S_1=\left(\begin{array}{ccccc}0& 1& 0& 0& 0\\ 1& 0& \sqrt{3/2}& 0& 0\\ 0& \sqrt{3/2}& 0& \sqrt{3/2}& 0\\ 0& 0& \sqrt{3/2}& 0& 1\\ 0& 0& 0& 1& 0\end{array}\right)S_2=\left(\begin{array}{ccccc}0& i& 0& 0& 0\\ i& 0& \sqrt{3/2}i& 0& 0\\ 0& \sqrt{3/2}i& 0& \sqrt{3/2}i& 0\\ 0& 0& \sqrt{3/2}i& 0& i\\ 0& 0& 0& i& 0\end{array}\right)$$
$$S_3=\left(\begin{array}{ccccc}2& 0& 0& 0& 0\\ 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 2\end{array}\right).$$
(401)
$$S_+=\left(\begin{array}{ccccc}0& 2& 0& 0& 0\\ 0& 0& \sqrt{6}& 0& 0\\ 0& 0& 0& \sqrt{6}& 0\\ 0& 0& 0& 0& 2\\ 0& 0& 0& 0& 0\end{array}\right)S_{}=\left(\begin{array}{ccccc}0& 0& 0& 0& 0\\ 2& 0& 0& 0& 0\\ 0& \sqrt{6}& 0& 0& 0\\ 0& 0& \sqrt{6}& 0& 0\\ 0& 0& 0& 2& 0\end{array}\right),$$
(402)
and
$$|\mathrm{2\hspace{0.25em}2}=\left(\begin{array}{c}1\\ 0\\ 0\\ 0\\ 0\end{array}\right)|\mathrm{2\hspace{0.25em}1}=\left(\begin{array}{c}0\\ 1\\ 0\\ 0\\ 0\end{array}\right)\text{etc.}$$
(403)
### C.2 4D Dirac matrices and spinors
Four-dimensional $`\gamma `$ matrices $`\gamma ^\mu `$ (with $`\mu =0,1,2,3`$) and the matrix $`\gamma ^5`$
$$\gamma ^5=i\gamma ^0\gamma ^1\gamma ^2\gamma ^3=\frac{i}{4}ϵ_{\mu \nu \alpha \beta }\gamma ^\mu \gamma ^\nu \gamma ^\alpha \gamma ^\beta $$
(404)
satisfy the anti-commutator relations
$$\{\gamma ^\mu ,\gamma ^\nu \}=2g^{\mu \nu },g^{\mu \nu }=\text{diag}(1,1,1,1,1).$$
(405)
If we expand this set with the commutator
$$\sigma ^{\mu \nu }=\frac{i}{2}[\gamma ^\mu ,\gamma ^\nu ]$$
(406)
we get the 5 dimensional Clifford algebra (with $`\mu ,\nu =0,1,2,3,5`$) $`\gamma `$ and $`\sigma `$ matrices satisfy commutation relations (in 5D)
$$[\gamma ^\mu ,\sigma ^{\alpha \beta }]=2i\left(g^{\mu \alpha }\gamma ^\beta g^{\mu \beta }\gamma ^\alpha \right).$$
(407)
Hermitean conjugates of $`\gamma ^\mu `$ are given by
$$(\gamma ^0)^{}=\gamma ^0(\gamma ^5)^{}=\gamma ^5(\gamma ^i)^{}=\gamma ^i,i=1,2,3$$
(408)
From this we can calculate the Hermitean conjugates for $`\sigma `$ matrices as well:
$$(\sigma ^{ij})^{}=\sigma ^{ij}(\sigma ^{0i})^{}=\sigma ^{0i}(\sigma ^{50})^{}=\sigma ^{50}(\sigma ^{5i})^{}=\sigma ^{5i}$$
(409)
We can see that either 4D subset of these matrices $`\sigma ^{\mu \nu }=\{\sigma ^{0i},\sigma ^{ij}\}`$ (with $`\gamma ^\mu =\{\gamma ^0,\gamma ^i\}`$) *or* $`\sigma ^{\mu \nu }=\{\sigma ^{5i},\sigma ^{ij}`$ ($`\gamma ^\mu =\{\gamma ^5,\gamma ^i\}`$) will satisfy the appropriate commutation relations so both can be identified with generators of Lorentz transformations. They correspond to different representations of 4D Dirac matrices.
#### C.2.1 Dirac or Parity representation of $`\gamma `$ matrices
Commutation relations don’t determine the $`\gamma `$ and $`\sigma `$ matrices completely. If $`\{\gamma ^\mu ,\sigma ^{\mu \nu }\}`$ is a set of matrices satisfying 4D Clifford algebra, any set related to this one by unitary transformation $`\{U^1\gamma ^\mu U,U^1\sigma ^{\mu \nu }U\}`$ will satisfy the same relations. Dirac originally proposed the representation
$$\gamma ^0=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\stackrel{}{\gamma }=\left(\begin{array}{cc}0& \stackrel{}{\sigma }\\ \stackrel{}{\sigma }& 0\end{array}\right)\gamma ^5=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)$$
(410)
where $`\sigma ^i`$ are Pauli matrices
$$\sigma ^1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\sigma ^2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)\sigma ^3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(411)
$`\sigma `$ matrices in this representation are
$$\sigma ^{0i}=i\left(\begin{array}{cc}0& \sigma ^i\\ \sigma ^i& 0\end{array}\right)=i\alpha ^i\sigma ^{ij}=ϵ_{ijk}\left(\begin{array}{cc}\sigma ^k& 0\\ 0& \sigma ^k\end{array}\right)$$
$$\sigma ^{50}=i\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)=i\gamma ^0\gamma ^5\sigma ^{5i}=i\left(\begin{array}{cc}\sigma ^i& 0\\ 0& \sigma ^i\end{array}\right)=i\gamma ^i\gamma ^5.$$
(412)
Lorentz boost matrices in this representation are
$`S_D(\omega )`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i}{2}}{\displaystyle \frac{\sigma ^{\mu \nu }}{2}}\omega _{\mu \nu }\right)`$ (413)
$`=`$ $`\mathrm{exp}\left\{{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}0& \stackrel{}{\omega }\stackrel{}{\sigma }\\ \stackrel{}{\omega }\stackrel{}{\sigma }& 0\end{array}\right)\right\}=\left(\begin{array}{cc}\mathrm{cosh}{\displaystyle \frac{\omega }{2}}& \widehat{\omega }\stackrel{}{\sigma }\mathrm{sinh}{\displaystyle \frac{\omega }{2}}\\ \widehat{\omega }\stackrel{}{\sigma }\mathrm{sinh}{\displaystyle \frac{\omega }{2}}& \mathrm{cosh}{\displaystyle \frac{\omega }{2}}\end{array}\right)`$ (418)
where $`\widehat{\omega }`$ is unit vector in the direction of $`\stackrel{}{\omega }`$, $`\mathrm{cosh}(\omega /2)=\sqrt{(E+m)/2m}`$, $`\mathrm{sinh}(\omega /2)=\sqrt{(Em)/2m}`$ and $`\widehat{\omega }\stackrel{}{\sigma }=\stackrel{}{p}\stackrel{}{\sigma }/\sqrt{E^2m^2}`$. Expressed through $`p^\mu `$ we get
$$S_D(\omega )=\frac{1}{\sqrt{2m(E+m)}}\left(\begin{array}{cc}E+m& \stackrel{}{p}\stackrel{}{\sigma }\\ \stackrel{}{p}\stackrel{}{\sigma }& E+m\end{array}\right)$$
(419)
In this representation parity matrix $`\gamma ^0`$ is diagonal, but helicity operator $`\gamma ^5`$ and Lorentz boosts $`S_D(\omega )`$ aren’t.
#### C.2.2 Weyl or chiral representation of $`\gamma `$ matrices
If we make the change $`\gamma ^0\gamma ^5`$ and $`\gamma ^5\gamma ^0`$ we get the chiral or Weyl representation
$$\gamma ^0=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\stackrel{}{\gamma }=\left(\begin{array}{cc}0& \stackrel{}{\sigma }\\ \stackrel{}{\sigma }& 0\end{array}\right)\gamma ^5=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)$$
(420)
$$\sigma ^{0i}=i\left(\begin{array}{cc}\sigma ^i& 0\\ 0& \sigma ^i\end{array}\right)=i\alpha ^i\sigma ^{ij}=ϵ_{ijk}\left(\begin{array}{cc}\sigma ^k& 0\\ 0& \sigma ^k\end{array}\right)$$
$$\sigma ^{50}=i\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)=i\gamma ^0\gamma ^5\sigma ^{5i}=i\left(\begin{array}{cc}0& \sigma ^i\\ \sigma ^i& 0\end{array}\right)=i\gamma ^i\gamma ^5.$$
(421)
In chiral representation we get Lorentz boost matrices
$`S_{Ch}(\omega )`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i}{2}}{\displaystyle \frac{\sigma ^{\mu \nu }}{2}}\omega _{\mu \nu }\right)=\mathrm{exp}\left\{{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}\stackrel{}{\omega }\stackrel{}{\sigma }& 0\\ 0& \stackrel{}{\omega }\stackrel{}{\sigma }\end{array}\right)\right\}`$ (424)
$`=`$ $`\left(\begin{array}{cc}\mathrm{cosh}{\displaystyle \frac{\omega }{2}}\widehat{\omega }\stackrel{}{\sigma }\mathrm{sinh}{\displaystyle \frac{\omega }{2}}& 0\\ 0& \mathrm{cosh}{\displaystyle \frac{\omega }{2}}+\widehat{\omega }\stackrel{}{\sigma }\mathrm{sinh}{\displaystyle \frac{\omega }{2}}\end{array}\right)`$ (427)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2m(E+m)}}}\left(\begin{array}{cc}E+m\stackrel{}{p}\stackrel{}{\sigma }& 0\\ 0& E+m+\stackrel{}{p}\stackrel{}{\sigma }\end{array}\right)`$ (430)
This representation has the advantage that all Lorentz transformations $`S_{Ch}`$ as well as chirality operator $`\gamma ^5`$ are diagonal, but parity operator $`\gamma ^0`$ isn’t. The connection between chiral and Dirac $`\gamma `$ matrices is
$$\begin{array}{ccccccc}\hfill \gamma _{ch}=U\gamma _DU^{}\varphi _{ch}=U\varphi _D& & & & & & \\ \hfill \gamma _D=U^{}\gamma _{ch}U\varphi _D=U^{}\varphi _{ch}& & & & & & \end{array}U=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)$$
(431)
or in particular
$$u_D(p,s)=\left(\begin{array}{c}u_R(p,s)+u_L(p,s)\\ u_R(p,s)u_L(p,s)\end{array}\right)v_D(p,s)=\left(\begin{array}{c}u_R(p,s)u_L(p,s)\\ u_R(p,s)+u_L(p,s)\end{array}\right)$$
(432)
where $`u_L(p,s)`$ and $`u_R(p,s)`$ are chiral (bi)spinors belonging to $`(1/2,0)`$ and $`(0,1/2)`$ representations of Lorentz group.
### C.3 Gordon identities
Using the fact that spinor satisfy equations
$$(p/m)u_{p,s}=0\overline{u}_{p,s}(p/m)=0(p/+m)v_{p,s}=0\overline{v}_{p,s}(p/+m)=0$$
(433)
multiplying the first equation in (433) by $`a/`$ (where $`a`$ is arbitrary four vector) from the left, and second from the right and adding them we get
$$\overline{u}_{p,s}[(p/m)a/+a/(q/m)]u_{q,r}=0$$
(434)
which can be rewritten as
$$2m\overline{u}_{p,s}a/u_{q,r}+\overline{u}_{p,s}(\{\frac{p/+q/}{2},a/\}+\{\frac{p/q/}{2},a/])u_{q,r}=0.$$
(435)
Evaluating the commutators and anti-commutators to
$$\{a/,b/\}=2ab[a/,b/]=2i\sigma ^{\mu \nu }a_\mu b_\nu ,$$
(436)
adding or subtracting the proper combinations of equations (433) and after differentiation with respect to $`a_\mu `$ we get Gordon identities:
$`\overline{u}_{p,s}\gamma ^\mu u_{q,r}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\overline{u}_{p,s}\left[(p+q)^\mu +i\sigma ^{\mu \nu }(pq)_\nu \right]u_{q,r}`$ (437)
$`\overline{v}_{p,s}\gamma ^\mu v_{q,r}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\overline{v}_{p,s}\left[(p+q)^\mu +i\sigma ^{\mu \nu }(pq)_\nu \right]v_{q,r}`$ (438)
$`\overline{u}_{p,s}\gamma ^\mu v_{q,r}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\overline{u}_{p,s}\left[(pq)^\mu +i\sigma ^{\mu \nu }(p+q)_\nu \right]v_{q,r}`$ (439)
$`\overline{v}_{p,s}\gamma ^\mu u_{q,r}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\overline{v}_{p,s}\left[(pq)^\mu +i\sigma ^{\mu \nu }(p+q)_\nu \right]u_{q,r}`$ (440)
### C.4 Direct product of representation and operators
If Operator $`𝒜`$ is an operator acting an a Hilbert space of dimension $`n`$ spanned by a complete set of vectors $`|\psi `$ and $``$ is an operator acting on a Hilbert space of dimension $`m`$ spanned by a complete set of vectors $`|\varphi `$, then the direct product of operators $`𝒜`$ and $``$ is defined to be
$$\left(𝒜\right)_{ik,jl}𝒜_{ij}_{kl}$$
(442)
which acts on the direct product of spaces
$$|\psi \varphi _{ij}|\psi _i|\varphi _j$$
(443)
of dimension $`mn`$ like
$$\left(𝒜\right)_{ik,jl}|\psi \varphi _{jl}\left(𝒜_{ij}_{kl}\right)|\psi _j|\varphi _l=\left(𝒜_{ij}|\psi _j\right)\left(_{kl}|\varphi _l\right).$$
(444)
For single space operators (operators in one space multiplied by identity operator in another space) we have
$`[𝒜\mathrm{𝟏},\mathrm{𝟏}]_{ik,jl}`$ $`=`$ $`(𝒜\mathrm{𝟏})_{ik,mn}(\mathrm{𝟏})_{mn,jl}(\mathrm{𝟏})_{ik,mn}(𝒜\mathrm{𝟏})_{mn,jl}`$ (445)
$`=`$ $`𝒜_{im}\mathrm{𝟏}_{kn}_{mj}\mathrm{𝟏}_{nl}_{im}\mathrm{𝟏}_{kn}𝒜_{mj}\mathrm{𝟏}_{nl}`$
$`=`$ $`(𝒜_{im}_{mj})(\mathrm{𝟏}_{kn}\mathrm{𝟏}_{nl})(_{im}𝒜_{mj})(\mathrm{𝟏}_{kn}\mathrm{𝟏}_{nl})`$
$`=`$ $`[𝒜,]_{ij}\mathrm{𝟏}_{kl}=([𝒜,]\mathrm{𝟏})_{ik,jl}`$
Following the same procedure, operators in different spaces commute:
$`[𝒜\mathrm{𝟏},\mathrm{𝟏}]_{ik,jl}`$ $`=`$ $`(𝒜\mathrm{𝟏})_{ik,mn}(\mathrm{𝟏})_{mn,jl}(\mathrm{𝟏})_{ik,mn}(𝒜\mathrm{𝟏})_{mn,jl}`$ (446)
$`=`$ $`𝒜_{im}\mathrm{𝟏}_{kn}\mathbf{\hspace{0.25em}1}_{mj}_{nl}\mathrm{𝟏}_{im}_{kn}𝒜_{mj}\mathrm{𝟏}_{nl}`$
$`=`$ $`(𝒜_{im}\mathrm{𝟏}_{mj})(\mathrm{𝟏}_{kn}_{nl})(\mathrm{𝟏}_{im}𝒜_{mj})(_{kn}\mathrm{𝟏}_{nl})`$
$`=`$ $`𝒜_{ij}_{kl}𝒜_{ij}_{kl}=0`$
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# Spin-1 Antiferromagnetic Heisenberg Chains in an External Staggered Field
## 1 Introduction
In recent times there has arisen a great experimental - and theoretical - interest in a new class of magnetic materials of the general composition $`R_2`$BaNiO<sub>5</sub>, with $`R`$ one of the magnetic rare-earth ions (typically: $`R=`$Nd or $`R=`$Pr). These materials are obtained by substitution from the reference compound Y<sub>2</sub>BaNiO<sub>5</sub>, a highly one-dimensional compound with negligible interactions among the spin-$`1`$ Ni<sup>2+</sup> linear chains. For this reason Y<sub>2</sub>BaNiO<sub>5</sub> is generally considered as an almost ideal example of an $`S=1`$ Haldane-gap system, with a spin gap of: $`\mathrm{\Delta }_0=0.41048(2)`$ in units of the AFM intrachain exchange coupling .
As the magnetic $`R^{3+}`$ ions order antiferromagnetically below a certain Néel temperature $`T_N`$, the $`R_2`$BaNiO<sub>5</sub>’s have been modelled, to a first approximation, as a set of $`S=1`$ chains with a negligible interchain coupling (as compared with the intrachain one), acted upon by an effective commensurate staggered field roughly proportional to the sublattice magnetization of the $`R^{3+}`$ lattice, and hence increasing when the temperature is decrease below $`T_N`$.
The above experimental scenario has motivated a renewal of theoretical activity on the model, that is nonetheless already quite interesting “per se”, of an integer-spin AFM Heisenberg chain coupled to an external commensurate staggered field, that can be described by the model Hamiltonian:
$$=\underset{i}{}\{J𝐒_i𝐒_{i+1}+(1)^i𝐇_s𝐒_i\}$$
(1)
where: $`𝐒_i^2=S(S+1)`$ with $`S`$ an integer (we set $`\mathrm{}=1`$ henceforth, and take $`S=1`$ for the $`Ni^{2+}`$ chains), $`J>0`$ and $`𝐇_s`$ is the external staggered field in appropriate units (see, e.g., Ref. for details).
An extensive DMRG study of the model of Eq.$`(1)`$ has been performed in Ref., where very accurate results were reported for the staggered magnetization curve, the spin gaps, the static correlation functions and the correlation lengths in both the longitudinal and transverse (with respect to the direction of the field) channels. The authors of Ref. made instead an analytic study of the model beginning with the familiar mapping of the Hamiltonian $`(1)`$ onto a nonlinear sigma-model (NL$`\sigma `$M). What they discussed very accurately was actually a related and somewhat more phenomenological model in which the strict NL$`\sigma `$M constraint is softened, then replacing the original NL$`\sigma `$M with a theory of the Ginzburg-Landau type parametrized by an adequate set of adjustable parameters (see Ref. for details).
One of the purposes of the present paper is to investigate carefully what are the resulting similarities and/or differences when the NL$`\sigma `$M constraint is not softened but enforced consistently at each level of approximation. We will report here only results at the tree-level of a loop expansion , i.e. essentially at the the mean-field (MFT) level, of the partition function of the model, supplementing them however with a stability analysis, and deferring a systematic evaluation of higher loop corrections, that are considerably more involved, to a forthcoming paper .
Our second purpose is to assess the validity in the present context of the so-called single-mode approximation (SMA) that did prove beyond doubt its validity previously but in rather different contexts . We will provide here the explicit proof of the fact that indeed the SMA is not applicable to discuss the elementary excitation spectrum in the longitudinal channel, a claim that we had already put forward some time ago , without providing however there an explicit proof.
The paper is organized as follows. In Sect.$`2`$ we state the essentials of the general formalism and derive the saddle-point approximation in the presence of a general external source field. This is needed in order to set up the consistent scheme of calculation of the propagators at the mean-field level that is reported in Sect.$`3`$. In Sect.$`4`$ we study the analytic structure of the propagators, and notably of the longitudinal one, at the physical saddle-point, i.e. when the source field becomes the staggered static field of Eq. $`(1)`$. The final Sect.$`5`$ is devoted to a discussion of our results and to a detailed comparison with previous theoretical approaches, as well as to a discussion of some as yet unsolved problems that are posed by the experimental scenario that has been outlined at the beginning. A useful sum rule for the propagators of the NL$`\sigma `$M field is derived in the Appendix.
## 2 Saddle-point approximation for a general <br>source field
Under the Haldane mapping : $`𝐒_i(1)^iS𝐧_i+𝐥_i`$, $`𝐧_i^2=1`$, where $`𝐧_i`$ represents the slowly-varying local staggered magnetization and $`𝐥_{i\text{ }}`$is the local generator of angular momentum, the Zeeman term of Eq.$`(1)`$ becomes : $`_i(1)^1𝐇_s𝐒_iS_i𝐧_i𝐇_s+_i(1)^1𝐥_i𝐇_s`$. In the continuum limit the second term becomes a total derivative that can be neglected if we adopt periodic boundary conditions on the chain.
Going then to the continuum limit, integrating out the fluctuation field $`𝐥`$ and implementing the NL$`\sigma `$M constraint $`𝐧^2(𝐱)=1`$ ($`𝐱=(x,\tau )`$ with $`\tau `$ the Euclidean time) with the aid of a Lagrange multiplier $`\lambda =\lambda (𝐱)`$ , we obtain the partition function: $`=Tr\{\mathrm{exp}[\beta ]\}`$ of the model in the continuum limit as the path-integral:
$$=[𝒟𝐧]\left[\frac{𝒟\lambda }{2\pi }\right]\mathrm{exp}(S_{eff})$$
(2)
where the effective action $`S_{eff}`$ is given by:
$$S_{eff}=𝑑𝐱\left\{_E(𝐱)S𝐇_s𝐧(𝐱)i\lambda (𝐱)(𝐧^2(x)1)\right\}$$
(3)
where: $`𝑑𝐱`$ $`=𝑑x_0^\beta 𝑑\tau `$ and the Euclidean Lagrangian is given by:
$$_E(𝐱)=\frac{1}{2gc}\left(c^2|_x𝐧|^2+|_\tau 𝐧|^2\right)$$
(4)
and the NL$`\sigma `$M mapping predicts: $`g=2/S`$ for the coupling constant and: $`c=2JSa`$ (with $`a`$ the lattice constant) for the spin-wave velocity.
Now we promote $``$ to a generating functional $`[𝐉]`$ by replacing $`S_{eff}`$ with:
$$S[𝐉]=𝑑𝐱\left\{_E(𝐱)S𝐉(𝐱)𝐧(𝐱)i\lambda (𝐱)(𝐧^2(x)1)\right\}$$
(5)
and we will set: $`𝐉=𝐇_s`$ only at the end of the calculations.
Altogether (after an integration by parts):
$$S[𝐉]=\frac{1}{2}𝑑𝐱𝑑𝐱^{}𝐧(𝐱)G_0^1(𝐱,𝐱^{})𝐧(𝐱^{})S𝑑𝐱𝐉(𝐱)𝐧(𝐱)+i𝑑𝐱\lambda (𝐱)$$
(6)
and $`G_0`$ solves, with the appropriate (Matsubara-Bose) boundary conditions the equation:
$$\frac{1}{gc}(c^2_x^2+_\tau ^2+2igc\lambda (𝐱))G_0(𝐱,𝐱^{})=\delta ^{(2)}(𝐱𝐱^{})$$
(7)
Performing now the linear shift: $`𝐧(𝐱)=𝐧^{}(x)+𝐚(𝐱)`$, with:
$$𝐚(𝐱)=S𝑑𝐱^{}G_0(𝐱,𝐱^{})𝐉(𝐱^{})$$
(8)
we obtain:
$`S[𝐉]`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑𝐱𝑑𝐱^{}𝐧^{}(𝐱)G_0^1(𝐱,𝐱^{})𝐧^{}(𝐱^{})}+i{\displaystyle 𝑑𝐱\lambda (𝐱)}`$ (9)
$`{\displaystyle \frac{1}{2}}S^2{\displaystyle 𝑑𝐱𝑑𝐱^{}𝐉(𝐱)G_0(𝐱,𝐱^{})𝐉(𝐱^{})}`$
Now we can integrate out the field $`𝐧^{}`$, obtaining:
$$[𝐉]\left[\frac{𝒟\lambda }{2\pi }\right]\mathrm{exp}(S[\lambda ;𝐉])$$
(10)
where:
$$S[\lambda ;𝐉]=\frac{3}{2}Tr\{\mathrm{ln}(G_0^1)\}+i𝑑𝐱\lambda (𝐱)\frac{1}{2}S^2𝑑𝐱𝑑𝐱^{}𝐉(𝐱)G_0(𝐱,𝐱^{})𝐉(𝐱^{})$$
(11)
We analyze now what are the general features of a saddle-point approximation made in the presence of an arbitrary (space-time dependent) source field $`𝐉(𝐱)`$. We will also evaluate here the propagators at the mean-field level. The saddle point will be determined by the equation:
$$\left(\frac{\delta S[\lambda ;𝐉]}{\delta \lambda (𝐱)}\right)_𝐉=0$$
(12)
where $`(..)_𝐉`$ means that we (functionally) differentiate while keeping $`𝐉`$ constant. As:
$$\frac{\delta G_0(𝐱^{},𝐱^{\prime \prime })}{\delta \lambda (𝐱)}=2iG_0(𝐱^{},𝐱)G_0(𝐱,𝐱^{\prime \prime })$$
(13)
and:
$$i\frac{\delta S}{\delta \lambda (𝐱)}=3G_0(𝐱,𝐱)+S^2𝑑𝐲𝑑𝐲^{}G_0(𝐲,𝐱)G_0(𝐱,𝐲^{})[𝐉(𝐲)𝐉(𝐲^{})]1$$
(14)
we find the saddle-point equation in the form:
$$3G_0(𝐱,𝐱)+S^2𝑑𝐲𝑑𝐲^{}G_0(𝐲,𝐱)G_0(𝐱,𝐲^{})[𝐉(𝐲)𝐉(𝐲^{})]=1$$
(15)
Eq.(15) will determine then a space-time dependent saddle-point that will be a functional of $`𝐉`$ as well: $`\lambda =\lambda ^{}[𝐱;𝐉]`$.
In Mean-Field Theory (MFT) one approximates $`[𝐉]`$ as:
$$[𝐉]\mathrm{exp}(S[\lambda ^{};𝐉])$$
(16)
and hence the connected two-point propagators:
$$G_c^{\alpha \beta }(𝐱,𝐱^{})=S^2\{T(n^\alpha (𝐱)n^\beta (𝐱^{})n^\alpha (𝐱)n^\beta (𝐱^{})\}=\frac{\delta ^2\mathrm{ln}[𝐉]}{\delta J^\alpha (𝐱)\delta J^\beta (𝐱^{})}$$
(17)
will be given by:
$$G_c^{\alpha \beta }(𝐱,𝐱^{})\frac{\delta ^2S[\lambda ^{};𝐉]}{\delta J^\alpha (𝐱)\delta J^\beta (𝐱^{})}$$
(18)
where now, when taking functional derivatives, one has to consider not only the explicit dependence of $`S`$ on $`𝐉`$ but also the implicit one through $`\lambda ^{}`$. We find then:
$$\frac{\delta S}{\delta J^\alpha (𝐱)}=𝑑𝐲\left(\frac{\delta S}{\delta \lambda (𝐲)}\right)_𝐉\frac{\delta \lambda (𝐲)}{\delta J^\alpha (𝐱)}+\left(\frac{\delta S}{\delta J^\alpha (𝐱)}\right)_\lambda $$
(19)
But the first term on the r.h.s. vanishes identically at the saddle point, and so:
$$\frac{\delta S}{\delta J^\alpha (𝐱)}=\left(\frac{\delta S}{\delta J^\alpha (𝐱)}\right)_\lambda =\frac{1}{2}S^2𝑑𝐲[G_0(𝐱,𝐲)+G_0(𝐲,𝐱)]J^\alpha (𝐲)$$
(20)
which provides also the mean-field equation for the (staggered) “magnetization” induced by the source field $`𝐉`$.
Proceeding one step further we find eventually:
$`G_c^{\alpha \beta }(𝐱,𝐱^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}S^2[G_0(𝐱,𝐱^{})+(𝐱𝐱^{})]\delta ^{\alpha \beta }`$ (21)
$`+S^2{\displaystyle }d𝐲d𝐲^{}[G_0(𝐱,𝐲)G_0(𝐲,𝐲^{})+(𝐱𝐲^{})]J^\alpha (𝐲)\left(i{\displaystyle \frac{\delta \lambda (𝐲^{})}{\delta J^\beta (𝐱^{})}}\right)`$
This is the general structure of the mean-field propagators for a general source field $`𝐉`$.
The functional derivative of $`\lambda `$ on the r.h.s. of Eq.(21) is determined by:
$$0=\frac{\delta }{\delta J^\alpha }\left(\frac{\delta S}{\delta \lambda }\right)_𝐉=\left(\frac{\delta ^2S}{\delta \lambda \delta \lambda ^{}}\right)_𝐉\frac{\delta \lambda ^{}}{\delta J^\alpha }+\left(\frac{\delta }{\delta J^\alpha }\left(\frac{\delta S[\lambda ;𝐉]}{\delta \lambda }\right)_𝐉\right)_\lambda $$
(22)
where the second term represents the variation of $`(\delta S/\delta \lambda )_𝐉`$ w.r.t. its explicit dependence on $`𝐉`$. As all the quantities in brackets have to be evaluated at the saddle point, this equation becomes an inhomogeneous linear integral equation for $`\delta \lambda /\delta J^\alpha `$ whose kernel is:
$`H(𝐱,𝐱^{})`$ $`=`$ $`\left({\displaystyle \frac{\delta ^2S}{\delta \lambda (𝐱)\delta \lambda (𝐱^{})}}\right)_𝐉|_{\lambda =\lambda ^{}}`$
$`=`$ $`6\mathrm{\Gamma }(𝐱,𝐱^{})+4S^2G_0(𝐱^{},𝐱){\displaystyle 𝑑𝐲𝑑𝐲^{}G_0(𝐲,𝐱^{})G_0(𝐱,𝐲^{})[𝐉(𝐲)𝐉(𝐲^{})]}`$
where $`\mathrm{\Gamma }`$ is the “polarization bubble”:
$$\mathrm{\Gamma }(𝐱,𝐱^{})=G_0(𝐱,𝐱^{})G_0(𝐱^{},𝐱)$$
(24)
The integral equation for $`\delta \lambda /\delta J^\alpha (𝐱)`$ reads then:
$$𝑑𝐲H(𝐱,𝐲)\left(i\frac{\delta \lambda (𝐲)}{\delta J^\alpha (𝐱^{})}\right)=2S^2G_0(𝐱,𝐱^{})𝑑𝐲G_0(𝐱,𝐲)J^\alpha (𝐲)$$
(25)
In the next Section we will specialize the results obtained here to the physical case $`𝐉=const.=𝐇_s`$, which will determine the physically relevant saddle point.
## 3 Results at the physical saddle point
When $`𝐉=const.=𝐇_s`$, the associated saddle point will correspond also to $`\lambda =const.`$ Setting then: $`2igc\lambda =const.=c^2\xi ^2`$, translational invariance will be restored and the following results can be easily derived :
1. The saddle-point condition is:
$$3G_0(\mathrm{𝟎})+S^2𝐇_s^2[\stackrel{~}{G_0}(\mathrm{𝟎})]^2=1$$
(26)
where $`\stackrel{~}{G_0}(\mathrm{𝟎})=\stackrel{~}{G_0}(𝐪=\mathrm{𝟎})`$ ($`𝐪=(q,\mathrm{\Omega }_n=2\pi n/\beta `$) and $`\stackrel{~}{G_0}(𝐪)`$, the Fourier transform of $`G_0(𝐱)`$, is given by:
$$\stackrel{~}{G_0}(𝐪)=\frac{gc}{\mathrm{\Omega }_n^2+c^2(q^2+\xi ^2)}$$
(27)
whence: $`\stackrel{~}{G_0}(0)=g\xi ^2/c.`$ Explicitly, at $`T=0`$ :
$$\frac{3g}{2\pi }\mathrm{ln}\left\{\mathrm{\Lambda }\xi +\sqrt{1+(\mathrm{\Lambda }\xi )^2}\right\}=1\left(\frac{Sg}{c}\right)^2𝐇_s^2\xi ^4$$
(28)
and the cutoff can disposed of by fitting it to the zero-field gap $`\mathrm{\Delta }_0=c\xi ^1`$ that is know from the DMRG studies. Eq.($`28`$) will determine then the field dependence of the correlation length $`\xi `$, and it is clear that $`\xi `$ will depend quadratically on the field.
2. The magnetization is given by:
$$𝐦_s=S𝐧_s=S^2\stackrel{~}{\mathrm{\Delta }}(\mathrm{𝟎})𝐇_s=\frac{gS^2\xi ^2}{c}𝐇_s$$
(29)
Comparison with the DMRG data of Ref. shows a slight overestimate of the values of the magnetization for small fields, but the agreement becomes better and better as the field increases (see Fig.$`1`$).
Note that, in view of this equation, the saddle-point condition can be written also as:
$$3G_0(\mathrm{𝟎})+\left(\frac{m_s}{S}\right)^2=1$$
(30)
3. While the transverse susceptibility is given by: $`\chi ^T=m_s/H_s=gS^2\xi ^2/c`$, the longitudinal one is given by the full derivative: $`\chi ^L=dm_s/dH_s`$, i.e. by:
$$\chi ^L=\chi ^T\left\{1+(2H_s/\xi )\frac{d\xi }{dH_s}\right\}=\chi ^T\left\{1+\frac{d(\mathrm{ln}(\xi ^2))}{d(\mathrm{ln}H_s)}\right\}$$
(31)
The derivative on the r.h.s. can be obtained by differentiating the saddle-point equation w.r.t. $`H_s`$, and the explicit expression for $`\chi ^L`$ is:
$$\chi ^L=\chi ^T\left\{1+\frac{2\pi }{3g}\left(\frac{2gS\xi ^2H_s}{c}\right)^2\frac{\sqrt{1+(\mathrm{\Lambda }\xi )^2}}{\mathrm{\Lambda }\xi }\right\}^1$$
(32)
(clearly exhibiting: $`\chi ^L<\chi ^T`$ always, the two coinciding only when $`H_s=0`$). In the limit $`\mathrm{\Lambda }\xi 1`$, Eq.$`(32)`$ reduces then to:
$$\chi ^L=\chi ^T\frac{1}{1+\frac{2\pi }{3g}(2gS\xi ^2H_s/c)^2}$$
(33)
4. In the translationally-invariant case, Eq.(25) for $`\delta \lambda /\delta J^\alpha `$ can be rewritten, using the equation for the magnetization, as:
$$𝑑𝐲H(𝐱𝐲)\left(i\frac{\delta \lambda (𝐲)}{\delta J^\alpha (𝐱^{})}\right)=2m_s^\alpha G_0(𝐱𝐱^{})$$
(34)
Making then the “Ansatz”:
$$i\frac{\delta \lambda (𝐱)}{\delta J^\alpha (𝐱^{})}=m_s^\alpha X(𝐱𝐱^{})$$
(35)
(Hence: $`\delta \lambda /\delta J^\alpha =0`$ in the directions orthogonal to the field) Eq.(34) reduces to the following equation for $`X`$:
$$𝑑𝐲H(𝐱𝐲)X(𝐲𝐱^{})=2G_0(𝐱𝐱^{})$$
(36)
that can be solved by Fourier transforming it, thus yielding:
$$\stackrel{~}{X}(𝐪)=2\frac{\stackrel{~}{G_0}(𝐪)}{\stackrel{~}{H}(𝐪)}$$
(37)
where, now:
$$\stackrel{~}{H}(𝐪)=6\stackrel{~}{\mathrm{\Gamma }}(𝐪)+4\left(\frac{m_s}{S}\right)^2\stackrel{~}{G_0}(𝐪)$$
(38)
5. The longitudinal connected propagator is given by:
$$G_c^L(𝐱𝐱^{})=S^2G_0(𝐱𝐱^{})2m_s^2𝑑𝐲G_0(𝐱𝐲)X(𝐲𝐱^{})$$
(39)
while:
6. The transverse propagator (that has no nonconnected parts) is simply given by:
$$G_c^T(𝐱𝐱^{})=S^2G_0(𝐱𝐱^{})$$
(40)
Going to Fourier space:
$$\stackrel{~}{G}_c^T(𝐪)=S^2\stackrel{~}{G_0}(𝐪)$$
(41)
which is (see Eq.($`27`$)) a free boson propagator that, when analytically continued to the real axis, has simple poles at: $`\omega =\pm \epsilon (q)`$ with:
$$\epsilon (q)=\sqrt{c^2q^2+\mathrm{\Delta }_0^2},\text{ }\mathrm{\Delta }_0=c\xi ^1$$
(42)
Notice also that:
$$\chi ^T=\stackrel{~}{G}_c^T(\mathrm{𝟎})$$
(43)
Explicitly, in Fourier space, one finds: $`\stackrel{~}{G_c^L}(𝐪)=S^2\stackrel{~}{G_0}(𝐪)2m_s^2\stackrel{~}{G_0}(𝐪)\stackrel{~}{X}(𝐪)`$ and, with some algebra:
$$\stackrel{~}{G}_c^L(𝐪)=S^2\stackrel{~}{G_0}(𝐪)\frac{3\stackrel{~}{\mathrm{\Gamma }}(𝐪)}{3\stackrel{~}{\mathrm{\Gamma }}(𝐪)+2(m_s/S)^2\stackrel{~}{G_0}(𝐪)}\stackrel{~}{G}_c^T(𝐪)\frac{3\stackrel{~}{\mathrm{\Gamma }}(𝐪)}{3\stackrel{~}{\mathrm{\Gamma }}(𝐪)+2(m_s/S)^2\stackrel{~}{G_0}(𝐪)}$$
(44)
$`\stackrel{~}{\mathrm{\Gamma }}(𝐪)`$ is the convolution of two $`\stackrel{~}{G_0}`$’s that can be evaluated explicitly as:
$`\stackrel{~}{\mathrm{\Gamma }}(𝐪)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(gc)^2{\displaystyle }{\displaystyle \frac{dk}{2\pi }}\mathrm{coth}(\beta \epsilon (k)/2){\displaystyle \frac{1}{\epsilon (k)\epsilon (k+q)}}\times `$ (45)
$`\left\{{\displaystyle \frac{\epsilon (k+q)+\epsilon (k)}{\mathrm{\Omega }_n^2+(\epsilon (k+q)+\epsilon (k))^2}}+{\displaystyle \frac{\epsilon (k+q)\epsilon (k)}{\mathrm{\Omega }_n^2+(\epsilon (k+q)\epsilon (k))^2}}\right\}`$
The analytic continuation in frequency ($`\mathrm{\Omega }_niz`$) has a branch-cut along the entire real axis for all $`\beta <+\mathrm{}`$. The discontinuity across the branch-cut vanishes however exponentially with temperature in the range: $`2\mathrm{\Delta }_0(q)<\mathrm{Re}(z)<+2\mathrm{\Delta }_0(q)`$, where: $`\mathrm{\Delta }_0(q)=c\sqrt{(q/2)^2+\xi ^2}=\epsilon (q/2)`$ and, right at $`T=0`$, the branch cut extends only from $`\mathrm{}`$ to $`2\mathrm{\Delta }_0(q)`$ and from $`+2\mathrm{\Delta }_0(q)`$ to $`+\mathrm{}`$. The reason for this is that the second term in curly brackets vanishes exponentially for $`\beta \mathrm{}`$. The same term vanishes also, irrespective of temperature, for $`q=0`$. For this reason we expect also $`\stackrel{~}{\mathrm{\Gamma }}(𝐪)`$ to be an essentially positive quantity at all temperatures. According also to Eq.$`(27)`$, this implies that $`\stackrel{~}{H}(𝐪)`$ is a positive-definite quantity, and hence that $`H(𝐱𝐱^{})`$ is a positive-definite kernel.
Now, the kernel $`H`$ determines actually also the stability of the saddle point. Indeed, by expanding the action (11) quadratically around the (physical) saddle point in terms of: $`\delta \lambda (𝐱)=\lambda (𝐱)\lambda ^{}`$, we obtain
$$S[\lambda ]=S[\lambda ^{}]+\frac{1}{2}𝑑𝐱𝑑𝐱^{}\delta \lambda (𝐱)H(𝐱𝐱^{})\delta \lambda (𝐱^{})$$
(46)
and the positivity of $`H`$ will then guarantee that the physical saddle point is indeed differentially stable, i.e. a local minimum.
We analyze now the asymptotic behaviour of the longitudinal propagator for both large and small values of $`|𝐪|`$.
While one sees immediately that: $`\stackrel{~}{\mathrm{\Delta }}(𝐪)|𝐪|^2`$ for large $`|𝐪|`$, in the same limit : $`\stackrel{~}{\mathrm{\Gamma }}(𝐪)\mathrm{ln}(|𝐪|^2)/|𝐪|^2`$, instead. Therefore: $`\stackrel{~}{G_c^L}(𝐪)`$ $``$ $`\stackrel{~}{G_c^T}(𝐪)`$ for large $`|𝐪|`$, and the (very) short-distance behaviour of the two propagators is the same. This implies:
$$\underset{𝐱0}{lim}G_c^L(𝐱)=\underset{𝐱0}{lim}G_c^T(𝐱)=S^2G_0(\mathrm{𝟎})$$
(47)
Recalling that, in the translationally-invariant case, the full longitudinal propagator $`G^L`$ is related to the connected one by: $`G^L=G_c^L+m_s^2`$, we see that, as a consequence: $`3G_0(\mathrm{𝟎})+(m_s/S)^2S^2\{2G^T(\mathrm{𝟎})+G^L(\mathrm{𝟎})\}`$, and that therefore the saddle-point condition can be again read simply as one implementing the constraint on the average, i.e. as:
$$𝐧^2=1$$
(48)
At the opposite end, when $`𝐪\mathrm{𝟎}`$:
$$\stackrel{~}{G}_c^L(\mathrm{𝟎})=\frac{3S^2\stackrel{~}{\mathrm{\Gamma }}(\mathrm{𝟎})\stackrel{~}{G_0}(\mathrm{𝟎})}{3\stackrel{~}{\mathrm{\Gamma }}(\mathrm{𝟎})+2(m_s/S)^2\stackrel{~}{G_0}(\mathrm{𝟎})}$$
(49)
that is markedly different from $`\stackrel{~}{G_c^T}(\mathrm{𝟎})`$, coinciding with the latter only for $`H_s0`$. The same will be true for the small-momentum behaviour of $`\stackrel{~}{G_c^L}(𝐪)`$. The long-distance behaviours of the two propagators will be then definitely different for $`H_s0`$, and so we expect quite different asymptotic behaviours at infinity, i.e. quite different correlation lengths .
On top of that, the equation (cfr, Eq.(43)): $`\chi ^L=\stackrel{~}{G_c^L}(\mathrm{𝟎})`$ provides us also with an independent expression for the longitudinal susceptibility, explicitly showing how it results from both one- and two-magnon contributions (the latter coming from the “polarization bubble” $`\stackrel{~}{\mathrm{\Gamma }}(𝐪)`$).
In explicit terms, we have, at $`T=0`$:
$$\stackrel{~}{\mathrm{\Gamma }}(\mathrm{𝟎})=\frac{(gc)^2}{8\pi }𝑑k\frac{1}{\epsilon (k)^3}=\frac{g^2}{4\pi c}\underset{0}{\overset{\mathrm{}}{}}𝑑k(k^2+\xi ^2)^{\frac{3}{2}}=\frac{(g\xi )^2}{4\pi c}$$
(50)
Inserting this expression (together with the known value of $`\stackrel{~}{G_0}(\mathrm{𝟎})`$) into the equation for $`\stackrel{~}{G}_c^L(\mathrm{𝟎})`$ yields back precisely Eq.$`(33)`$ that had been obtained by letting $`\mathrm{\Lambda }\mathrm{}`$ in the previous equation for the longitudinal susceptibility wherever this did not lead to divergent results, which is precisely the attitude that has been taken in the present calculation. All this shows that (at least) there are no inconsistencies in the MFT approach that has been adopted here.
In the next Section we will discuss in detail the structure of the analytic continuation of the propagators to the complex frequency plane and to the real axis.
## 4 The analytic structure of the propagators
At the present level of approximation the transverse propagator is just a free-boson propagator. Analytic continuation is straightforward, leading to:
$$\stackrel{~}{G_c^T}(q,z)=\frac{gcS^2}{\epsilon ^2(q)z^2}$$
(51)
($`\epsilon (q)=c\sqrt{q^2+\xi ^2}`$). Going to the real axis from above: $`z=\omega +i\eta ,\eta 0^+`$:
$$\mathrm{Im}\stackrel{~}{G_c^T}(q,\omega )=\frac{\pi gcS^2}{2\epsilon (q)}\left\{\delta (\omega \epsilon (q))\delta (\omega +\epsilon (q))\right\}$$
(52)
The spectral weight function is then fully exhausted by single poles at $`\omega =\pm \epsilon (q)`$ which is the structure required for the applicability of the SMA. The relation, which is a direct consequence of the SMA: $`\chi _T=Sgc/(\mathrm{\Delta }_T)^2`$, with: $`\mathrm{\Delta }_T=\mathrm{\Delta }_0=c\xi ^1`$, is obeyed exactly, at this level of approximation, in the transverse channel.
Let’s turn now to the longitudinal propagator, and let’s begin by looking at the polarization bubble $`\stackrel{~}{\mathrm{\Gamma }}(𝐪)=\stackrel{~}{\mathrm{\Gamma }}(q,\mathrm{\Omega }_n)`$. We shall consider for simplicity only the $`T=0`$ limit in which the second term on the r.h.s. of Eq.$`(45)`$ can be neglected. Then:
$$\stackrel{~}{\mathrm{\Gamma }}(𝐪)=\frac{1}{2}(gc)^2\frac{dk}{2\pi }\frac{1}{\epsilon (k)\epsilon (k+q)}\frac{\epsilon (k+q)+\epsilon (k)}{\mathrm{\Omega }_n^2+(\epsilon (k+q)+\epsilon (k))^2}$$
(53)
that we write for short as:
$$\stackrel{~}{\mathrm{\Gamma }}(𝐪)=\frac{1}{2}(gc)^2\frac{dk}{2\pi }\frac{A(k,q)}{\mathrm{\Omega }_n^2+E(k,q)^2}$$
(54)
where:
$$E(k,q)=\epsilon (k+q)+\epsilon (k)$$
(55)
and:
$$A(k,q)=\frac{E(k,q)}{\epsilon (k)\epsilon (k+q)}$$
(56)
Note that :
$$E(k,q)=E(kq,q)$$
(57)
and the same will hold true for $`A(k,q)`$.
Performing now the analytic continuation and going to the real axis from above:
$$\stackrel{~}{\mathrm{\Gamma }}(q,\omega )=\mathrm{\Gamma }_1(q,\omega )+i\mathrm{\Gamma }_2(q,\omega )$$
(58)
where:
$$\mathrm{\Gamma }_1(q,\omega )=\frac{1}{2}(gc)^2\frac{dk}{2\pi }𝒫\left\{\frac{A(k,q)}{E^2(k,q)\omega ^2}\right\}$$
(59)
(“$`𝒫`$” standing for Cauchy principal part) and:
$$\mathrm{\Gamma }_2(q,\omega )=\frac{1}{4}(gc)^2𝑑k\frac{1}{ϵ(k)ϵ(k+q)}[\delta (\omega E(k,q))\delta (\omega +E(k,q))]$$
(60)
that is odd in $`\omega `$ (which implies that $`\mathrm{\Gamma }_1(q,\omega )`$ will be an even function of $`\omega `$) and positive for positive $`\omega `$.
To evaluate $`\mathrm{\Gamma }_2`$ explicitly it will be enough to consider the positive frequency part. Notice that, for any fixed $`q`$: $`\mathrm{min}\{E(k,q)\}=2\mathrm{\Delta }_0(q)`$ ($`\mathrm{\Delta }_0(q)=\epsilon (q/2)`$, $`\mathrm{\Delta }_0(0)=\mathrm{\Delta }_0`$), and therefore: $`\mathrm{\Gamma }_2(q,\omega )=0`$ for $`|\omega |<2\mathrm{\Delta }_0(q)`$, as we know already. Otherwise, it is easy to see graphically that the equation: $`\omega =E(k,q)`$ has two solutions at $`k=k^{}(q,\omega )=k_0q/2`$ and at $`k=k^{}q=k_0q/2`$, where: $`k_0=(\omega /2c)\sqrt{(\omega ^24\mathrm{\Delta }_0^2(q))/(\omega ^2c^2q^2)}`$.
Then we obtain easily, as: $`\epsilon (k)\epsilon (k+q)(E(k,q)/k)=c^2[k\epsilon (k+q)+`$$`(k+q)\epsilon (k)]`$:
$$\mathrm{\Gamma }_2(q,\omega )=\frac{g^2}{4C(q)}\theta (\omega ^24\mathrm{\Delta }_0^2(q))sgn(\omega )$$
(61)
where: $`C(q)=|k^{}\epsilon (k^{}+q)+(k^{}+q)\epsilon (k^{})|`$ and: $`k^{}=k^{}(q,|\omega |)`$. An explicit analytic expression for $`\mathrm{\Gamma }_2`$ can then be written down in general, but it is not especially illuminating, although it can be very useful for numerical calculations. It simplifies greatly for $`q0`$, where we get simply: $`C(q=0)=(|\omega |/2c)\sqrt{\omega ^24\mathrm{\Delta }_0^2}`$.
The (integrable) square-root singularity at the edges of the branch cuts is present at finite $`q`$ as well, and indeed, for $`\omega ^22\mathrm{\Delta }_0^2(q)`$, we obtain, to leading order: $`C(q)\alpha (q)\sqrt{\omega ^24\mathrm{\Delta }_0^2(q)}+𝒪(\omega ^24\mathrm{\Delta }_0^2(q))`$ with: $`\alpha (q)=(\mathrm{\Delta }_0(q)/2c\mathrm{\Delta }_0)\{|\omega |q^2c^2/2\mathrm{\Delta }_0(q)\}.`$
By exploiting the parity of $`\mathrm{\Gamma }_2,\mathrm{\Gamma }_1`$ will be given then, via dispersion relations, by:
$$\mathrm{\Gamma }_1(q,\omega )=2\underset{2\mathrm{\Delta }_0(q)}{\overset{\mathrm{}}{}}\frac{d\omega ^{}}{\pi }\omega ^{}\mathrm{\Gamma }_2(q,\omega ^{})𝒫(\frac{1}{\omega ^2\omega ^2})$$
(62)
which shows that $`\mathrm{\Gamma }_1(q,\omega )`$ will be strictly positive for $`|\omega |<2\mathrm{\Delta }_0(q)`$, a result that will prove to be useful shortly.
Let’s consider now the full longitudinal propagator:
$$\stackrel{~}{G}_c^L(q,\mathrm{\Omega }_n)=gcS^2\frac{(3\stackrel{~}{\mathrm{\Gamma }}(q,\mathrm{\Omega }_n)/2gc)}{(3\stackrel{~}{\mathrm{\Gamma }}(q,\mathrm{\Omega }_n)/2gc)[\mathrm{\Omega }_n^2+\epsilon ^2(q)]+(m_s/S)^2}$$
(63)
Performing the analytic continuation, omitting specification of the label $`q`$ and defining:
$$G(z)=\stackrel{~}{G}_c^L(q,\mathrm{\Omega }_n)/gcS^2,\text{ }\mathrm{\Gamma }(z)=3\stackrel{~}{\mathrm{\Gamma }}(q,\mathrm{\Omega }_n)/2gc,\text{ }\delta =(m_s/S)^2,\text{ }ϵ=\epsilon (q)$$
(64)
we are led to study the analytic structure of a function of the form:
$$G(z)=\frac{\mathrm{\Gamma }(z)}{\mathrm{\Gamma }(z)(ϵ^2z^2)+\delta }$$
(65)
($`0\delta <1,`$ $`ϵ<2\mathrm{\Delta }_0`$). $`\mathrm{\Gamma }(z)`$ will be given by:
$$\mathrm{\Gamma }(z)=\frac{d\omega ^{}}{\pi }\frac{\mathrm{\Gamma }_2(\omega ^{})}{\omega ^{}z}$$
(66)
with $`\mathrm{\Gamma }_2(\omega )`$ having all the properties that have been listed above (odd in $`\omega `$, positive for positive $`\omega `$ and vanishing for $`|\omega |2\mathrm{\Delta }_0`$ (cfr. Eq.($`61`$)), thus producing a branch-cut in $`\mathrm{\Gamma }(z)`$ for real $`z=\omega `$ and $`2\mathrm{\Delta }_0|\omega |<+\mathrm{}`$). We will write: $`\mathrm{\Gamma }(\omega +i0^+)=\mathrm{\Gamma }_1(\omega )+\mathrm{\Gamma }_2(\omega )`$ on the real axis. The analytic properties of $`G(z)`$ will be determined in turn by its spectral weight function.
Going to the real axis we find, for $`z=\omega +i\eta ,\eta >0`$:
$$G(\omega +i\eta )=G_1(\omega )+iG_2(\omega )$$
(67)
where, defining:
$$A(\omega )=(ϵ^2\omega ^2+\eta ^2)\mathrm{\Gamma }_1(\omega )+2\eta \omega \mathrm{\Gamma }_2(\omega )+\delta $$
(68)
$$B(\omega )=2\eta \omega \mathrm{\Gamma }_1(\omega )(ϵ^2\omega ^2+\eta ^2)\mathrm{\Gamma }_2(\omega )$$
(69)
$`G_1`$and $`G_2`$ are given by:
$$G_1(\omega )=\frac{A(\omega )\mathrm{\Gamma }_1(\omega )B(\omega )\mathrm{\Gamma }_2(\omega )}{A^2(\omega )+B^2(\omega )}$$
(70)
and:
$$G_2(\omega )=\frac{A(\omega )\mathrm{\Gamma }_2(\omega )+B(\omega )\mathrm{\Gamma }_1(\omega )}{A^2(\omega )+B^2(\omega )}$$
(71)
Even for $`\eta 0`$, $`G_2`$ needs not vanish when $`\mathrm{\Gamma }_2`$ does. In particular, let’s inspect its structure for $`|\omega |<2\mathrm{\Delta }_0.`$ Sending $`\eta `$ to $`0`$ inside the $`\mathrm{\Gamma }`$’s (and only there) we have, for $`|\omega |<2\mathrm{\Delta }_0`$: $`A(\omega )=(ϵ^2\omega ^2+\eta ^2)\mathrm{\Gamma }_1(\omega )+\delta ,`$ $`B(\omega )=2\eta \omega \mathrm{\Gamma }_1(\omega )`$, leading to:
$$G_1(\omega )=\frac{ϵ^2\omega ^2+\eta ^2+\delta /\mathrm{\Gamma }_1(\omega )}{(ϵ^2\omega ^2+\eta ^2+\delta /\mathrm{\Gamma }_1(\omega ))^2+4\eta ^2\omega ^2}$$
(72)
and:
$$G_2(\omega )=\frac{2\eta \omega }{(ϵ^2\omega ^2+\eta ^2+\delta /\mathrm{\Gamma }_1(\omega ))^2+4\eta ^2\omega ^2}$$
(73)
From the structure of $`G_2(\omega )`$ it is clear that:
$$\underset{\eta 0}{lim}G_2(\omega )=\pi sgn(\omega )\delta (f(\omega ));\text{ }f(\omega )=\omega ^2ϵ^2\frac{\delta }{\mathrm{\Gamma }_1(\omega )}$$
(74)
Remembering that $`\mathrm{\Gamma }_1`$ is an even function of $`\omega `$, we can write: $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_1(\omega ^2),`$ and $`f(\omega )=0`$ will have solutions at $`\omega =\pm ϵ_L`$ with:
$$ϵ_L^2=ϵ^2+\frac{\delta }{\mathrm{\Gamma }_1(ϵ_L^2)}$$
(75)
that reduces (as it should) to $`ϵ^2=\epsilon ^2(q)`$ when $`\delta 0`$. Moreover:
$$\frac{df}{d\omega }=2\omega \frac{df}{d\omega ^2}=2\omega \left[1+\frac{\delta }{\mathrm{\Gamma }_1^2(\omega ^2)}\frac{d\mathrm{\Gamma }_1(\omega ^2)}{d\omega ^2}\right]$$
(76)
According to known formulas, then:
$$\underset{\eta 0}{lim}G_2(\omega )=\gamma \frac{\pi }{2ϵ_L}\{\delta (\omega ϵ_L)\delta (\omega +ϵ_L)\}$$
(77)
where:
$$\gamma =\left\{\left[1+\frac{\delta }{\mathrm{\Gamma }_1^2(\omega ^2)}\frac{d\mathrm{\Gamma }_1(\omega ^2)}{d\omega ^2}\right]_{\omega =ϵ_L}\right\}^1$$
(78)
This proves of course that the longitudinal propagator has (in the range we are examining) simple poles on the real axis at: $`\omega =\pm ϵ_L`$, and the prefactor $`\gamma `$ will give the reduction of the quasiparticle weight w.r.t. the pure bosonic case.
As, for small fields, $`\delta H_s^2`$, $`\gamma `$ will approach $`1`$ quadratically in the field when $`H_s0`$. A numerical plot of the relative quasiparticle weight $`\gamma `$ at $`q=0`$ is presented in Fig. $`2`$. It is a steadily decreasing function of the field, and the quadratic regime near $`H_s=0`$ is confined to a very narrow region of fields. For higher fields there is an intermediate region in which $`\gamma `$ is almost linear in the field, and we find numerically that for $`S=1`$ it saturates in the high-field limit to:
$$\underset{H_s\mathrm{}}{lim}\gamma =\underset{\delta 1}{lim}\gamma 0.279$$
(79)
More than $`70\%`$ of the quasiparticle weight is then lost when the system evolves towards saturation. No significant changes are expected for $`q0`$.
To complete the analysis we have to investigate the range $`|\omega |>2\mathrm{\Delta }_0`$ of $`G_2(\omega )`$ where $`\mathrm{\Gamma }_2(\omega )`$ does not vanish. With some long but straightforward algebra we find:
$$G_2(\omega )=\frac{\delta \mathrm{\Gamma }_2(\omega )}{(ϵ^2\omega ^2)^2(\mathrm{\Gamma }_2(\omega ))^2+(\delta +(ϵ^2\omega ^2)\mathrm{\Gamma }_1(\omega ))^2}$$
(80)
for $`|\omega |>2\mathrm{\Delta }_0`$. Therefore, in this range of frequencies $`G_2`$ will vanish as $`H_s0`$ and we will recover the simple pole structure with the longitudinal and transverse propagators becoming equal. The longitudinal pole will survive also up to saturation, but with a strongly field-dependent strength.
That as the field increases the spectral weight that is lost from the pole gets transferred to the two-magnon continuum $`(80)`$ (and viceversa when the field decreases) is dictated, e.g., by the sum rule:
$$\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\frac{d\omega }{\pi }\omega \mathrm{Im}G_2(\omega )=1$$
(81)
The sum rule $`(81)`$ is just one of the general sum rules connected with the moment expansions of the spectral weight functions that are related to equal-time expectation values of multiple commutators and that have been known for a long time in many-body theory . A proof of the sum rule adapted to the specific context of the NL$`\sigma `$M will be given in the Appendix<sup>1</sup><sup>1</sup>1Note however that while (cfr. Eq.$`(52)`$) all moments exist for the transverse propagator, when $`G_2`$ is given by Eqs.$`(77)`$ and $`(80)`$ only the first moment will exist and all the others will turn out to be divergent. This is just an artifact of the mean-field approximation (see, e.g., the discussion of a similar problem in Ref...
The pole at $`ϵ_L`$ represents the longitudinal magnon. It will be a well defined excitation as long as $`ϵ_L^2(q)<4\mathrm{\Delta }_0^2(q)`$, which we will prove to be the case. It will be higher in energy (as we have proved previously that $`\mathrm{\Gamma }_1(q,\omega )>0`$ for $`|\omega |<2\mathrm{\Delta }_0(q)`$) than the two (degenerate) transverse magnons that both have energy $`ϵ`$, and will become degenerate with the latter when $`\delta 0`$ (the limit in which the full $`SO(3)`$ invariance is restored).
Just as in the transverse case ($`ϵ=\epsilon (q)=\sqrt{c^2q^2+\mathrm{\Delta }_T^2}`$, $`\mathrm{\Delta }_T=c\xi ^1=\mathrm{\Delta }_0`$), we can define a longitudinal gap $`\mathrm{\Delta }_{L\text{ }}`$via: $`\mathrm{\Delta }_{L\text{ }}=ϵ_L(q=0)`$, i.e.:
$$\mathrm{\Delta }_L^2=\mathrm{\Delta }_T^2+\frac{\delta }{\mathrm{\Gamma }_1(0,\mathrm{\Delta }_L)}$$
(82)
In general:
$$\mathrm{\Gamma }_1(0,\omega )=\frac{g^2c}{2}\underset{\eta 0}{lim}\underset{2\mathrm{\Delta }_0(q)}{\overset{+\mathrm{}}{}}\frac{d\omega ^{}}{\pi }\frac{1}{\sqrt{(\omega ^24\mathrm{\Delta }_0^2)}}\frac{\omega ^2\omega ^2}{(\omega ^2+\omega ^2+\eta ^2)^24\omega ^2\omega ^2}$$
(83)
As we are assuming here: $`\mathrm{\Delta }_L<2\mathrm{\Delta }_T`$, the limiting procedure becomes trivial and we obtain:
$$\mathrm{\Gamma }_1(0,\mathrm{\Delta }_L)=\frac{3}{4}g\underset{2\mathrm{\Delta }_T}{\overset{+\mathrm{}}{}}\frac{d\omega }{\pi }\frac{1}{\sqrt{(\omega ^24\mathrm{\Delta }_T^2)}}\frac{1}{\omega ^2\mathrm{\Delta }_L^2}$$
(84)
Eq.($`82`$) can be rewritten in dimensionless form as:
$$\left(\frac{\mathrm{\Delta }_L}{\mathrm{\Delta }_T}\right)^2=1+\frac{4\delta }{3gF(\mathrm{\Delta }_L/\mathrm{\Delta }_T)}$$
(85)
where:
$$F(y)=\underset{2}{\overset{+\mathrm{}}{}}\frac{dx}{\pi }\frac{1}{\sqrt{x^24}}\frac{1}{x^2y^2};\text{ }|y|<2$$
(86)
Explicitly:
$$F(y)=\frac{1}{2\pi y\sqrt{1y^2/4}}\left\{\frac{\pi }{2}\mathrm{tan}^1\left[\frac{2}{y}\sqrt{1y^2/4}\right]\right\}$$
(87)
and, defining: $`y=\mathrm{\Delta }_L/\mathrm{\Delta }_T`$ we obtain, with: $`g=2/S`$:
$$y^2=1+\frac{8m_s^2yS\sqrt{1y^2/4}}{3\left\{1\frac{2}{\pi }\mathrm{tan}^1\left[\frac{2}{y}\sqrt{1y^2/4}\right]\right\}}$$
(88)
The numerical results for the longitudinal gap are reported for $`S=1`$ in Fig. 3, where we report also the results for the transverse gap. Here too the results are in excellent qualitative agreement with the DMRG results, with some small quantitative discrepancies in the low-field regime. We also find that near $`H_s=0`$ the longitudinal gap increases roughly as three times the transverse one (actually: $`lim_{H_s0}(\mathrm{\Delta }_L\mathrm{\Delta }_0)/(\mathrm{\Delta }_T\mathrm{\Delta }_0)3.58`$). For $`H_s\mathrm{}`$ instead:
$$y^2=1+\frac{8yS\sqrt{1y^2/4}}{3\left\{1\frac{2}{\pi }\mathrm{tan}^1\left[\frac{2}{y}\sqrt{1y^2/4}\right]\right\}}$$
(89)
and, for $`S=1`$:
$$\underset{H_s\mathrm{}}{lim}\frac{\mathrm{\Delta }_L}{\mathrm{\Delta }_T}1.855$$
(90)
i.e. for high fields $`\mathrm{\Delta }_L`$ tends to increase slightly less than twice $`\mathrm{\Delta }_T`$. This is in agreement with the DMRG results of Ref.. Notice that, however, the limiting form of Eq.$`(89)`$ tells us immediately that the ratio of the saturation values of the gaps will tend exactly to two in the large-$`S`$ limit, and indeed it is not difficult to see that:
$$\underset{H_s\mathrm{}}{lim}\frac{\mathrm{\Delta }_L}{\mathrm{\Delta }_T}2\frac{a^2}{S^2}$$
(91)
for $`S1`$, with $`a`$ a numerical constant of order $`0.5`$. The relative quasiparticle weight $`\gamma `$ can also be shown to be of order $`S^2`$ in the same limit, i.e. it will vanish when the magnon poles reaches the edge of the continuum.
## 5 Discussion and conclusions
We summarize here our results, comparing them at the same time with those obtained within other approaches. We will list and discuss only a few relevant points, namely:
$`i)`$ In zero field the excitation spectrum consists of the well-known degenerate triplet of massive Haldane bosons with energy $`\epsilon (q)=c\sqrt{q^2+\xi ^2}`$ with a gap $`\mathrm{\Delta }_0=c\xi ^1.`$ For finite fields, instead:
$`ii)`$ The staggered magnetization curve (Eq.$`(29)`$ and Fig.$`1`$) turns out to agree well with the DMRG results of Ref.. As can be deduced directly from Eq. $`(29)`$ and as was discussed in more detail in Ref., the low- and high-field behaviours of the staggered magnetization are respectively:
$$m_s\chi ^TH_s+𝒪(H_s^3)$$
(92)
for $`H_s0`$ and:
$$m_sS\left[1\frac{A}{\sqrt{H_s}}+𝒪(H_s^1)\right]$$
(93)
with $`A`$ a numerical constant (see for more details) for large $`H_s`$. The value obtained in of $`\chi ^T=23.74/J`$ is somewhat higher than the DMRG result of $`\chi ^T=18.5/J`$. We will resume this point shortly below.
Eq.$`(93)`$ shows that the staggered magnetization saturates only asymptotically in the large-field limit. This is what had to be expected, of course, and is a direct consequence of our implementing in a consistent way the NL$`\sigma `$M constraint. The authors of Ref. found instead a magnetization curve that (see Fig.$`1`$) agrees better with the DMRG data than ours in the initial part (i.e. for low fieds), but that disagrees more and more as the field is increased. Even worse, their magnetization saturates at a finite value of the staggered fields and keeps growing indefinitely beyond saturation afterwards. Full saturation implies of course that the fully polarized Néel state should become the exact ground state of the Hamiltonian $`(1)`$, which known to be true only asymptotically for $`H_s\mathrm{}`$. Moreover, an $`\mathrm{𝑜𝑣𝑒𝑟}`$saturated magnetization is clearly not an admissible result.
It appears therefore that softening the NL$`\sigma `$M constraint as done in Ref. can be an admissible procedure in the zero-field limit, that can be also given (see and references therein) some legitimation in the same limit, but that can become more and more dangerous as the field increases, leading eventually to rather unphysical results.
$`iii)`$ The transverse channel is saturated by a single well-defined magnon pole with energy $`ϵ_T(q)=\epsilon (q)`$ and a gap $`\mathrm{\Delta }_T=\mathrm{\Delta }_0.`$The strength of the pole has a rather weak field dependence and the pole persists up to the highest fields. The relation: $`\chi _T=Sgc/(\mathrm{\Delta }_T)^2`$, which is typical of the SMA, is exactly obeyed in this channel at the mean-field level. Both our results and those of Ref. agree quite well with the DMRG results.
$`iv)`$ In the longitudinal channel instead we find also a well defined magnon pole with energy $`ϵ_L(q)`$ and a gap $`\mathrm{\Delta }_L>\mathrm{\Delta }_T`$, but the magnon propagator acquires also a (two-magnon) branch-cut for nonzero staggered fields. The strength of the pole of the propagator at $`ϵ_L`$ is strongly field-dependent and as the field increases it is steadily transferred to the continuum, in agreement with the sum rule $`(81)`$. Near saturation, and for $`S=1`$, the strength of the pole is reduced to less than $`30\%`$ of its zero-field value, while it vanishes completely at saturation in the large-$`S`$ limit. In the same limit the longitudinal pole disappears into the continuum at saturation, while it remains below the continuum (hence a well-defined excitation, although with a strongly reduced intensity) for finite $`S`$. The relation analogous to case $`iii)`$, namely: $`\chi _L=Sgc/(\mathrm{\Delta }_L)^2`$ is badly violated in this case (cfr. Fig.$`3`$) except in the zero-field limit, and this proves that, due to a non negligible contribution from the two-magnon continuum, the SMA is not applicable at all in the longitudinal channel in the presence of a nonvanishing (staggered) field, as was done in (see the close critical comparison with the DMRG results that was made in Ref., and the SMA curve that is reported for comparison in Fig.$`3`$). This makes a great difference with the zero-field case, where there is quite convincing evidence that two- and/or multi-magnon excitations carry a negligible spectral weight, which turns out to be actually exactly zero at the mean-field level.
In conclusion, we would like to point to two somewhat unsatisfactory aspects of the theoretical scenario that has been outlined here and elsewhere -.
First of all, and as we have already remarked, there is excellent qualitative and even semiquantitative agreement between our NL$`\sigma `$M analysis and the DMRG results over the whole range of fields. There remain however some quantitative discrepancies. As to these, we can only stress the fact that the NL$`\sigma `$M mapping is basically a semiclassical expansion starting from the large-$`S`$ limit that is then continued to lower values of the spin. Being so, and if it has to have any validity at all, it should definitely be able to capture all the essential features of the low-energy physics of models such as that of Eq.$`(1)`$, which we believe we have proved to be the case for $`S=1`$. It should then lead definitely to qualitative and semiquantitative agreement, but it cannot be expected to yield also, and it would be actually a totally unexpected bonus if it did, precise quantitative agreement with, e.g., the DMRG results. We believe that, apart from other refinements , not much more than that can be expected from the NL$`\sigma `$M mapping.
In view of what has been just said, however, the quantitative agreement should improve for increasing values of the spin. A test of the NL$`\sigma `$M, in conjunction with a DMRG analysis, on higher-spin chains can be of some interest .
Coming now to the experimental scenario outlined in the Introduction, while essentially all the theoretical models presented so far (including ours) yield quite reasonable agreement between theory and experiment as far as the staggered magnetization curve and the transverse gap are concerned, recent measurements of the longitudinal gap in Nd<sub>2</sub>BaNiO<sub>5</sub> point to some discrepancies between theory and experiment. Namely, the longitudinal gap is found to survive, with $`\mathrm{\Delta }_L>\mathrm{\Delta }_T`$, also above the Néel temperature, i.e. also when $`H_s=0`$, which points to (single-ion) anisotropies that, although could be accomodated very easily within the NL$`\sigma `$M approach, thus leading to an additional spin-gap in the longitudinal channel, are not explicitly accounted for by the models considered so far. Moreover, below $`T_N`$, the longitudinal gap appears to grow more slowly than the transverse one, and this remains a so far unresolved puzzle.
More measurements on compunds of the class $`R_2`$BaNiO<sub>5</sub> are called for to ascertain whether this is a general fact or it is specific to Nd<sub>2</sub>BaNiO<sub>5</sub>. If the former were the case, this would imply that a model involving just a single chain in an effective staggered field, interesting in itself and worth of theoretical interest in its own, although it can account for most of the experimental observations, is not enough to provide a complete description of all the observed low-energy-physics features of such compounds, and that more sophisticated models are needed.
## 6 Appendix. A sum rule for the NL$`\sigma `$M <br>propagators
To cope with the notation adopted in Sect.$`4`$ after Eq.$`(64)`$, let’s consider (setting: $`g=c=S=1`$), in real time now, a classical Lagrangian of the form:
$$(x,t)=\frac{1}{2}\left\{|_t𝐧|^2|_x𝐧|^2V(𝐧)\right\}$$
(94)
with $`V(𝐧)`$ an arbitrary interaction term, e.g. one implementing in some appropriate limit the NL$`\sigma `$M constraint $`𝐧^2(x,t)=1`$. The canonical momentum density is defined by:
$$\pi (x,t)=\frac{}{(_t𝐧)}=_t𝐧$$
(95)
and the Hamiltonian density will be given, as usual, by: $`(x,t)=\pi _t𝐧,`$ with canonical (equal-time) Poisson brackets (PB’s):
$$\{n^\alpha (x,t),\pi ^\beta (y,t)\}=\delta ^{\alpha \beta }\delta (xy)$$
(96)
with all the other PB’s vanishing.
Canonical quantization (here too we set: $`\mathrm{}=1`$ and use for simplicity the same notation for the classical variables and the corresponding quantum operators) is accomplished by replacing the nonvanishing PB’s with the commutators:
$$[n^\alpha (x,t),\pi ^\beta (y,t)]=i\delta ^{\alpha \beta }\delta (xy)$$
(97)
Then, we obtain immediately, taking expectation values in the ground state:
$$i_t[n^\alpha (x,t),n^\beta (0,0)]|_{t=0}=i[\pi ^\alpha (x,0),n^\beta (0,0)]=\delta (x)$$
(98)
As the r.h.s. involves only equal-time commutators, this relation will be true irrespective of the form of the interaction term as long as the latter depends only on $`𝐧`$ and not on its time derivative.
Introducing the Fourier transform $`A^{\alpha \beta }(k,\omega )`$ defined via:
$$[n^\alpha (x,t),n^\beta (0,0)]=\frac{dk}{2\pi }\frac{d\omega }{2\pi }\mathrm{exp}\{i[kx\omega t]\}A^{\alpha \beta }(k,\omega )$$
(99)
Eq.$`(98)`$ will imply the sum rule:
$$\frac{d\omega }{2\pi }\omega A^{\alpha \beta }(k,\omega )=\delta ^{\alpha \beta }$$
(100)
Let’s test it in the case: $`V(𝐧)=m^2𝐧^2`$, with $`m`$ fixed, e.g., by the condition: $`𝐧^2(x,t)=1`$. Then the field can be quantized as<sup>2</sup><sup>2</sup>2Note that, in order to reproduce correctly the commutation relations (97) between the fields and the conjugate momenta, the signs of the arguments of the exponentials had to be taken here as the opposites of those of Ref..:
$$n^\alpha (x,t)=\frac{dk}{4\pi \epsilon (k)}\left\{a_k\mathrm{exp}[i(kx\epsilon (k)t)]+a_k^{}\mathrm{exp}[i(kx\epsilon (k)t)]\right\}$$
(101)
($`\epsilon (k)=\sqrt{k^2+m^2}`$), and:
$$[n^\alpha (x,t),\pi ^\beta (y,t)]=i\delta ^{\alpha \beta }\delta (xy)[a_k^\alpha ,a_q^\beta ]=4\pi \delta ^{\alpha \beta }\epsilon (k)\delta (kq)$$
(102)
One finds then easily:
$$[n^\alpha (x,t),n^\beta (0,0)]=\delta ^{\alpha \beta }\frac{dk}{2\pi }\mathrm{exp}(ikx)\frac{1}{2\epsilon (k)}[\mathrm{exp}(i\epsilon (k)t)\mathrm{exp}(i\epsilon (k)t)]$$
(103)
whence:
$$A^{\alpha \beta }(k,\omega )=\delta ^{\alpha \beta }\frac{\pi }{\epsilon (k)}\{\delta (\omega \epsilon (k))\delta (\omega +\epsilon (k))\}$$
(104)
and the sum rule can be checked at once.
To compare with Eq.$`(81)`$, let’s recall first of all that our definition (Eq.$`(17)`$) of the Euclidean propagators differs by a sign from that that is usually adopted . Accordingly, we will define the retarded Green functions of the components of the $`𝐧`$ field as:
$$G_R^{\alpha \beta }(x,t)=i\theta (t)[n^\alpha (x,t),n^\beta (0,0)]$$
(105)
and their Fourier transform are then given by:
$$\stackrel{~}{G}_R^{\alpha \beta }(k,\omega )=\frac{d\omega ^{}}{2\pi }\frac{A^{\alpha \beta }(k,\omega ^{})}{\omega \omega ^{}+i\delta }|_{\delta 0^+}$$
(106)
Hence:
$$\mathrm{Im}\stackrel{~}{G}_R^{\alpha \beta }(k,\omega )=\frac{1}{2}A^{\alpha \beta }(k,\omega ^{})$$
(107)
In view of the standard relationship between retarded (in real time) and causal (in Euclidean time) Green functions, the sum rule that we have just derived extends to the Euclidean propagators introduced in the text and becomes then the sum rule $`(81)`$.
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# Rotations and 𝑒, 𝜈 Propagators, Part I
## 1 Introduction
A free spin 1/2 particle, such as an electron, moves through space-time carrying with it an internal spin. A spin of 1/2 is only one of many choices allowed by space-time symmetry. The theorem that only integer or half integer spins are possible is well known and its agreement with observation is widely accepted.
But isn’t the logic just backwards? The rotation group of the internal spin space mimics a subgroup of the Lorentz group of space-time symmetries. One knows that going from a subgroup to the full group naturally imposes constraints in addition to the requirements of the subgroup. A moment’s reflection confirms this: consider the left-right symmetric objects that can be drawn on a sheet of paper and objects with six-sided symmetry. The left-right symmetric objects can have more variety. (Also see Problem 1.) Hence it is the space-time group, the Lorentz group, that imposes more rules on its elements than the internal spin space’s rotation group. Shouldn’t it be space-time symmetry that is obtained by constraining spin?
In this paper we start with the more general, the more primitive group - the rotation group of a Euclidean space. For completeness, elementary results are collected in Sec. 2 from the well known theory of the $`2\times 2`$ matrix representation of rotations in Euclidean space of three or four dimensions. Two 2-vectors make ‘pairs’ whose properties are the focus of the work. Orthogonality and the normalization of a selected basis of pairs appears in Sec. 3. In Sec. 4 the projection operators are obtained for the basis pairs. The momentum (as-yet-undefined functions) and space occur as parameters of a delta function, which is the part of the projection operator that selects the continuous parameters of a specific basis.
The projection operators have rotational symmetry. To achieve spacetime symmetry, we add two projection matrices and define space, time, energy, and momentum appropriately. An important step in this process, which appears in Sec. 5, is making the space-time 4-vector in the delta function phase be the same as the space-time 4-vector in the two operator summed matrix. The common 4-vector is the energy-momentum with the phase and matrix typically denoted $`p_\mu x^\mu `$ and $`p_\mu \gamma ^\mu ,`$ respectively. Comparison with well-known formulas shows that the final two summed projection operators are just the electron and positron propagators in the absence of interactions.
The essential constraints applied to the rotation based quantities are (i) the specific choice of basis, (ii) the definitions of energy and momentum, and (iii) the choice of projection operators to combine. The calculation follows, with some modification, a classic derivation of propagators by Feynman . Furthermore the hypothesis that positive energy states propagate forward in time and negative energy states propagate backwards in time is explained as irreversible rotation in the Euclidean space.
In Part II we apply the same procedure to a different basis and derive neutrino propagators. In Part III we consider arbitrary bases and show that the $`uv`$ basis here, (9), and the basis considered in Part II are the only bases that allow space time symmetric projection matrices (propagators). Also in Part III is a discussion of the 2-dimensional representations of the rotation group in a Euclidean space of four dimensions and the relationship between planes of rotation and wave functions. The two bases in Parts I and II are special because of a special property of four dimensional Euclidean space.
## 2 Rotation Matrix and 2-Vectors
In this section well known results are recalled for the 2-dimensional representations of the rotation group, which are also known as spinor representations. We consider rotations in a Euclidean space of three or four dimensions. In Part III it is shown that the process and the results imply that four dimensional Euclidean space is more suitable than three dimensional space.
Consider the $`2\times 2`$ matrix $`R`$ with complex-valued elements
$$R=e^{in^k\sigma ^k\theta /2}=\sigma ^4\mathrm{cos}(\theta /2)+in^k\sigma ^k\mathrm{sin}(\theta /2),$$
(1)
where the three $`n^k`$ form a real 3-vector with unit magnitude, summation is assumed over $`k`$ $`\{1,2,3\},`$ $`\theta `$ is real, and the sigma matrices are
$$\sigma ^1=\left(\begin{array}{ccc}0& & 1\\ 1& & 0\end{array}\right)\sigma ^2=\left(\begin{array}{ccc}0& & i\\ i& & 0\end{array}\right)\sigma ^3=\left(\begin{array}{ccc}1& & 0\\ 0& & 1\end{array}\right)\sigma ^4=\left(\begin{array}{ccc}1& & 0\\ 0& & 1\end{array}\right).$$
(2)
The set of all matrices like $`R`$ forms a group under matrix multiplication. For small $`\theta `$ the matrix group is isomorphic to the rotation group in three dimensions with $`R`$ representing the rotation about the origin through an angle $`\theta `$ that preserves the axis in the direction of the unit 3-vector $`n^k`$ in a 3-dimensional Euclidean space, $`E_{3\mathrm{d}}`$. In a 4-dimensional Euclidean space the same $`R`$ can represent a rotation in another plane that intersects the first only at the origin.
The $`2\times 2`$ matrix $`R`$ acts on ordered sets of two complex numbers, ‘2-vectors’. One basis for 2-vectors consists of the two unit eigenvectors of $`R`$
$$u^+=e^{i\alpha }\left(\begin{array}{c}\mathrm{cos}(\rho /2)\mathrm{exp}(i\varphi /2)\\ \mathrm{sin}(\rho /2)\mathrm{exp}(+i\varphi /2)\end{array}\right)u^{}=e^{i\beta }\left(\begin{array}{c}\mathrm{sin}(\rho /2)\mathrm{exp}(i\varphi /2)\\ \mathrm{cos}(\rho /2)\mathrm{exp}(+i\varphi /2)\end{array}\right),$$
(3)
where $`\rho `$ and $`\varphi `$ are the polar and azimuthal angles for the unit 3-vector for the rotation plane
$$n^k=(\mathrm{sin}\rho \mathrm{cos}\varphi ,\mathrm{sin}\rho \mathrm{sin}\varphi ,\mathrm{cos}\rho )$$
(4)
and $`\alpha `$ and $`\beta `$ are arbitrary phases. By (1), (2), (3), and (4), we have
$$Ru^+=e^{+i\theta /2}u^+Ru^{}=e^{i\theta /2}u^{}.$$
(5)
Call $`u^+`$ ‘spin up’ and $`u^{}`$ ‘spin down.’ By (2), (3), and (4), we get
$$n^k\sigma ^ku^+=+u^+n^k\sigma ^ku^{}=u^{}u^\pm \sigma ^ku^\pm =\pm n^ku^\pm \sigma ^4u^\pm =1.$$
(6)
By (3) $`u^+`$ and $`u^{}`$ are orthonormal,
$$u^+u^+=1u^{}u^{}=1u^+u^{}=0u^{}u^+=0.$$
(7)
Hence any 2-vector $`u`$ can be written as $`u`$ = $`au^++bu^{}`$ where $`a`$ = $`u^+u`$ and $`b`$ = $`u^{}u.`$ The following matrix combinations prove useful,
$$u^+u^++u^{}u^{}=\sigma ^4u^+u^+u^{}u^{}=n^k\sigma ^k.$$
(8)
## 3 Pairs of 2-vectors
In this section, we seek to use the properties of pairs of 2-vectors to obtain energy, momentum and polarization matrices. The immediate justification for using pairs of 2-vectors is simply that four component wave functions are standard and they work. The properties of four dimensional Euclidean space are seen to justify such wave functions in Part III.
Consider the set of all ordered pairs of 2-vectors, hereafter simply called ‘pairs.’ Thus each pair has four complex components. One basis is
$$u_+^+=\left(\begin{array}{c}e^{w/2}u^+\\ e^{w/2}u^+\end{array}\right)u_{}^{}=\left(\begin{array}{c}e^{w/2}u^{}\\ e^{w/2}u^{}\end{array}\right),$$
(9)
$$v_+^+=\left(\begin{array}{c}e^{w/2}u^+\\ e^{w/2}u^+\end{array}\right)v_{}^{}=\left(\begin{array}{c}e^{w/2}u^{}\\ e^{w/2}u^{}\end{array}\right).$$
This basis is special because in each pair the upper and lower 2-vectors are eigenvectors of $`R`$ with the same eigenvalue.
By (7) and (9), the basis pairs are normalized to $`2\mathrm{cosh}(w)`$ and are mutually orthogonal,
$$i^{}i=2\mathrm{cosh}wi^{}j=0,i,j\{u_+^+,u_{}^{},v_+^+,v_{}^{}\},$$
(10)
where there is no sum over $`i`$ in the left expression and $`i`$ is not the same as $`j`$ in the middle expression.
Let $`R_+^+u_+^+`$ indicate applying the $`2\times 2`$ matrix $`R`$ to the upper and lower 2-vectors of $`u_+^+.`$ By (5), we get a rotated basis with a common phase factor,
$$R_+^+u_+^+=e^{+i\theta /2}u_+^+R_{}^{}u_{}^{}=e^{+i\theta /2}u_{}^{}R_+^+v_+^+=e^{+i\theta /2}v_+^+R_{}^{}v_{}^{}=e^{+i\theta /2}v_{}^{},$$
(11)
where $`R_{}^{}`$ means applying $`R^1`$ to both the upper 2-vector and the lower 2-vector. The rotated basis (11) satisfies the orthogonality and normalization conditions (10).
There are other bases. Given $`u_+^+`$ and requiring an orthogonal basis, we must have $`v_+^+.`$ But $`u_{}^{}`$ could have the same ratio $`e^w`$ of upper to lower 2-vector as does $`u_+^+;`$ just change $`w`$ $`w`$ in $`u_{}^{}`$ and $`v_{}^{}.`$ The basis in (9) is almost uniquely determined by requiring orthogonality and requiring that the basis gives the matrix expressions which display space-time symmetry. (See Part III.)
## 4 Projection Operators for Pairs of 2-vectors
A projection operator replicates selected parts of a quantity and destroys the rest. Given a linear combination of the rotated basis pairs, one can project out the coefficient of any one of the four basis pairs. The projection operator must select the ratio parameter $`w`$, the rotation plane $`n^k`$, the eigenvector $`u^+`$ or $`u^{}`$, and the rotation angle $`\theta .`$
Select the $`u_+^+`$ term from a sum of basis pairs. An arbitrary 4-component object $`\psi ,`$ can be written as a linear combination of the basis pairs (9),
$$\psi au_+^++bu_{}^{}+cv_+^++dv_{}^{},$$
(12)
where $`a`$ = $`u_{+}^{+}{}_{}{}^{}\psi /(2\mathrm{cosh}w).`$ To make the transition to space-time symmetry below it is convenient to include a factor of $`\gamma ^4.`$ See (46) for the definition of $`\gamma ^4.`$ We have the projection matrix
$$K\gamma ^4\frac{1}{2\mathrm{cosh}w}u_+^+u_+^+K\gamma ^4\psi =au_+^+$$
(13)
The projection matrices for the other base pairs $`u_{}^{},`$ $`v_+^+,`$ and $`v_{}^{}`$ are similar.
Select the ratio parameter $`w`$ and axis $`n^k`$. By (3), (4), and (9), the basis pairs are functions of the ratio parameter $`w`$ and the unit 3-vector $`n^k.`$ Since $`n_{}^{k}{}_{}{}^{2}`$ = 1, in the set of four variables $`w`$ and $`n^k`$ only three are independent and we can use a three dimensional delta function to select specific values of these variables. However, it is not conventional to write the delta functions of $`w`$ and $`n^k,`$ instead the delta functions are written in terms of the ‘momentum’ $`p^j,`$ which are three functions of $`w`$ and $`n^k,`$ $`p^j(w,n^k).`$ This means that $`w`$ and $`n^k`$ are determined by $`p^k,`$ i.e. $`w(p^j)`$ and $`n^k(p^j),`$ and we can use delta functions for the three momentum components, $`p^jp_{}^{j}{}_{}{}^{},`$ instead of delta functions for $`ww^{}`$ and $`n^kn_{}^{k}{}_{}{}^{}.`$ The momentum function $`p^j`$ remains undetermined for now until (33) and (41) in Sec. 5 where the requirements of space-time symmetry are considered.
Any function of $`p^j,`$ such as the basis pair $`u_+^+`$ = $`u_+^+(p^j),`$ can be made the amplitude of a plane wave,
$$\psi _p(1)e^{ip^kx_1^k}u_+^+(p^j).$$
(14)
Call the phase $`p^kx^k`$ the ‘plane wave phase.’ By integrating with the three dimensional delta function, we get
$$d^3x_1K_p^{}(2,1)\psi _p(1)=e^{ip^kx_2^k}u_+^+(p^j)=\psi _p(2),$$
(15)
where
$$K_p^{}(2,1)\frac{d^3p^{}}{(2\pi )^3}e^{ip^kx_2^k}e^{ip^kx_1^k}.$$
(16)
By (15) the delta function has no effect on the plane wave except to reparameratize the phase, i.e. $`x_1^k`$ $`x_2^k.`$
Rotation angle $`\theta `$. Rotation of $`u_+^+`$ or $`v_+^+`$ by $`R`$ through an arbitrary angle $`\theta `$ produces a phase factor $`\psi _\theta `$ = $`e^{i\theta /2}`$ by (11). The arbitrary phase can be replicated by multiplication with another phase factor $`K_\theta (2,1),`$
$$\psi _\theta (1)e^{i\theta _1/2}K_\theta (2,1)e^{i\theta _2/2}e^{i\theta _1^{}/2}$$
(17)
$$K_\theta (2,1)\psi _\theta (1)=(e^{i\theta _2/2}e^{i\theta _1^{}/2})e^{i\theta _1/2}=e^{i(\theta _2\theta _1^{}+\theta _1)}=e^{i\theta _2/2}=\psi _\theta (2).$$
(18)
There is no loss of generality in choosing $`\theta _1`$ = $`\theta _1^{}`$ because $`\theta _2`$ can have any real value and therefore represent any angle in the new function $`\psi _\theta (2)`$ = $`e^{i\theta _2/2}`$, just as $`\theta _1`$ can be any angle in the original $`\psi _\theta (1).`$
The plane wave phase $`p^kx^k`$ differs from the eigenvalue phase $`\theta /2.`$ Let $`\mathrm{\Delta }`$ be the difference so that
$$\frac{\theta }{2}=\mathrm{\Delta }+p^kx^k.$$
(19)
Thus we can write $`\psi _\theta (1)`$ and $`K_\theta `$ as
$$\psi _\theta (1)=e^{i\mathrm{\Delta }_1}e^{ip^kx_1^k}K_\theta (2,1)=e^{i\mathrm{\Delta }_2^{}}e^{ip^kx_2^k}e^{i\mathrm{\Delta }_1^{}}e^{ip^kx_1^k}.$$
(20)
Since the rotation angle $`\theta `$ can have any real value, by (19) the difference $`\mathrm{\Delta }`$ can have any real value at each point $`x^k.`$
Projection Operators. Now we project the positive or negative ratio, spin up or spin down portion out of an arbitrary 4-component wave function $`\psi `$ by combining the various $`4\times 4`$ projection matrices, delta functions, and phase adjustments.
Expand the wave function $`\psi (1)`$ over the basis pairs (11).
$$\psi (1)=aR_+^+u_+^++bR_{}^{}u_{}^{}+cR_+^+v_+^++dR_{}^{}v_{}^{}=$$
(21)
$$=e^{i\theta _1/2}(au_+^++bu_{}^{}+cv_+^++dv_{}^{})=e^{i\theta _1/2}\psi _0$$
where $`a`$ = $`(R_+^+u_+^+)^{}\psi /(2\mathrm{cosh}w)`$ and $`\psi _0`$ is $`\psi (1)`$ for $`\theta _1`$ = 0. The $`R_+^+u_+^+`$ projection operator combines the $`K`$s from (13), (17), and (20). We have
$$K(2,1,R_+^+u_+^+)\gamma ^4\frac{d^3p^{}}{(2\pi )^3}\frac{1}{2\mathrm{cosh}(w)^{}}e^{i\theta _2^{}}u_+^+u_{+}^{+}{}_{}{}^{}e^{i\theta _1^{}/2},$$
(22)
where $`\theta _i^{}`$ = $`\mathrm{\Delta }_i^{}+p^kx_i^k.`$ By (13), (15), and (18), we get
$$d^3x_1K(2,1,R_+^+u_+^+)\gamma ^4\psi (1)=ae^{i\theta _2/2}u_+^+=aR_+^+u_+^+=\psi (2)_{b=c=d=0},$$
(23)
where $`\theta _2/2`$ = $`\mathrm{\Delta }_2^{}+\mathrm{\Delta }_1^{}\mathrm{\Delta }_1+p^kx_2^k`$ = $`\mathrm{\Delta }_2+p^kx_2^k,`$ i.e. we make $`\mathrm{\Delta }_2`$ obey $`\mathrm{\Delta }_2\mathrm{\Delta }_1`$ = $`\mathrm{\Delta }_2^{}\mathrm{\Delta }_1^{}.`$ Equation (23) illustrates a projection operator at work: $`K(2,1,R_+^+u_+^+)\gamma ^4`$ projects out the $`R_+^+u_+^+`$ portion of $`\psi (1).`$ The projection operators for the other basis pairs, $`K(2,1,R_{}^{}u_{}^{})\gamma ^4,`$ $`K(2,1,R_+^+v_+^+)\gamma ^4,`$ $`K(2,1,R_{}^{}v_{}^{})\gamma ^4`$ are given by the expression (22) with the basis pair $`R_+^+u_+^+`$ replaced by $`R_{}^{}u_{}^{},`$ $`R_+^+v_+^+,`$ and $`R_{}^{}v_{}^{},`$ respectively.
The projection operator $`K(2,1,R_+^+u_+^+)`$ in (22) is invariant under rotations when both $`p^k`$ and $`x^k`$ transform as vectors leaving $`p^kx^k`$ and $`\theta `$ invariant. To show that $`u_+^+u_{+}^{+}{}_{}{}^{}`$ is a rotation invariant, see (30) below. Making spacetime invariant projection operators using the four rotation invariant projection operators, like $`K(2,1,R_+^+u_+^+),`$ is a task for the next section.
## 5 Space-Time Symmetry, the Electron Propagator
In this section two projection operators are obtained and written in a way that makes their invariance under space-time transformations evident. For phases, we show how the quantity $`\mathrm{\Delta }`$ in (19) becomes $`Et,`$ where $`E`$ is the ‘energy’ and $`t`$ is time. We see what it means to require that positive energy states must propagate forward in time and negative energy states must propagate backwards in time, i.e. the positron hypothesis is explained. For matrices, we find that we can combine the two projection operators for the positive ratio pairs $`u_+^+`$ and $`u_{}^{},`$ thereby making a positive energy matrix. Then the two negative ratio operators combine to make a negative energy matrix.
Time. Time can be introduced by applying an integral expression for the phase factor $`e^{i\mathrm{\Delta }}`$ ,
$$e^{i\mathrm{\Delta }}=\frac{i}{\pi }_{\mathrm{}}^{\mathrm{}}𝑑a\frac{e^{ia\mathrm{\Delta }}}{a^21+iϵ}$$
(24)
where the value of $`\mathrm{\Delta }`$ is positive and $`ϵ`$ is small and positive so that the pole at $`a`$ = 1 is included when the integral is evaluated by contour integration. For brevity, we do not consider $`\mathrm{\Delta }<`$ 0 except to note that, for $`\mathrm{\Delta }_2^{}<0`$, change the overall sign and change the sign of $`ϵ`$ in (24).
By (19), (22), and (23), we have
$$\psi (2)_{b=c=d=0}=d^3x_1K(2,1,R_+^+u_+^+)\gamma ^4\psi (1)$$
$$=d^3x_1\frac{d^3p^{}}{(2\pi )^3}\frac{1}{2\mathrm{cosh}(w)^{}}e^{i\theta _2^{}/2}u_+^+u_{+}^{+}{}_{}{}^{}e^{i\theta _1^{}/2}e^{i\theta _1/2}\gamma ^4\gamma ^4\psi _0$$
$$=d^3x_1[\frac{d^3p^{}}{(2\pi )^3}\frac{1}{2\mathrm{cosh}w^{}}e^{i\theta _2^{}/2}u_+^+u_{+}^{+}{}_{}{}^{}\gamma ^4]e^{i(\mathrm{\Delta }_1^{}\mathrm{\Delta }_1)}e^{i(p^kp^k)x_1^k}\gamma ^4\psi _0.$$
(25)
Let $`I`$ be the quantity in (25) in brackets. By (19) and (24) and for $`\mathrm{\Delta }_2^{}>`$ 0, we get
$$I=\frac{d^3p^{}}{(2\pi )^3}\frac{1}{2\mathrm{cosh}w^{}}e^{i\mathrm{\Delta }_2^{}}e^{ip^kx_2^k}u_+^+u_{+}^{+}{}_{}{}^{}\gamma ^4$$
$$=\frac{d^3p^{}}{(2\pi )^3}\frac{1}{2\mathrm{cosh}w^{}}(\frac{i}{\pi }_{\mathrm{}}^{\mathrm{}}𝑑a\frac{e^{ia\mathrm{\Delta }_2^{}}}{a^21+iϵ})e^{ip^kx_2^k}u_+^+u_{+}^{+}{}_{}{}^{}\gamma ^4$$
$$=i\frac{d^4p^{}}{(2\pi )^4}e^{i(p^4t_2p^kx_2^k)}\frac{(\pm m)u_+^+u_{+}^{+}{}_{}{}^{}\gamma ^4}{p_{}^{4}{}_{}{}^{2}m^2\mathrm{cosh}^2w^{}+iϵ}$$
(26)
where we introduce a fourth component of momentum and a fourth component of displacement, i.e. energy and time,
$$p_{}^{4}{}_{}{}^{}\pm m\mathrm{cosh}(w)^{}aa\mathrm{\Delta }_2^{}=p_{}^{4}{}_{}{}^{}t_2,$$
(27)
where $`m`$ is a positive constant tentatively called the ‘bare mass.’
Surface Integral. Next, in (25) identify the integral over the 3-space $`x_1^k`$ with a surface integral in space-time,
$$\psi (2)_{b=c=d=0}=_Sd^4x_1Ie^{i(\mathrm{\Delta }_1^{}\mathrm{\Delta }_1)}e^{i(p^kp^k)x_1^k}N_\mu \gamma ^\mu \psi _0$$
(28)
where $`x^4`$ = $`t,`$ $`\mu `$ $`\{1,2,3,4\},`$ $`N^\mu `$ = $`\{N^k,N^4\}`$ is the unit normal to the three dimensional surface of integration $`S`$ in four dimensional space-time, and the space-time summation convention is used, $`N_\mu \gamma ^\mu `$ = $`N^4\gamma ^4N^k\gamma ^k.`$ In (23) and (25), the integration is over the three dimensional surface $`x_1^k`$ in the special space-time reference frame with the normal in the time direction, $`N^\mu `$ = $`\{0,0,0,1\}.`$
Space-time Transformation. The linear transformations of $`p^\mu `$ and $`x^\mu `$ that preserve the sum $`p_\mu x^\mu `$ = $`p^4tp^kx^k`$ are the space-time transformations, i.e. Lorentz transformations. Only homogeneous transformations are considered here. We denote the coefficients of one such transformation by $`\mathrm{\Lambda }_\nu ^\mu ,`$
$$P^\mu =\mathrm{\Lambda }_\nu ^\mu p^\nu P_\mu P^\mu =P_{}^{\mathrm{\hspace{0.17em}4}}{}_{}{}^{2}P_{}^{k}{}_{}{}^{2}=p_{}^{4}{}_{}{}^{2}p_{}^{k}{}_{}{}^{2}.$$
(29)
Thus the phases, and hence also the rotation angle $`\theta ,`$ in (26) are space-time invariants because they contain scalar products such as $`p_{\mu }^{}{}_{}{}^{}x_{}^{\mu }{}_{}{}^{}.`$
Matrices. The projection operator $`K(2,1,R_+^+u_+^+)\gamma ^4`$ rewritten with (25) and (26) still fails to have space-time symmetry because of the quantity $`mu_+^+u_{+}^{+}{}_{}{}^{}\gamma ^4`$ in (26). We can express the matrix $`mu_+^+u_{+}^{+}{}_{}{}^{}\gamma ^4`$ as a sum of sixteen linearly independent $`4\times 4`$ matrices such as the set of gamma matrices in Appendix A. By (8), (9), and (49), we get
$$mu_+^+u_{+}^{+}{}_{}{}^{}\gamma ^4=\frac{1}{2}[m1m\mathrm{sinh}(w)n^i\gamma ^i+m\mathrm{cosh}(w)\gamma ^4]+$$
(30)
$$+\frac{i}{2}[m\mathrm{cosh}(w)n^j\gamma ^j\gamma ^5m\mathrm{sinh}(w)\gamma ^4\gamma ^5+mn^k\gamma ^5\gamma ^4\gamma ^k],$$
where we drop the primes.
Let $`n^k`$ and $`\gamma ^k`$ transform as 3-vectors under rotations. One can show that $`\gamma ^5`$, (46), is invariant under rotations. Thus the matrix (30) is a rotation invariant. While the matrix (30) is a rotation invariant, it is not a space-time invariant; see Problem 3.
The expansion of $`mu_{}^{}u_{}^{}{}_{}{}^{}\gamma ^4`$ over the same set of gammas gives
$$mu_{}^{}u_{}^{}{}_{}{}^{}\gamma ^4=\frac{1}{2}[m1m\mathrm{sinh}(w)n^i\gamma ^i+m\mathrm{cosh}(w)\gamma ^4]+$$
(31)
$$\frac{i}{2}[m\mathrm{cosh}(w)n^j\gamma ^j\gamma ^5m\mathrm{sinh}(w)\gamma ^4\gamma ^5+mn^k\gamma ^5\gamma ^4\gamma ^k].$$
When we add the two expressions, we get
$$m(u_+^+u_{+}^{+}{}_{}{}^{}+u_{}^{}u_{}^{}{}_{}{}^{})\gamma ^4=m\mathrm{cosh}(w)\gamma ^4m\mathrm{sinh}(w)n^k\gamma ^k+m1.$$
(32)
Let $`\gamma ^k`$ and $`\gamma ^4`$ form a 4-vector, i.e. transform as in (29). See for why this can be allowed. Let $`m`$ and the unit matrix $`1`$ be scalars under space-time transformations. If $`m\mathrm{sinh}(w)n^k`$ and $`m\mathrm{cosh}w`$ also transform as the components of a 4-vector, then the expression (32) is a space-time invariant.
Energy-momentum. We now make the choice of momentum functions and, at the same time, we choose the energy function. Note that both $`\{m\mathrm{sinh}(w)n^k,m\mathrm{cosh}w\}`$ and $`\{p^k,p^4\}`$ = $`\{p^k,\pm m\mathrm{cosh}(w)a\}`$ are 4-vectors. We define the momentum functions
$$p^km\mathrm{sinh}(w)n^kp^4=+m\mathrm{cosh}(w)a.$$
(33)
Then the pole in (24) occurs when $`p^4`$ has the value $`E,`$
$$Em\mathrm{cosh}w\mathrm{\Delta }=Et.$$
(34)
The energy-momentum components connect the invariant phase $`\theta /2`$ = $`Et`$ $`p^kx^k`$ and the matrix (32), which we can now rewrite as
$$m(u_+^+u_{+}^{+}{}_{}{}^{}+u_{}^{}u_{}^{}{}_{}{}^{})\gamma ^4=E\gamma ^4p^k\gamma ^k+m1.$$
(35)
The phases $`\mathrm{\Delta }_1^{}\mathrm{\Delta }_1`$ and $`(p^kp^k)x_1^k`$ in (25) can be treated similarly. Since the propagated wave function $`\psi (2)`$ is the initial wave function for the next propagator, the phase for $`\psi (1)`$ should have the same form as $`\psi (2).`$ It follows that, as in (27), we can write $`\mathrm{\Delta }_1^{}`$ = $`E^{}t_1`$ and $`\mathrm{\Delta }_1`$ = $`Et_1`$ so that the integral over $`t_1`$ = $`x_1^4`$ makes a delta function $`\delta (E^{}E).`$ In the special reference frame of Sec. 4, the basis pairs are functions of $`p^k,`$ i.e. $`w`$ and $`n^k,`$ so the delta function for time does not constrain the basis pairs. But in a more general frame the four dimensional delta function is needed to select the basis pairs corresponding to the $`w`$ and $`n^k`$ of the special frame.
$`Em`$ Propagator. The space-time symmetry of the sum of the two positive ratio operators, $`K(2,1,R_+^+u_+^+)\gamma ^4`$ and $`K(2,1,R_{}^{}u_{}^{})\gamma ^4,`$ is evident when the above results are combined. By (25), (26), (28), (33), (34), and (35), we get
$$\psi (2)_{c=d=0}=d^3x_1(K(2,1,R_+^+u_+^+)+K(2,1,R_{}^{}u_{}^{}))\gamma ^4\psi (1)=$$
(36)
$$=i\frac{d^4x_1d^4p^{}}{(2\pi )^4}e^{ip_\mu ^{}x_2^\mu }\frac{p_\nu ^{}\gamma ^\nu +m1}{p_\eta ^{}p^\eta m^2+iϵ}e^{ip_\tau ^{}x_1^\tau }N_\sigma \gamma ^\sigma \psi (1),$$
where we used another integral formula,
$$e^{i\mathrm{\Delta }}=\frac{i}{\pi }_{\mathrm{}}^{\mathrm{}}𝑑a\frac{ae^{ia\mathrm{\Delta }}}{a^21+iϵ},$$
(37)
to replace $`E`$ $`Ea`$ = $`p^4`$ in the matrix $`(E\gamma ^4p^k\gamma ^k+m1).`$
Spacetime symmetry is evident in (36) because scalar products like $`p_\mu p^\mu `$ are invariant under space-time transformations and the integrals are over four dimensions.
Experimental confirmation. The fraction in (36) is sometimes written $`(p_\mu ^{}\gamma ^\mu m)^1`$ because
$$(p_\mu ^{}\gamma ^\mu m)(p_\nu ^{}\gamma ^\nu +m)=(p_\sigma ^{}p_{}^{\sigma }{}_{}{}^{}m^2).$$
(38)
The projection operator (36) can be shown to be the Green’s function for the inverted matrix operator $`(p_\mu ^{}\gamma ^\mu m).`$ And one can show that the Green’s function is related to the Dirac equation. This justifies the use of the terms ‘energy’ and ‘momentum’ for $`E`$ and $`p^k`$ because the formulas have been used to describe electrons and other spin 1/2 particles. In particular, expression (36) is equivalent to the standard electron propagator for positive energy states.
$`Em`$ Propagator. For the basis pairs $`v_+^+`$ and $`v_{}^{}`$ we can follow the steps that lead to (32). We have
$$m(v_+^+v_{+}^{+}{}_{}{}^{}+v_{}^{}v_{}^{}{}_{}{}^{})\gamma ^4=(m1m\mathrm{sinh}(w)n^k\gamma ^km\mathrm{cosh}(w)\gamma ^4).$$
(39)
Then, for $`\mathrm{\Delta }_2^{}>`$ 0, we get
$$\psi (2)_{a=b=0}=d^3x_1[K(2,1,R_+^+v_+^+)+K(2,1,R_{}^{}v_{}^{})]\gamma ^4\psi (1)=$$
$$=d^3x_1\frac{d^3p^{}}{(2\pi )^3}\frac{1}{2m\mathrm{cosh}w^{}}e^{i\mathrm{\Delta }_2^{}}e^{ip^kx_2^k}(m\mathrm{cosh}(w^{})\gamma ^4m\mathrm{sinh}(w^{})n^k\gamma ^k+m)\gamma ^4\psi (1)$$
$$=i\frac{d^4x_1d^4p^{}}{(2\pi )^4}e^{ip_\mu ^{}x_2^\mu }\frac{p_\nu ^{}\gamma ^\nu +m}{p_\tau ^{}p^\tau m+iϵ}e^{ip_\eta ^{}x_1^\eta }N_\sigma \gamma ^\sigma \psi (1)$$
(40)
Since the $`\mathrm{cosh}w^{}`$ term is negative in the matrix factor $`(m\mathrm{cosh}(w^{})\gamma ^4m\mathrm{sinh}(w^{})n^k\gamma ^k+m),`$ we need to choose the negative sign in (27) and we get new definitions for $`p_{}^{4}{}_{}{}^{}`$ and $`t_2`$ in (40). We define new momentum and energy functions
$$p^k=m\mathrm{sinh}(w)n^kp^4=m\mathrm{cosh}(w)a.$$
(41)
Note that we must choose the negative sign in the $`\pm m`$ term in (26) and (27). The pole in (24) occurs when $`p^4`$ has the value $`E,`$
$$Em\mathrm{cosh}w\mathrm{\Delta }=Et.$$
(42)
Positron Hypothesis. We now consider the positron hypothesis, i.e. positive energy states propagate forward in time while negative energy states propagate backwards in time. By (19), (33), and (34), we get an expression for time $`t`$ for positive energy states
$$t=\frac{1}{\mathrm{cosh}(w)}(\mathrm{sinh}(w)n^kx^k\frac{\theta }{2m})(Em>0).$$
(43)
And, by (19), (41), and (42), we get an expression for time $`t`$ when the energy is negative,
$$t=\frac{1}{\mathrm{cosh}(w)}(\frac{\theta }{2m}\mathrm{sinh}(w)n^kx^k)(Em<0).$$
(44)
In the special reference frame of Sec. 4 let the time $`t_1`$ be zero on the three dimensional surface of integration $`x_1^k.`$ Identify $`x_2^k`$ with $`x_1^k`$ for $`t_2`$ = 0. By (43), if $`\theta `$ decreases then $`t_2`$ is positive everywhere, $`t_2>`$ 0 for positive energy. And, and by (44), if $`\theta `$ decreases then $`t_2`$ is negative everywhere, $`t_2>`$ 0 for negative energy. We have, at any point $`x_2^k,`$
$$\mathrm{positive}\mathrm{energy}:Em\mathrm{and}\theta \mathrm{decreasing}t_2>0\mathrm{and}t_2\mathrm{increasing}$$
$$\mathrm{negative}\mathrm{energy}:Em\mathrm{and}\theta \mathrm{decreasing}t_2<0\mathrm{and}t_2\mathrm{decreasing}.$$
(45)
Thus the positron hypothesis is equivalent to requiring that the rotation angle $`\theta `$ decreases with time, i.e. the rotation angle always decreases as the absolute value of the time increases.
Remarks. One should note that the functions chosen for $`p^k`$ and $`E`$ in (33), (34), (41), and (42) have immediate impact on the properties of the space and time variables $`x^\mu `$ with which they appear in the delta function phase.
Since the $`uv`$ basis (9) and $`\psi `$ are originally written as functions of $`w`$ and $`n^k`$ and the delta functions are in terms of $`p^\mu ,`$ the two different definitions of energy (34) and (42) require the $`uv`$ basis and $`\psi `$ to be two different functions of $`p^4.`$
In this part, we have obtained the standard QED propagator for a free spin 1/2 particle moving forward in time, (36), and the QED propagator for a free spin 1/2 particle moving backward in time, (40). The propagators are expressed in terms of quantities defined on a rotation group in a Euclidean space.
## Appendix A Sixteen Gammas
The set of $`4\times 4`$ matrices is spanned by a basis consisting of sixteen linearly independent $`4\times 4`$ matrices, which can be chosen in many ways. We choose a convenient set of matrices, called the ‘gammas’, to express the results in the text simply.
Start by defining five gammas in terms of the sigmas (2),
$$\gamma ^k=\left(\begin{array}{ccc}0& & \sigma ^k\\ \sigma ^k& & 0\end{array}\right)\gamma ^4=\left(\begin{array}{ccc}0& & \sigma ^4\\ \sigma ^4& & 0\end{array}\right)\gamma ^5=\gamma ^4\gamma ^1\gamma ^2\gamma ^3,$$
(46)
where $`k`$ $`\{1,2,3\}.`$ Then the set of sixteen gammas is
$$\gamma _A=\{1;\gamma ^k;\gamma ^4;\gamma ^k\gamma ^4;\gamma ^5\gamma ^4\gamma ^k;\gamma ^4\gamma ^5;\gamma ^k\gamma ^5;\gamma ^5\},$$
(47)
where $`1`$ is the unit matrix and $`A`$ takes successive integer values starting with $`\gamma _0`$ = $`1`$ to $`\gamma _{15}`$ = $`\gamma ^5.`$
The square of each matrix (47) is either plus or minus the unit matrix and the trace of each matrix is zero except for $`\gamma _0.`$ We have
$$\gamma _{A}^{}{}_{}{}^{2}=ϵ^A1\mathrm{trace}(\gamma _0)=\mathrm{trace}(1)=4\mathrm{trace}(\gamma _A)_{A0}=0,$$
(48)
where $`ϵ^A`$ = $`\pm 1.`$ We can use these properties to obtain the coefficients $`\alpha ^A`$ in the expansion of a $`4\times 4`$ matrix $`M`$. We get
$$M=\alpha ^A\gamma ^A\alpha ^A=\frac{ϵ^A}{4}\mathrm{trace}(\gamma _A.M),$$
(49)
where $`A`$ is summed in the left expression but $`A`$ is not summed in the expression on the right.
## Appendix B Problems
1. Consider the set $`S`$ of all spheres centered on the origin of a given inertial frame $`F.`$ Each sphere is invariant under rotations about the origin in $`F.`$ What subset $`s`$ of $`S`$ contains spheres that are also spheres in the sets $`S^{}`$ of spheres in all other space-time reference frames $`F^{}\mathrm{?}`$ \[Hence space-time symmetry is more restrictive than rotational symmetry.\]
2. Verify (24).
3. (a) Derive (30). (b) Show that $`mu_+^+u_{+}^{+}{}_{}{}^{}\gamma ^4`$ is invariant under rotations and also under boosts in the direction of $`p^k,`$ but not invariant under boosts perpendicular to $`p^k.`$ Assume that $`\gamma ^\mu `$ transforms as a 4-vector, i.e. like $`p^\mu `$ in (29).
4. The proper time $`\tau `$ is the ordinary time $`t`$ in a space-time reference frame with $`p^k`$ = 0. How does the proper time along a path $`x^\mu `$ that is proportional to $`p^\mu `$ depend on the rotation angle $`\theta \mathrm{?}`$
5. Let $`A`$ = $`(u_+^+u_{+}^{+}{}_{}{}^{}+u_{}^{}u_{}^{}{}_{}{}^{})/(2\mathrm{cosh}w)`$ and $`B`$ = $`(v_+^+v_{+}^{+}{}_{}{}^{}+v_{}^{}v_{}^{}{}_{}{}^{})/(2\mathrm{cosh}w).`$ (a) Show that $`A`$ and $`B`$ are idempotent and divisors of zero, i.e. $`A^2`$ = $`A,`$ $`B^2`$ = $`B,`$ and $`AB`$ = 0. Also show that $`A+B`$ = 1. (b) Then use (33), (34), and (35) and (39), (41), and (42) to rewrite $`A`$ and $`B`$ in terms of $`E,`$ $`p^k,`$ and $`\gamma ^\mu .`$
| Glossary | |
| --- | --- |
| Euclidean Space Rotations (Internal Spin Space) | |
| $`R`$ | $`2\times 2`$ rotation matrix |
| $`\theta `$ | rotation angle |
| $`n^k`$ | unit 3-vector |
| $`u^+,`$ $`u^{}`$ | eigenvectors of $`R`$ with eigenvalues $`e^{+i\theta /2},`$ $`e^{i\theta /2}`$ |
| $`u_+^+,`$ $`u_{}^{},`$ $`v_+^+,`$ $`v_{}^{}`$ | pairs of eigenvectors; a basis for the set of all pairs |
| $`e^{+w},`$ $`e^w`$ | ratio of upper to lower 2-vector in a pair |
| Space-time Quantities | |
| $`m`$ | mass is not defined in terms of rotation quantities |
| $`E`$ = $`m\mathrm{cosh}w`$ | energy |
| $`p^k`$ = $`m\mathrm{sinh}(w)n^k`$ | momentum |
| $`x^k`$ | space occurs as parameters in a 3d delta function |
| | that selects $`w`$ and $`n^k`$ |
| $`\tau `$ = $`\theta /(2m)`$ | proper time |
| $`Et`$ = $`p^kx^km\theta /2`$ | time arises from the difference between |
| | delta function phase and eigenphase |
| $`\psi `$ | wave function for the state with definite |
| | momentum and energy, i.e. definite $`w`$ and $`n^k`$ |
| $`K\gamma ^4,`$ $`K_p^{},`$ $`K_\theta ,`$ $`K(2,1,R_+^+u_+^+)\gamma ^4`$ | parts of projection operators used |
| | to select and replicate various parts of $`\psi `$ |
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# The W4 chimney/superbubble
## 1. The original observations: morphology and energy source
The pilot project of the Canadian Galactic Plane Survey (CGPS) brought to light a conical void in the HI above the W4 HII region (Normandeau, Taylor, & Dewdney 1996; hereafter NTD96). This cone is perpendicular to the plane and opens up towards higher Galactic latitudes. At its apex is the OCl 352 cluster containing nine O-stars, two of which are of very early type. Given this, it is unlikely that there has been a supernova in the cluster.
Within the lower density region there are HI “streamers” and two elongated molecular clouds, all of which point away from OCl 352 (Heyer et al. 1996 and Taylor et al. 1998). These features combine to suggest that this is a Galactic chimney blown by the stellar winds of the O-stars, with gas streaming upward toward the halo.
## 2. Is it really a chimney?
### 2.1. H$`\alpha `$ imaging
Narrow band H$`\alpha `$ observations by Dennison, Topasna, & Simonetti (1997; hereafter DTS97) suggest the presence of a faint cap at $`b7^{}`$, corresponding to a height of approximately 200 pc above the star cluster. Their data, scaled by $`\mathrm{sin}(b)`$ to highlight the faint higher latitude emission, are shown in the left-hand panel of Figure 1.
### 2.2. Modelling
In NTD96, the structure’s approximate age was derived from the Weaver et al. (1977) formalism for a bubble expanding in a uniform medium. However, the evolution of the superbubble must be affected by the general decrease in density with increasing latitude. Basu, Johnstone, & Martin (1999; hereafter BJM99) modelled the wind blown bubble and the ionisation structure, using both the data from NTD96 and from DTS97. For the wind-blown bubble, they used the Kompaneets model (Kompaneets 1960) for a shock wave propagating in an exponential atmostphere.
The most surprising result from the dynamical modelling is the implication of a very small scale height ($`H`$), namely 25 pc. This value was obtained by matching the aspect ratio of the superbubble as seen in H$`\alpha `$, assuming a distance of 2.35 kpc. The authors also point out that a small scale height is unavoidable in any model because the current maximum radius of the bubble must be significantly greater than $`H`$ in order for the bubble to have become so elongated.
## 3. New HI data
Data from the CGPS pilot project (Normandeau, Taylor, & Dewdney 1997) were combined with observations of six new fields taken with the Dominion Radio Astrophysical Observatory’s Synthesis Telescope. The right-hand panel of Figure 1 shows the resulting mosaic for velocities near –43.4 km s$`^1`$, again scaled by $`\mathrm{sin}(b)`$ to highlight the weaker high latitude emission.
### 3.1. Scale height
Using the data from complementary low resolution observations, which extend further up in latitude, almost to +10, and fitting an exponential to the decay in column density, a scale height of roughly 140 pc is found for the HI in the vicinity of the superbubble. Apart from the contradiction with the prediction by BJM99, this result is not particularly surprising: scale heights for the neutral medium are in the 100-200 pc range. It does suggest that the local medium into which the superbubble grew is quite different even from the relatively nearby gas.
### 3.2. Open or closed?
In the new data, there is no evidence of a cap at high latitudes in the HI. The eastern HI wall of the superbubble is clearly visible up to 5.5, at which point it curves slightly inward and disappears. The western wall is only well defined up to 3.4. An open geometry in HI images and a closed one in data showing ionised gas are not mutually exclusive as pointed out by BJM99. The superbubble’s shell could be sufficiently thin at high latitudes that while it closes the shell and prevents streaming of gas towards higher latitudes, it does not trap the ionizing radiation which then obliterates the HI at higher latitudes.
Extending approximately from (134.2, +6.0) to (134.2, +7.3), there is a small filament of HI pointing upward, away from the plane, which may be the signature of recent break out. This low level feature is present in four channels of the mosaic, from –41.76 km s$`^1`$to –46.70 km s$`^1`$. It is suggestive that it is perpendicular to the tangent to and at the (high longitude) extremity of an equally faint arc of HI . This arc follows the natural line that flows from the low longitude wall to the high longitude one as seen in the HI and seems to mark the boundary beyond which the HI emission is less at these velocities. The faint arc is below the latitude of the cap claimed by DTS97 and therefore below the upper boundary of the best fit Kompaneets model by BJM99. The upper tip of the vertical filament is above the latitude claimed for the ionized cap suggested by the DTS97 data.
### 3.3. Shape
It is noticeable that the OCl 352 cluster is at the base of the HI cone, as was described in NTD96 whereas the Kompaneets model by BJM99, which was fit primarily to the H$`\alpha `$ image, shows a bubble extending to lower latitudes, to the base of the ionized gas loop which forms the lower half of W4.
This cone shape with an energy source at the apex is not unique, something similar is seen at the base of the Aquila supershell (Maciejewski et al. 1996), but it is unclear how such a structure could form. The W4 cone extends approximately 200 pc upward from the OCl 352 cluster, based on the well-defined upper longitude wall. This is comparable to the cone at the base of the Aquila supershell which extends roughly 175 pc for an assumed distance of 3.3 kpc.
Tenorio-Tagle, Rózyczka, & Yorke (1985) present models of supernova remnants crossing large density discontinuities which result in shapes that are more conical than ellipsoidal. The opening angle of the structure above W4 is more regular however, it is more nearly a cone than are the models.
## 4. Conclusions and future endeavours
The classification of the structure above the W4 HII region, whether a superbubble or a chimney, remains undecided at this time. It may be an object in transition between the two phases of evolution. However, it is clear that there is no impediment to ionizing radiation escaping towards higher latitudes and therefore the OCl 352 cluster can contribute to maintaining the Reynolds layer of ionized gas via this conduit.
Along with the HI data, the DRAO observations also yield images of the radio continuum at two frequencies, with full polarimetric information at 1420 MHz. It is hoped that the latter data set will help to confirm or infirm the marginal detection of a possible cap in the H$`\alpha `$ emission.
The shape of the structure along with the scale height of the adjacent gas present challenges for modelling. The scale height inferred from the aspect ratio is surprisingly small and may reflect a local enhancement. The walls are very straight, making the structure conical rather than ellipsoidal as are the results of most models. The location of the presumed energy source at the apex of the cone rather than within it is also unusual.
## References
Basu, S., Johnstone, D., Martin, P. G. 1999, ApJ, 516, 843
Dennison, B., Topasna, G. A., Simonetti, J. H. 1997, ApJ, 474, L31
Heyer, M. H., Brunt, C., Snell, R. L., Howe, J., Schloerb, F. P., Carpenter, J. C., Normandeau, M., Taylor, A. R., Cao, Y., Terebey, S., Beichman, C.A. 1996, ApJ, 464, L175
Kompaneets, A. S. 1960, Sov. Phys. Dokl., 5, 46
Maciejewski, W., Murphy, E. M., Lockman, F. J., Savage, B. D. 1996, ApJ, 469, 238
Normandeau, M., Taylor, A. R., Dewdney, P. E. 1996, Nature , 380, 687
Normandeau, M., Taylor, A. R., Dewdney, P. E. 1997, ApJS, 108, 279
Taylor, A. R., Irwin, J. A., Matthews, H. E., Heyer, M. H. 1998, ApJ, 513, 339
Tenorio-Tagle, G., Rózyczka, M., Yorke, H. W. 1985, A&A, 148, 52
Weaver, R., McCray, R., Castor, J., Shapiro, P., Moore, R. 1977, ApJ, 218, 377
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# Statistical Physics of Vehicular Traffic and Some Related Systems
## 1 Introduction
The concepts and techniques of statistical physics are being used nowadays to study several aspects of complex systems many of which, till a few decades ago, used to fall outside the traditional domain of physical systems . Physical-, chemical-, earth-, biological- and social–sciences as well as technology meet at this frontier area of inter-disciplinary research. Flow of vehicular traffic and granular matter , folding of proteins , formation and growth of bacterial colonies , biological evolution of species and transactions in financial markets are just a few examples of exotic phenomena in such systems. Most of these systems are interesting not only from the point of view of Natural Sciences for fundamental understanding of how Nature works but also from the points of view of applied sciences and engineering for the potential practical use of the results of the investigations. Our review of the current status and future trends of research on the theory of vehicular traffic (and some related systems) will, we hope, convince you that, indeed, the results of recent studies of complex systems have been ”a conceptual revolution, a paradigm shift that has far reaching consequences for the very definition of physics” .
For almost half a century physicists have been trying to understand the fundamental principles governing the flow of vehicular traffic using theoretical approaches based on statistical physics . The approach of a physicist is usually quite different from that of a traffic engineer. A physicist would like to develop a model of traffic by incorporating only the most essential ingredients which are absolutely necessary to describe the general features of typical real traffic. The theoretical analysis and computer simulation of these models not only provide deep insight into the properties of the model but also help in better understanding of the complex phenomena observed in real traffic. Our aim in this review is to present a critical survey of the progress made so far towards understanding the fundamental aspects of traffic phenomena from the perspective of statistical physics.
There are two different conceptual frameworks for modeling traffic. In the ”coarse-grained” fluid-dynamical description, the traffic is viewed as a compressible fluid formed by the vehicles but these individual vehicles do not appear explicitly in the theory. In contrast, in the so-called ”microscopic” models of vehicular traffic attention is explicitly focused on individual vehicles each of which is represented by a ”particle”; the nature of the interactions among these particles is determined by the way the vehicles influence each others’ movement. In other words, in the ”microscopic” theories vehicular traffic is treated as a system of interacting ”particles” driven far from equilibrium. Thus, vehicular traffic offers the possibility to study various fundamental aspects of the dynamics of truly nonequilibrium systems which are of current interest in statistical physics .
In order to provide a broad perspective, we describe both ”macroscopic” and ”microscopic” approaches although we put more emphasis on the latter. Sometimes the phenomenological equations of traffic flow in the ”macroscopic” models can be obtained from microscopic considerations in the same spirit in which macroscopic or phenomenological theories of matter are derived from their molecular-theoretic description.
At present, even within the conceptual framework of ”microscopic” approach, there are several different types of mathematical formulations of the dynamical evolution of the system. For example, the probabilistic description of vehicular traffic in the kinetic theory is developed by appropriately modifying the kinetic theory of gases. On the other hand, a deterministic description of the motion of individual vehicles is provided by the so-called car-following theories which are based on the basic principles of classical Newtonian dynamics. In contrast, the so-called particle-hopping models describe traffic in terms of a stochastic dynamics of individual vehicles. We explain the guiding principles behind all these formulations. But we discuss in detail the results obtained mainly from the investigations of the recently developed ”particle-hopping” models which are usually formulated using the language of cellular automata (CA) . At present, there is no traffic model yet which can account for all aspects of vehicular traffic. In this review we consider a wide variety of CA models which describe various different types of traffic phenomena.
We map the particle-hopping models of vehicular traffic onto some other model systems; these mappings indicate the possibility of exploiting powerful techniques, used earlier for other systems, to study traffic models and, sometimes, enable us to obtain results for traffic models directly from the known results for models of other systems. We present pedagogical summaries of the statistical mechanical treatments of the CA models of traffic. We critically examine the regimes of validity of the approximation schemes of analytical calculations which we illustrate with explicit calculations in those limiting cases where these usually yield exact results. The results of the theoretical analysis of these models are compared with those obtained from computer simulations and, wherever possible, with the corresponding empirical results from real traffic. We also compare vehicular traffic with many other similar physical systems to show the ubiquity of some physical phenomena.
Computer simulations are known to provide sufficiently accurate quantitative results when analytical treatments require approximations which are too crude to yield results of comparable accuracy. In this review we demonstrate how computer simulations often help in getting deep insight into various phenomena involved in traffic and in qualitative understanding of the basic principles governing them thereby avoiding potentially hazardous experiments with real traffic. Computer simulations of the ”microscopic” models of traffic have not only attracted the attention of a growing number of statistical physicists in the recent years, but have also been received positively by many traffic engineers. The ongoing research efforts to utilize computer simulations of the microscopic models for practical applications in planning and design of transportation networks have been reviewed very recently by Nagel et al. .
Our review is complimentary to those published in the recent years by Helbing and by Nagel et al. . A large number of important papers on traffic published in recent years are based on the particle-hopping models. But these works have received very little attention in . We discuss the methods and results for the particle-hopping models in great detail in this review after explaining the basic principles of all the theoretical approaches. Moreover, we focus almost exclusively on the fundamental principles from the point of view of statistical physics while Nagel et al. emphasize practical applications which are directly relevant for traffic engineering.
At this point a skeptic may raise a serious question: ”can we ever predict traffic phenomena with statistical mechanical theories without taking into account effects arising from widely different human temperaments and driving habits of the individual drivers?” We admit that, unlike the particles in a gas, a driver is an intelligent agent who can ”think”, make individual decisions and ”learn” from experience. Besides, the action of the driver may also depend on his/her physical and mental states (e.g., sorrow, happiness, etc.). It is also true that the behavior of each individual driver does not enter explicitly into the ”microscopic” models of traffic. Nevertheless, as we shall show in this review, many general features of traffic can be explained in general terms with these models provided the different possible behavioural effects are captured collectively through a probabilistic description which requires only a few phenomenological parameters. Similar strategies have been suggested also for capturing the behavioural effects of some individual traders in financial markets collectively through probabilistic descriptions . These probabilistic descriptions make the dynamics of the models intrinsically stochastic.
Throughout this article the terms vehicle (or, car) and driver are used interchangeably, although each of these terms usually refers to the composite unit consisting of the vehicle and the driver. The main questions addressed by physicists are posed as problems in Section 2. As a motivation we present some relevant empirically observed general features of real traffic as well as their plausible phenomenological explanations in Section 3. The following sections review the different theoretical approaches. The classification of the models into different classes is not unique. Mostly we choose a classification according to the use of discrete or continuous space, time and state variables. The conceptual basis of the older theoretical approaches, namely, the fluid-dynamical theories, the kinetic theories and the car-following theories, are explained, in the Sections 4,5 and 6, respectively, where the corresponding recent developments are also summarized. These model classes are continuous in space, time and state variables. Some coupled map lattice models of traffic are considered in Section 7. They are discrete in time. CA models are discrete in space, time and state variables. The model suggested by Nagel and Schreckenberg (NaSch) is the minimal model of traffic on highways; the theoretical results on various aspects of this model are discussed in Section 8 where the nature of the spatio-temporal organization of vehicles are also investigated and the fundamental question of the (im-)possibility of any dynamical phase transition in the NaSch model is addressed. Various generalizations and extensions of the NaSch model (including those for multi-lane traffic) are reviewed in Section 9. The occurrence of self-organized criticality in the so-called cruise-control limit of the NaSch model is pointed out. It is demonstrated how additional ”slow-to-start” rules of CA can give rise to metastability, hysteresis and phase-separation in the generalized NaSch models, in qualitative agreement with empirical observations. In Section 10 the formation and ”coarsening” of platoons of vehicles are investigated in an appropriate generalization of the NaSch model with one type of quenched randomness; your attention is drawn to the formal analogy between this phenomenon and the Bose-Einstein condensation. In this section the effects of other kinds of quenched disorder on the nature of the steady-states of the NaSch model are also considered. In Section 11 we present some other CA models of highway traffic which are not directly related to the NaSch model. The Biham-Middleton-Levine (BML) model is the earliest CA model of traffic in idealized networks of streets in cities; it exhibits a first order phase transition. A critical review of this model is presented in Section 12 together with a list of its generalizations which have been reported so far. Furthermore a marriage of the NaSch description of traffic and the BML model, which has led to the development of a novel model of city traffic, is explained. A brief status report of the ongoing efforts to make practical use of the theoretical models for traffic engineering is also presented. The similarities between various particle-hopping models of traffic and some other systems far from equilibrium are pointed out in Section 13 followed by the concluding Section 14 where your attention is also drawn towards challenging open questions. Several Appendices deal mostly with more technical aspects of some important calculations, but are not necessary for an understanding of the main text.
## 2 Fundamental and practical questions
The aim of basic research in traffic science is to discover the fundamental laws governing traffic systems. The main aim of traffic engineering is on planning, design and implementation of transportation network and traffic control systems. Statistical physicists have been contributing to traffic science by developing models of traffic and drawing general conclusions about the basic principles governing traffic phenomena by studying these models using the tools of statistical physics. Moreover, using these models, statistical physicists have also been calculating several quantities which may find practical applications in traffic engineering. Furthermore, several groups of statistical physicists are currently also engaged in developing strategies for fast on-line simulation and traffic control so as to optimize traffic flow; significant contributions to this traditional domain of traffic engineering can reduce the financial burden on the governments.
### 2.1 Some fundamental questions
Because of the apparent similarities between the ”microscopic” models of traffic and macroscopic samples of ionic conductors in the presence of external electric field, the tools of statistical mechanics seem to be the natural choice for studying these models. However, the actual calculation of even the steady-state properties of traffic from the ”microscopic” models is a highly difficult problem because (apart from the human element involved) (a) the vehicles interact with each other and (b) the system is driven far from equilibrium, although it may attain a nonequilibrium steady-state.
In principle, the time-independent observable properties of large pieces of matter can be calculated within the general framework of equilibrium statistical mechanics, pioneered by Maxwell, Boltzmann and Gibbs, provided the system is in thermal equilibrium. Of course, in practice, it may not be possible to carry out the calculations without making approximations because of the interactions among the constituents of the system. Some time-dependent phenomena, e.g., fluctuation and relaxation, can also be investigated using the Linear Response Theory provided the system is not too far from equilibrium. Unfortunately, so far there is no general theoretical formalism for dealing with systems far from equilibrium. Moreover, the condition of detailed balance does not hold although a condition of pairwise balance holds for some special systems driven far from equilibrium.
The dynamical phases of systems driven far from equilibrium are counterparts of the stable phases of systems in equilibrium. Some of the fundamental questions related to the nature of these phases are as follows.
(i) What are the various dynamical phases of traffic? Does traffic exhibit phase-coexistence, phase transition, criticality or self-organized criticality and, if so, under which circumstances?
(ii) What is the nature of fluctuations around the steady-states of traffic? Analogous phenomenon of the fluctuations around stable states in equilibrium is by now quite well understood.
(iii) If the initial state is far from a stationary state of the driven system, how does it evolve with time to reach a truly steady-state? Analogous phenomena of equilibration of systems evolving from metastable or unstable initial states through nucleation (for example, in a supersaturated vapour) or spinodal decomposition (for example, in a binary alloy) have also been studied earlier extensively .
(iv) What are the effects of quenched (static or time-independent) disorder on the answers of the questions posed in (i)-(iii) above?
### 2.2 Some practical questions
Let is first define some characteristic quantitative features of vehicular traffic. The flux $`J`$, which is sometimes also called flow or current, is defined as the number of vehicles crossing a detector site per unit time . The distance from a selected point on the leading vehicle to the same point on the following vehicle is defined as the distance-headway . The time-headway is defined as the time interval between the departures (or arrivals) of two successive vehicles recorded by a detector placed at a fixed position on the highway . The distributions of distance-headways and time-headways are regarded as important characteristic of traffic flow. For example, larger headways provide greater margins of safety whereas higher capacities of the highway require smaller headways.
Let us now pose some questions which are of practical interest in traffic engineering.
(a) What is the relation between density $`c`$ and flux $`J`$? In traffic engineering, this relation is usually referred to as the fundamental diagram.
(b) What are the distributions of the distance-headway and time-headway?
(c) How should on- and off-ramps be designed?
(d) Does an additional lane really lead to an improvement?
(e) What are the effects of a new road on the performance of the road network?
(f) What type of signaling strategy should be adopted to optimize the traffic flow on a given network of streets and highways?
(g) The generalized travelling salesman problem: Is the shortest trip also the fastest?
## 3 Some empirical facts and phenomenological explanations
For several reasons, it is difficult to obtain very reliable (and reproducible) detailed empirical data on real traffic. First of all, unlike controlled experiments performed in the conventional fields of research in physical sciences, it is not possible to perform such laboratory experiments on vehicular traffic. In other words, empirical data are to be collected through passive observations rather than active experiments. Secondly, unambiguous interpretation of the collected data is also often a subtle exercise because traffic states depend on several external influences, e.g. the weather conditions. The systematic investigation of traffic flow has a quite long history . Although we now have a clear understanding of many aspects of real traffic several other controversial aspects still remain intellectual challenges for traffic scientists. In this section we give an overview of some of the well understood experimental findings, which are relevant for our theoretical analysis in the following sections. Moreover, wherever possible, we provide phenomenological explanations of these empirically observed traffic phenomena. Furthermore, we shall also mention some of the more recent empirical observations for which, at present, there are no generally accepted explanations.
### 3.1 Acceleration noise
In general, because of the different human temperaments and driving habits, different drivers react slightly differently to the same conditions on a highway, even when no other vehicle influences its motion. Consequently, even on an empty stretch of a highway, a driver can neither maintain a constant desired speed nor accelerate in a smooth fashion. In addition to the type of the highway (i.e., the surface conditions, frequency of the curves, etc.) the driver-to-driver fluctuation of the acceleration also depends on the density of vehicles on the highway. The root-mean-square deviation of the acceleration of the vehicles is a measure of the so-called acceleration noise. The distributions of the accelerations have been measured since mid nineteen fifties and are well documented .
### 3.2 Formation and characterization of traffic jams
Traffic jam is the most extensively studied traffic phenomenon. Traffic jams can emerge because of various different reasons. Most often traffic jams are observed at bottlenecks, e.g. lane-reductions or crossings of highways . At bottlenecks the capacity of the road is locally reduced thereby leading to the formation of jams upstream traffic. Downstream the bottleneck, typically, a free-flow region is observed. In addition, traffic accidents, which also lead to a local reduction of the capacity of the highway, can give rise to traffic jams.
Fig. 1 shows as an example empirical data of the velocity at a three-lane highway close to an on- and off-ramp . The data show that downstream the bottleneck (at detector C) no slow vehicle has been recorded. In contrast, in the merging regime near detector A vehicles often have to move slowly. In between the on- and off-ramps the vehicles move with larger velocities compared to those in location A although the number of vehicles passing detector B is maximal. Therefore, the on-ramp causes a local reduction of the capacity of the highway.
Perhaps, what makes the study of traffic jams so interesting is that jams often appear, as if, from nowhere (apparently without obvious reasons) suddenly on crowded highways; these so-called ”phantom jams”<sup>1</sup><sup>1</sup>1”Stau aus dem Nichts” in german are formed by spontaneous fluctuations in an otherwise streamlined flow. Direct empirical evidence for this spontaneous formation of jams was presented by Treiterer by analyzing a series of aerial photographs of a multi-lane highway. In Fig. 2 the picture from is redrawn. Each line represents the trajectory of an individual vehicle on one lane of the highway<sup>2</sup><sup>2</sup>2The discontinuous trajectories correspond to vehicles changing the lane.. The space-time plot (i.e., the trajectories $`x(t)`$ of the vehicles) shows the formation and propagation of a traffic jam. In the beginning of the analysed time vehicles are well separated from each other. Then, due to fluctuations, a dense region appears which, finally, leads to the formation of a jam. The jam remains stable for a certain period of time but, then, disappears again without any obvious reason. This figure clearly establishes not only the spontaneous formation of traffic jam but also shows that such jams can propagate upstream (opposite to the direction of flow of the vehicles). Moreover, it is possible that two or more jams coexist on a highway.
A more detailed analysis of traffic jams in absence of hindrances has been given by Kerner and Rehborn who pointed out the following characteristic features of wide jams. They found that the upstream velocity and, therefore, the outflow from a jam is approximately constant. The outflow from a jam and the velocity of the jam fronts are now regarded as two important empirical parameters of highway traffic which can be used for calibrating theoretical models.
### 3.3 Flux-density relation
Obviously, traffic flow phenomena strongly depend on the occupancy of the road. What type of variation of flux and average velocity $`v`$ with density $`c`$ can one expect on the basis of intuitive arguments? So long as $`c`$ is sufficiently small, the average speed $`v`$ is practically independent of $`c`$ as the vehicles are too far apart to interact mutually. Therefore, at sufficiently low density of vehicles, practically ”free flow” takes place. However, from the practical experience that vehicles have to move slower with increasing density, one expects that at intermediate densities,
$$\frac{dv}{dc}0,$$
(1)
when the forward movement of the vehicles is strongly hindred by others because of the reduction in the average separation between them. A faster-than-linear monotonic decrease of $`v`$ with increasing $`c`$ can lead to a maximum in the flux $`J=cv`$ at $`c=c_m`$; for $`c<c_m`$, increasing $`c`$ would lead to increasing $`J`$ whereas for $`c>c_m`$ sharp decrease of $`v`$ with increase of $`c`$ would lead to the overall decrease of $`J`$. However, contrary to this naive expectation, in recent years some nontrivial variation of flux with density have been observed. The nature of the variation of the flux with the density is still not clearly understood since the details of the complex experimental setup can strongly influence the empirical results.
Fig. 3 shows typical time averaged local measurements of the density and flow which have been obtained from the Queen Elizabeth Way in Ontario (Canada) . At low densities the data indicate a linear dependence of the flow on the density. In contrast strong fluctuations of the flow at large densities exist which prevents a direct evaluation of the functional form at high densities.
In order to use the empirical results for a theoretical analysis it is often more convenient to use the mean-values of the flow at a given density. Fig. 4 shows a collection of possible forms of averaged fundamental diagrams consistent with empirical data . While the discontinuity of the fundamental diagram now seems to be well established no clear answer can be given to the question on the form of the diagram in the free-flow or high-density regime. In the low density regime linear as well as non-linear functional forms of the fundamental diagrams have been suggested. For the high density branch no consistent picture for the high density branch exists. Here the results strongly depend on the specific road network.
In several situations it has been observed that $`J`$ does not depend uniquely on $`c`$ in an intermediate regime of density; it indicates the existence of hysteresis effects and meta-stable states. In the context of traffic flow, hysteresis effects have the following meaning: if a measurement starts in the free-flow regime, an increase of the density leads to an increase of the flow. However, beyond a certain density, a further increase of the density leads to a discontinuous reduction of the stationary flow (”capacity drop”) and jams emerge. The corresponding fundamental diagram (Fig. 4 A, B) then has the so-called ’inverse-$`\lambda `$ form’. Figure 5 shows an experimental verification of a hysteresis loop at a transition from a free flow to a congested state.
Recent empirical observations which have been obtained near a crossing of highways exhibit a flat plateau (i.e., a density-indepedent flux) over an intermediate regime of density of the vehicles.
### 3.4 Microscopic states of traffic flow and phase transitions
The results for the flux-density relation already suggest the existence of at least two different dynamical phases of vehicular traffic on highways, namely a free-flow phase and a congested phase. In the free-flow regime all vehicles can move with high speed close to the speed limit. The nature of the congested traffic is still under debate. Careful empirical observations in the recent years indicate the existence of two different congested phases, namely, the synchronized phase and the stop-and-go traffic phase . Vehicles move rather slowly in the synchronized states, as compared to the free-flow states, but the flux in the synchronized states can take a values close to the optimum value because of relatively small headways. Besides, non-trivial strong correlations between the density on different lanes exist in the synchronized state which actually motivates the notation synchronized traffic. The stop-and-go traffic differs from the synchronized states in the sense that every vehicle inside the jams come to a complete halt for a certain period of time.
Following Kerner three different types of synchronized traffic can be distringuished by the time-dependent behaviour of the density and flow . In synchronized traffic of type (i) constant values of density and flow can be observed during a long period of time. In synchronized traffic of type (ii) patterns of density and flow quite similar to free flow states have been observed. The differences between synchronized states of type (ii) and free-flow are given by the reduced average velocities and the alignment of the speeds on different lanes in synchronized traffic. Moreover irregular patterns of time-traced measurements of the flow have been found in synchronized traffic (see Fig. 6) of type (iii).
### 3.5 Time- and distance-headways
The flux $`J`$ can be written as $`J=N/T`$ where $`T=_{i=1}^Nt_i`$ is the sum of the time-headways recorded for all the $`N`$ vehicles. Hence, $`J=1/T_{av}`$ where $`T_{av}=(1/N)_it_i`$ is the average TH. Therefore, the TH distribution contains more detailed informations on traffic flow than that available from the flux alone. With the variation of density $`c`$ of the vehicles, $`T_{av}`$ exhibits a minimum at $`c=c_m`$ where the flux is maximum .
The results discussed in the preceeding subsections are based on time averaged local measurements. But it is also very useful to analyze the single-vehicle data directly . The single-vehicle data allows calculation of the time-headway distributions . All the time-headway distributions in the free-flow regime show a two peak structure. The first peak at $`\mathrm{\Delta }t=0.8\mathrm{sec}`$ represents the global maximum of the distribution. On a microscopic level these short time-headways correspond to platoons of some vehicles traveling very fast – their drivers are taking the risk of driving ”bumper-to-bumper” with a rather high speed. These platoons are the reason for the occurrence of high-flow states in free traffic. The corresponding states exhibit metastability, i.e. a perturbation of finite magnitude and duration is able to destroy such a high-flow state . Additionally, a second peak emerges at $`\mathrm{\Delta }t=1.8\mathrm{sec}`$ which can be associated with a typical drivers’ urge to maintain a temporal headway of $`2\mathrm{sec}`$ (which is the safe distance recommended in driving schools).
Surprisingly, the small time headways have much less weight in congested traffic. Only the peak at $`\mathrm{\Delta }t=1.8\mathrm{sec}`$ is recovered, where the time headway distribution now takes the maximum value. But nevertheless, the small time headways ($`\mathrm{\Delta }t<1.8\mathrm{sec}`$) contribute significantly in synchronized traffic. In stop-and-go traffic only the $`1.8\mathrm{sec}`$-peak remains and short time-headways are surpressed. The asymptotic behavior is rather unsystematic and reflects the dynamics of vehicles inside the jams.
Another important result characterizing the microscopic states is the dependence of the velocity of individual vehicles on the distance headway (see Fig. 8). This function is also of great importance for theoretical approaches, e.g. it is used as input for the so-called optimal velocity model. In the free-flow regime it is evident that the asymptotic velocity is reached already for small distance-headways. The slope of the velocity function is much lower than that in the free-flow regime. Surprisingly, the asymptotic velocity depends strongly on the microscopic state as well as on the density, e.g., in dense traffic low velocities of the vehicles are also observed even when large distance headways are available.
## 4 Fluid-dynamical Theories of vehicular traffic
When viewed from a long distance, say, an aircraft, flow of fairly heavy traffic appears like a stream of a fluid. Therefore, a ”macroscopic” theory of traffic can be developed, in analogy with the hydrodynamic theory of fluids, by treating traffic as an effectively one-dimensional compressible fluid (a continuum) . We follow the convention that traffic is flowing from left to right.
Suppose, $`c(x;t)`$ and $`J(x;t)`$ are the ”coarse-grained” density and flux at an arbitrary location $`x`$ at an arbitrary instant of time $`t`$. The equation of continuity for the fluid representing traffic is
$$\frac{c(x;t)}{t}+\frac{J(x;t)}{x}=\underset{i=1}{\overset{J_{in}}{}}\alpha _i(xx_i;t)\underset{j=1}{\overset{J_{out}}{}}\beta _j(xx_j;t)$$
(2)
where the first and the second terms on the right hand side take care of the sources and sinks, respectively, at the $`J_{in}`$ on-ramps situated at $`x_i`$ ($`i=1,2,\mathrm{},J_{in}`$) and $`J_{out}`$ off-ramps situated at $`x_j`$ ($`j=1,2,\mathrm{},J_{out}`$). We can write $`\alpha _i(xx_i;t)`$ and $`\beta _j(xx_j;t)`$ as
$$\alpha _i(xx_i;t)=\alpha _i^0(t)\varphi _i(xx_i)\text{and}\beta _j(xx_j;t)=\beta _j^0(t)\varphi _j(xx_j)$$
(3)
where $`\varphi _i(xx_i)`$ and $`\varphi _j(xx_j)`$ describe the spatial distribution of the incoming and outgoing flux, respectively, while $`\alpha _i^0(t)`$ and $`\beta _j^0(t)`$ account for the corresponding temporal variations.
In the following subsections, for simplicity, we shall consider a given stretch of highway with no entries or exits. In such special situations the equation of continuity reduces to the simpler form
$$\frac{c(x;t)}{t}+\frac{J(x;t)}{x}=0.$$
(4)
One cannot get two unknowns, namely, $`c(x;t)`$ and $`J(x;t)`$ (or, equivalently, $`v(x;t)`$) by solving only one equation, namely (4), unless they are related to each other. In order to proceed further, one needs another independent equation, say, for $`v(x;t)`$; we shall write down such an equation later in subsection 4.4. An alternative possibility, which Lighthill and Whitham adopted in their pioneering work, is to assume that $`J(x;t)`$ is determined primarily by the local density $`c(x;t)`$ so that $`J(x;t)`$ can be treated as a function of only $`c(x;t)`$. Consequently, the number of unknown variables is reduced to one as, according to this assumption, the two unknowns $`c(x;t)`$ and $`J(x;t)`$ are not independent of each other.
### 4.1 Lighthill-Whitham theory and kinematic waves
As a first approximation, let us begin with Lighthill-Whitham assumption that
$$J(x;t)=j(c(x;t))$$
(5)
where $`j(c)`$ is a function of $`c`$. The functional relation (5) between density and flux cannot be calculated within the framework of the fluid-dynamical theory; this must be either taken as a phenomenological relation extracted from empirical data or derived from more microscopic considerations. In general, the flux-density curve implied by equation (5) need not be identical with the fundamental diagram in the steady-state.
Under the assumption (5), the $`x`$-dependence of the local flux $`J(x;t)`$ arises only from the $`x`$-dependence of $`c(x;t)`$. Alternatively, since $`J(x;t)=c(x;t)v(x;t)`$, assuming $`v(x;t)`$ to depend only on $`c(x;t)`$ the $`x`$-dependence of $`v(x;t)`$ arises only from the $`x`$-dependence of $`c(x;t)`$. Using (5) the equation of continuity (4) can be expressed as
$$\frac{c(x;t)}{t}+\frac{c(x;t)}{x}\left[v(x;t)+c(x;t)\frac{dv}{dc}\right]=\frac{c(x;t)}{t}+v_g\frac{c(x;t)}{x}=0$$
(6)
where $`v_g=dJ/dc`$. Note that the equations (5) and (6) form the complete system of dynamical equations governing traffic flow in this first approximation. However, the equation (6) is non-linear because, in general, $`v_g=dJ/dc=v(c)+c\frac{dv(c)}{dc}`$ depends on $`c`$. If $`v_g`$ were a constant $`v_0`$, independent of $`c`$, equation (6) would become linear and the general solution would be of the form $`c(x;t)=f(xv_0t)`$ where $`f`$ is an arbitrary function of its argument. In that case, the solution of any particular problem would be found by merely matching the function $`f`$ to the corresponding given initial and boundary conditions. Such a solution describes a density wave motion as an initial density profile would get translated by a distance $`v_0t`$ in a time interval $`t`$ without any change in its shape. However, the non-linearity of the equation (6) gives rise to subtleties which are essential to capture at least some aspects of real traffic.
The solution of the nonlinear equation (6) is of the general form
$$c(x;t)=F(xv_gt),$$
(7)
where $`F`$ is an arbitrary function of its arguments. If we define a wave to be ”recognizable signal that is transferred from one part of a medium to another with a recognizable velocity of propagation” then the solutions of the form (7) can be regarded as a ”density wave”. There are several similarities between the density wave and the more commonly encountered waves like, for example, acoustic or elastic waves, But, the acoustic or elastic waves are solutions of linearized partial differential equations whereas the equation (6) is nonlinear and, hence, $`v_g`$ is $`c`$-dependent. Besides, the waves of the type (7) are called kinematic waves to emphasize their purely kinematic origin, in contrast to the dynamic origin of the accoustic and elastic waves. From the initial given density profile $`c(x;0)`$ the profile $`c(x;\mathrm{\Delta }t)`$ at time $`\mathrm{\Delta }t`$ can be obtained by moving each point on the initial profile a distance $`v_g(c)t`$ to the right; obviously, the distance moved is different for different values of $`c`$. The time-evolution of the density profile can be shown graphically on the space-time diagram (i.e., the $`xt`$ plane) where an arbitrary point $`x_0`$ on the $`t=0`$ axis moves along a straight line of slope $`v_g(c)`$ if the initial density at $`x_0`$ is $`c`$. These straight lines are referred to as characteristics; different characteristics corresponding to different $`c`$ have different slopes $`v_g(c)`$.
The speed $`v_g(c)`$ of the density wave should not be confused with $`v(c)`$, the actual speed of the continuum fluid representing traffic. In fact, at any instant of time $`v(x;t)`$ can be obtained from the corresponding density profile $`c(x;t)`$ by using the relation $`v(x;t)=j(c(x;t))/c(x;t)`$. Moreover, since $`v_g=v(c)+c\frac{dv(c)}{dc}`$ and since $`\frac{dv(c)}{dc}<0`$, the speed of the density wave is less than that of the fluid. Therefore, the density wave propagates backward relative to the traffic and the drivers are thereby warned of density fluctuations ahead downstream. Furthermore, the density wave moves forward or backward relative to the road, depending on whether $`c<c_m`$ or $`c>c_m`$ where $`c_m`$ corresponds to the maximum in the function $`j(c)`$.
When $`J(c)`$ is convex, i.e., $`d^2J/dc^2<0`$, we have $`dv_g/dc<0`$; consequently, higher values of $`c`$ propagate slower than lower values of $`c`$ thereby distorting the initial density profile. On the other hand, when $`dv_g/dc>0`$ higher values of $`c`$ propagate faster and the distortion has the opposite tendency as compared to the case of $`dv_g/dc<0`$. In both the situations the distortion of the initial density profile is caused by the $`c`$-dependence of $`v_g`$ which arises from the nonlinearity of the equation (6). The distortion of the density profile with time can also be followed on the space-time diagram. If $`dv_g/dc<0`$, in regions of decreasing density (i.e., $`c(x_1)>c(x_2)`$ for $`x_1<x_2`$) the characteristics move away from each other whereas, in regions of increasing density, the characteristics move towards each other.
When two characteristic lines on the space-time diagram intersect the density would be double-valued at the point of intersection. We can avoid this apparently impossible scenario by the following interpretation: When two characteristic lines intersect a shock wave is generated. By definition, a shock represents a mathematical discontinuity in $`c`$ and, hence, also in $`v`$. The speed of a shock wave is given by
$$v_s=\frac{J(c^+)J(c^{})}{c^+c^{}}$$
(8)
where $`c^+`$ and $`c^{}`$ are, respectively, the densities immediately in front (downstream) and behind (upstream) the shock while $`J(c^+)`$ and $`J(c^{})`$ represent the corresponding downstream and upstream fluxes, respectively. Note that the shock wave moves downstream (upstream) if $`v_s`$ is positive (negative). Often the shock is weak in the sense that the relative discontinuity $`(c^+c^{})/c^{}`$ is small and in such cases the shock wave speed tends to $`v_g=dJ/dc`$. As a shock separates a section of high an low densities of the model, it corresponds to a section of a highway where a free-flow and a congested regime is present. In particular for large differences between $`c^+`$ and $`c^{}`$ the velocity of the shock can be interpreted as the velocity of a backwards moving jam.
One advantage of the kinematic approach outlined above over any dynamic approach is that the dynamical equation, which will be given in Sec. 4.4, is difficult to derive from basic first principles and usually involve quite a few phenomenological parameters and even a phenomenological function. On the other hand, the only input needed for the kinematic approach is the phenomenological function $`J(c)`$ which can be obtained from empirical data.
### 4.2 Diffusion term in Lighthill-Whitham theory and its effects
Let us now make improvement over the original Lighthill-Whitham theory, which is based on the first approximation (5). We now assume that the local flux $`J(x;t)`$ is determined not only by the local density $`c(x;t)`$ but also by the gradient of the density. In other words, we replace the assumption (5) by
$$J(c)=j(c)D\frac{c}{x}$$
(9)
where $`D`$ is a positive constant. Note that, for fixed $`c(x;t)`$ (and, hence, fixed $`j(c)`$), a positive (negative) density gradient leads to a lower (higher) flux as the drivers are expected to reduce (increase) the speed of their vehicles depending on whether approaching a more (less) congested region. Using the relation (9) in the equation of continuity (4) we now get
$$\frac{c(x;t)}{t}+v_g\frac{c(x;t)}{x}=D\frac{^2c(x;t)}{x^2}$$
(10)
where $`v_g(c)=\frac{dj(c)}{dc}`$. The equation (10) reduces to the equation (6) when $`D=0`$. The nonlinearity and diffusion have opposite effects: the term $`v_g(c)\frac{c}{x}`$ leads to ”steepening” and ultimate ”breaking” of the wave whereas the term $`D^2c/x^2`$ smoothens out the profile. Nonvanishing $`D`$ also leads to a non-zero width of the shock wave.
### 4.3 Greenshields model and Burgers equation
So far in the preceding subsections we have not considered any specific form of the function $`j(c)`$ relating flux with density. One can start with the simplest (differentiable) approximation capturing the basic form of the fundamental diagram,
$$J=v_{max}c(1c).$$
(11)
Note that $`v_{max}`$ in (11) is a phenomenological parameter and it is interpreted to be the maximum average speed for $`c0`$. In traffic science and engineering, one usually uses $`1c/c_{jam}`$ instead of $`1c`$ in the equation (11) and the corresponding form of the relation between $`J`$ and $`c`$ is known as the Greenshields model. Substituting (11) into the equation (10) we get
$$\frac{c(x;t)}{t}+v_{max}\frac{c(x;t)}{x}2v_{max}c\frac{c(x;t)}{x}=D\frac{^2c}{x^2}$$
(12)
Introducing the linear transformation of variables
$$x=v_{max}t^{}x^{};t=t^{}$$
(13)
one gets
$$\frac{c(x;t)}{t^{}}+2v_{max}c\frac{c(x;t)}{x^{}}=D\frac{^2c}{x_{}^{}{}_{}{}^{2}}$$
(14)
which is the (deterministic) Burgers equation . Note that the transformation (13) takes one from the space-fixed coordinate system $`(x,t)`$ to a coordinate system $`(x^{},t^{})`$ that moves with uniform speed $`v_{max}`$; so, vehicles moving with speed $`v_{max}`$ with respect to the coordinate system $`(x,t)`$ do not move at all with respect to the coordinate system $`(x^{},t^{})`$.
The advantage of this route to the theory of traffic flow is that the Burgers equation (14) can be transformed further into a diffusion equation, thereby getting rid of the nonlinearity, through a nonlinear transformation called the Cole-Hopf transformation . Since it is straightforward to write down the formal solution to the diffusion equation, one can see clearly the role of the coefficient $`D`$ and the nature of the solutions in the limit $`D0`$.
If equation (6) is assumed to be the only equation governing traffic flow then an inhomogeneous initial state can lead to a shock wave but the amplitude of the shock wave decreases with time and eventually the shock wave fades out leading to a homogeneous steady state in the limit $`t\mathrm{}`$. Leibig has studied how a random initial distribution of steps in the density profile evolves with time in this theory. No traffic jam forms spontaneously from a state of uniform density at this level of sophistication of the fluid-dynamical approach.
### 4.4 Navier-Stokes-like momentum equation and consequences
Corresponding to the assumption (9) we can write a velocity equation
$$v(x;t)=v(c(x;t))\frac{D}{c}\frac{c(x;t)}{x}$$
(15)
where $`v(c)=\frac{j(c(x;t))}{c(x;t)}`$. In the kinematic approach discussed so far in the preceding subsections it is implicitly assumed that, following any change in the local density (and density gradient) is followed by an immediate response (without delay) of the velocity field. For a more realistic description, the local speed should be allowed to relax after a non-zero delay time $`t_d`$. So, it seems natural to treat the right hand side in (15) as a desired local velocity at $`x`$ and write the total derivative $`dv/dt`$ of $`v`$ with respect to time as
$$c\frac{v}{t}+v\frac{v}{x}=\frac{c}{t_d}\left[v_{safe}(c)v\right]D\frac{c(x;t)}{x}$$
(16)
where the function $`v_{safe}(c)`$ is identical to $`v(c)`$. Note that $`v_{safe}(c)`$ is a monotonically decreasing function of $`c`$, i.e, $`dv_{safe}/dc<0`$. Equation (16) is an additional dynamical equation describing the time-dependence of the velocity $`v(x;t)`$.
Now let us interpret the two terms on the right hand side of (16). The phenomenological function $`v_{safe}(c)`$ gives the safe speed, corresponding to the vehicle density $`c`$, that is achieved in time-independent and homogeneous traffic flow and $`t_d`$ is the corresponding average relaxation time. Next, note that the term $`D\frac{c(x;t)}{x}`$ takes into account the natural tendency of the drivers to accelerate (decelerate) if the density gradient is negative (positive), i.e. if the density in front becomes smaller (larger); therefore, it can be interpreted as proportional to the pressure gradient in the fluid describing traffic. In addition to these terms, another term proportional to $`\frac{^2v}{x^2}`$ is also added to the right hand side of the velocity equation; this tends to reduce spatial inhomogeneities of the velocity field and is usually interpreted as the analogue of the viscous dissipation term in the Navier-Stokes equation.
Thus, finally, in the fluid-dynamical approach, a complete mathematical description of the vehicular traffic on highways is provided by two equations, namely, the equation of continuity (2) and the Navier-Stokes-like velocity equation
$`c\left[{\displaystyle \frac{v}{t}}+v{\displaystyle \frac{v}{x}}\right]=D{\displaystyle \frac{c}{x}}+\mu {\displaystyle \frac{^2v}{x^2}}+{\displaystyle \frac{c}{t_d}}[v_{safe}(c)v]`$ (17)
where $`D,\mu `$ and $`t_d`$ are phenomenological constants.
### 4.5 Fluid-dynamical theories for multi-lane highways and city traffic
One can describe the traffic on two-lane highways by two equations each of the same form (2) and where the source term in the equation for lane $`1`$ (lane $`2`$) takes into account the vehicles which enter into it from the lane $`2`$ (lane $`1`$) while the sink term takes into account those vehicles entering the lane $`2`$ (lane $`1`$) from the lane $`1`$ (lane $`2`$).
A lattice hydrodynamic theory for city traffic has been formulated recently . This fluid-dynamical model is motivated by the CA model, developed by Biham et al. , which will be discussed in detail later in this review. Instead of generalizing the Navier-Stokes equation (17) a simpler form of the velocity equation has been assumed.
### 4.6 Some recent results of the fluid-dynamical theories and their physical implications
The fluid-dynamical model of vehicular traffic has been studied numerically by discretizing the partial differential equations (2) and (17) together with appropriate initial and boundary conditions. Both periodic boundary conditions and open boundary conditions with time-independent external flux $`\alpha _i^0(t)=\beta _i^0(t)=\gamma `$ have been considered.
In the fluid-dynamical theory, based on the equation of continuity (2) and the Navier-Stokes-like equation (17) traffic jams can appear spontaneously, even if the initial density profile $`c(x;0)`$ deviates very little from the homogeneous state (fig. 9). In order to understand the physical mechanism of the formation of local cluster of vehicles, let us consider a local increase of the density $`\mathrm{\Delta }c(x)>0`$ at some location $`x`$. Since $`dv_{safe}/dc<0`$, the local increase of the density leads to a decrease of $`v_{safe}`$. This decrease in the safe velocity forces drivers to reduce their average velocity $`v`$ sharply if $`|dv_{safe}/dc|`$ is large enough. On the other hand, it follows from the equation of continuity that the local decrease of $`v`$ gives rise to further increase of $`c`$ around $`x`$ and, consequently, further subsequent decrease of $`v(x)`$ in this location. This avalanche-like process, which tends to increase the amplitude of the local fluctuation of the density around the homogeneous state, competes against other processes, like diffusion and viscous dissipation, which tend to decrease inhomogeneities.
To our knowledge, the first attempt to understand the physical mechanism of synchronized traffic within the framework of the fluid-dynamical formalism was made by Lee et al. <sup>3</sup><sup>3</sup>3A brief summary of that work can be found in .. On a finite stretch of highway, of length $`L`$, they installed an on-ramp and an off-ramp on the model highway with a separation of $`L/2`$ between them. They chose the spatial distribution of the external flux $`\varphi (x)`$ in equation (3) as
$$\varphi (x)=(2\pi \sigma ^2)^{1/2}\mathrm{exp}(x^2/2\sigma ^2)$$
(18)
with $`\sigma =56.7m`$. They also assumed the form $`v_{safe}(c)=v_0(1c/\widehat{c})/[1+E(c/\widehat{c})^\theta ]`$ for the safe velocity with adjustable parameters $`v_0,E,\theta ,\widehat{c}`$. Lee et al. first allowed the system to reach the steady state after applying a weak time-independent flux $`\alpha _i^0(t)=\beta _i^0(t)=\gamma `$ and simulating the time-evolution of the traffic by solving simultaneously the equation of continuity (2) and the Navier-Stokes-like equation (17) with a specific set of chosen values for the parameters $`t_d,D,\mu ,v_0,E,\theta ,\widehat{c}`$, etc. Since they chose $`\gamma <\gamma _c`$, an initially homogeneous traffic reaches a steady ”free-flow” where homogeneous regions with different densities are separated from each other by narrow ”transition layers” near the ramps ( no stable ”free-flow” exists if $`\gamma \gamma _c`$). Then, they applied a pulse of additional flux $`\delta q`$ at the on-ramp for a short duration $`\delta t`$. After a transient period, which depends on the parameters of the model, the system was found to settle in a limit cycle in which the local density and local flux oscillate periodically and the oscillations are localized near the on-ramp. The discontinuous change of the spatio-temporally averaged velocity induced by the localized perturbations of finite amplitude, associated hysteresis effects and the stability of the limit cycle were found to be qualitatively similar to some of the empirically observed characteristics of synchronized flow in real traffic; therefore, Lee et al. identified the limit cycle observed in their theoretical investigation as the synchronized state of vehicular traffic. They drew analogy between this state and a ”self-excited oscillator” . However, this mechanism of the synchronized state is not yet accepted as the true and only possible explanation of the phenomena associated with the synchronized state observed empirically.
Meanwhile similar results have been obtained for a gas-kinetic based traffic model , also using on- and off-ramps in order to explain the transition from free-flow to synchronised states. This work was completed in a recent paper, where a phase diagram was calculated, which dependends on the on-ramp activity and the flow on the highway . Summarising the recent results of the macroscopic traffic models, there seems to be evidence that on-and off-ramps play an important role for a theoretical explanation of synchronised traffic. Nevertheless some experimental features are still not captured by these approaches. E.g. the empirical results show that for synchonised traffic of type (iii) no correlations between density and flow exist , in contrast to the regular patterns of the oscillating states found in simulations of the macroscopic models.
Despite its success in capturing many aspects of traffic flow the fluid-dynamical approach has its limitations; for example, viscosity of traffic is not a directly measurable quantity. Nevertheless, the fluid dynamical approach is being pursued not only by some physicists but also by several members of the traffic engineering community .
## 5 Kinetic theories of vehicular traffic
In the kinetic theory, traffic is treated as a gas of interacting particles where each particle represents a vehicle. The various different versions of the kinetic theory of vehicular traffic have been developed by modifying the kinetic theory of gases.
Recall that in the kinetic theory of gases $`f(\stackrel{}{r},\stackrel{}{p};t)d^3rd^3p`$ denotes the number of molecules which, at time $`t`$, have positions lying within a volume element $`d^3r`$ about $`\stackrel{}{r}`$ and momenta lying within the momentum-space element $`d^3p`$ about $`\stackrel{}{p}`$. The Boltzmann equation, which describes the time-evolution of the distribution $`f(x,v;t)`$, is given by
$$\left[\frac{f}{t}+\frac{\stackrel{}{p}}{m}_r+\stackrel{}{F}_p\right]f(\stackrel{}{r},\stackrel{}{p};t)=\left(\frac{f}{t}\right)_{coll}$$
(19)
where the symbols $`_r`$ and $`_p`$ denote gradient operators with respect to $`\stackrel{}{r}`$ and $`\stackrel{}{p}`$, respectively, while $`\stackrel{}{F}`$ is the external force. The term $`\left(\frac{f}{t}\right)_{coll}`$ represents the rate of change of $`f`$, with time, which is caused by the mutual collisions of the molecules.
In the first of the following two subsections we present the earliest version of the kinetic theory of vehicular traffic which was introduced by Prigogine and coworkers by modifying some of the key concepts in the kinetic theory of gases and by writing down an equation analogous to the Boltzmann equation (19). In the subsequent subsection we discuss the kinetic theory developed later by Paveri-Fontana to cure the defects from which the Prigogine theory was found to suffer.
### 5.1 Prigogine model
Suppose $`f(x,v;t)dxdv`$ denotes the number of vehicles, at time $`t`$, between $`x`$ and $`x+dx`$, having actual velocity between $`v`$ and $`v+dv`$. In addition, Prigogine and coworkers introduced a desired distribution $`f_{des}(x,v)`$ which is a mathematical idealization of the goals that the population of the drivers collectively strives to achieve. The actual distribution may deviate from the desired distribution because of various possible influences, e.g., road conditions, weather conditions or interaction with other vehicles, etc. They also argued that some of these influences cease after some time while the interactions with the other vehicles persist for ever. For example, only a short stretch of the road surface may be icy and strong winds or rain may stop after a short duration; in such situations $`f`$ can relax to $`f_{des}`$ over a relaxation time $`\tau _{rel}`$ provided mutual interactions of the vehicles is negligibly small. On the basis of these arguments, Prigogine and coworkers suggested that the analogue of the Boltzmann equation for the traffic should have the form
$$\frac{f}{t}+v\frac{f}{x}=\left(\frac{f}{t}\right)_{rel}+\left(\frac{f}{t}\right)_{int}$$
(20)
where $`\left(\frac{f}{t}\right)_{rel}`$ accounts for the relaxation of $`f`$ towards $`f_{des}`$ in the absence of mutual interactions of the vehicles while $`\left(\frac{f}{t}\right)_{int}`$ accounts for the changes of $`f`$ arising from mutual interactions among the vehicles. Note that the term $`\left(\frac{f}{t}\right)_{int}`$ on the right hand side of (20) may be interpreted as the analogue of the term $`\left(\frac{f}{t}\right)_{coll}`$ in the equation (19) whereas the term $`\left(\frac{f}{t}\right)_{rel}`$ in equation (20) may be interpreted as the counterpart of the term $`\stackrel{}{F}_pf(\stackrel{}{r},\stackrel{}{p};t)`$ in the equation (19).
Prigogine and coworkers wrote down an explicit form for the term $`\left(\frac{f}{t}\right)_{int}`$ by generalizing that for the term $`\left(\frac{f}{t}\right)_{coll}`$ in the kinetic theory of gases. We shall consider this term in the next subsection. In order to write down a simple explicit form of the relaxation term in the equation (20) they assumed that
(i) the collective relaxation, which would cause the actual distribution to tend towards the desired distribution, involves only a single relaxation time $`\tau _{rel}`$ so that
$$\left(\frac{f}{t}\right)_{rel}=\frac{ff_{des}}{\tau _{rel}}$$
(21)
and
(ii) the desired speed distribution $`F_{des}(v)`$ remains independent of the local concentration $`c(x;t)`$ so that
$$f_{des}(x,v;t)=c(x;t)F_{des}(v)$$
(22)
Therefore, a more explicit form of the Boltzmann-like equation (20) in the Prigogine theory is given by
$$\frac{f}{t}+v\frac{f}{x}=\frac{f(x,v;t)c(x;t)F_{des}(v)}{\tau _{rel}}+\left(\frac{f}{t}\right)_{int}$$
(23)
Note that, in the absence of mutual interactions of the vehicles, the distribution $`f(x,v;t)`$ would relax exponentially with time. The concept of desired distribution $`f_{des}(x,v;t)`$ and this scenario of collective relaxation of $`f`$ towards $`f_{des}`$ has subsequently come under severe criticism . Analyzing a set of ”ideal experiments” in the light of the Prigogine theory, Paveri-Fontana showed that the results obtained from the Boltzmann-like equation (23) are physically unsatisfactory.
More recently, Lehmann has attempted to revive the Prigogine approach by reformulating it as a semi- phenomenological theory where the distribution $`f(x,v;t)`$ is assumed to follow the simpler form
$$\frac{f}{t}+v\frac{f}{x}=\frac{ff_{des}(v,c)}{\tau _{rel}}$$
(24)
and the effects of the interactions are taken into account implicitly through a density-dependent desired distribution function $`f_{des}(v,c)`$ which has to determined empirically.
### 5.2 Paveri-Fontana model
In order to remove the conceptual as well as mathematical drawbacks of the Prigogine model of the kinetic theory of vehicular traffic, Paveri-Fontana argued that each vehicle, in contrast to the molecules in a gas, has a desired velocity towards which its actual velocity tends to ”relax” in the absence of ”interaction” with other vehicles. Thus, Paveri-Fontana model is based on a scenario of relaxation of the velocities of the individual vehicles rather than a collective relaxation of the distribution of the velocities.
In mathematical language, Paveri-Fontana introduced an additional phase-space coordinate, namely, the desired velocity. Suppose, $`g(x,v,v_{des};t)dxdvdv_{des}`$ denotes the number of vehicles at time $`t`$ between $`x`$ and $`x+dx`$, having actual velocity between $`v`$ and $`v+dv`$ and desired velocity between $`v_{des}`$ and $`v_{des}+dv_{des}`$. The one-vehicle actual velocity distribution function
$`f(x,v;t)={\displaystyle 𝑑v_{des}g(x,v,v_{des};t)}`$ (25)
describes the probability of finding a vehicle between $`x`$ and $`x+dx`$ having actual velocity between $`v`$ and $`v+dv`$ at time $`t`$. Similarly, the one-vehicle desired velocity distribution function
$`f_0(x,v_{des};t)={\displaystyle 𝑑vg(x,v,v_{des};t)}`$ (26)
describes the probability of finding a vehicle between $`x`$ and $`x+dx`$ having desired velocity between $`v_{des}`$ and $`v_{des}+dv_{des}`$. The local density of the vehicles $`c(x;t)`$ at the position $`x`$ at time $`t`$ can be obtained from
$$c(x;t)=_0^{\mathrm{}}𝑑v_{des}_0^{\mathrm{}}𝑑vg(x,v,v_{des};t).$$
Similarly, the corresponding average actual speed $`v(x;t)`$ and the average desired speed $`v_{des}(x;t)`$ are defined as
$$v(x;t)=\frac{_0^{\mathrm{}}𝑑v_{des}_0^{\mathrm{}}𝑑vvg(x,v,v_{des};t)}{c(x;t)}.$$
$$v_{des}(x;t)=\frac{_0^{\mathrm{}}𝑑v_{des}_0^{\mathrm{}}𝑑vv_{des}g(x,v,v_{des};t)}{c(x;t)}.$$
Finally, the local flux $`J(x;t)`$ is defined as $`J(x;t)=c(x;t)v(x;t)`$.
Now let us assume that the desired velocity of each individual driver is independent of time, i.e., $`dv_{des}/dt=0`$. Of course, the drivers may also adapt to the changing traffic environment and their desired velocities may change accordingly. In principle, these features can be incorporated into the kinetic theory at the cost of increasing complexity of the formalism.
Next, let us also assume that, in the absence of interaction with other vehicles, an arbitrary vehicle reaches the desired velocity exponentially with time, i.e., $`dv/dt=(v_{des}v)/\tau `$ where $`\tau `$ is a relaxation time. The Boltzmann-like kinetic equation for $`g(x,v,v_{des};t)`$ can be written as
$`\left[{\displaystyle \frac{}{t}}+v{\displaystyle \frac{}{x}}\right]g+{\displaystyle \frac{}{v}}\left[{\displaystyle \frac{v_{des}v}{\tau }}g\right]=\left({\displaystyle \frac{g}{t}}\right)_{int}`$ (27)
In order to write down an explicit form of the ”interaction term” we have to model the interactions among the vehicles. First of all, we model the vehicles as point-like objects. We consider the scenario where a fast vehicle, when hindered by a slow leading vehicle, either passes or slows down to the velocity of the lead vehicle. Let us now make some further simplifying assumptions:
(i) The slowing down takes place with a probability $`1P_{pass}`$ where $`P_{pass}`$ is the probability of passing.
(ii) If the fast vehicle passes the slower leading vehicle, its own velocity remains unchanged.
(iii) The velocity of the slower leading vehicle remains unchanged, irrespective of whether the faster following vehicle passes or slows down.
(iv) The slowing down process is instantaneous, i.e., the braking time is negligibly small.
(v) It is adequate to consider only two-vehicle interactions; there is no need to consider three-vehicle (or multi-vehicle) interactions.
(vi) The postulate of ”vehicular chaos”, which is the analogue of the postulate of ”molecular chaos” in the kinetic theory of gases, holds, so that the two-vehicle distribution function $`g_2(x,v,v_{des},x^{},v^{},v_{des}^{};t)`$ can be approximated as a product of two one-particle distributions $`g(x,v,v_{des};t)`$ and $`g(x^{},v^{},v_{des}^{};t)`$, i.e., $`g_2(x,v,v_{des},x^{},v^{},v_{des}^{};t)g(x,v,v_{des};t)g(x^{},v^{},v_{des}^{};t)`$. Thus, the equation (27) can be written explicitly as
$`\left[{\displaystyle \frac{}{t}}+v{\displaystyle \frac{}{x}}\right]g`$ $`+`$ $`{\displaystyle \frac{}{v}}\left[{\displaystyle \frac{v_{des}v}{\tau }}g\right]`$ (28)
$`=`$ $`f(x,v;t){\displaystyle _v^{\mathrm{}}}𝑑v^{}(1P_{pass})(v^{}v)g(x,v^{},v_{des};t)`$
$`g(x,v,v_{des};t){\displaystyle _0^v}𝑑v^{}(1P_{pass})(vv^{})f(x,v^{};t)`$
where the form of the ”interaction term” on the right-hand side of the equation (28) follows from the assumptions (i)-(vi) above. The first term on the right hand side of (28) describes the ”gain” of probability $`g(x,v,v_{des};t)`$ from the interaction of vehicles of actual velocity $`v^{}`$ with slower leading vehicle of actual velocity $`v`$ while the second term describes the loss of the probability $`g(x,v,v_{des};t)`$ arising from the interaction of vehicles of actual velocity $`v`$ with even slower leading vehicle of actual velocity $`v^{}`$.
The stationary homogeneous solution $`g(v,v_{des})`$ is, by definition, independent of $`x`$ and $`t`$. But, to our knowledge, so far it has not been possible to get even this solution of the Boltzmann-like integro-differential equation (28) by solving it analytically even for the simplest possible choice of the desired distribution function although numerical solutions provide some insights into the regimes of validity of the equation (28) and gives indications as to the directions of further improvements of the Paveri-Fontana model. For example, the finite sizes of the vehicles must be taken into account at high densities. Besides, the assumption (iv) of instantaneous relaxation has also been relaxed in a more recent extension .
Normally passing would require more than one lane on the highway. Therefore, the models discussed so far in the context of the kinetic theory may be regarded, more appropriately, as quasi-one-dimensional. These neither deal explicitly with $`g_i(x,v,v_{des};t)`$ for the individual lanes (labeled by $`i`$) nor take into account the process of lane-changing. Besides, all the vehicles were assumed to be of the same type. Now, in principle, we can generalize the formalism of the kinetic theory of traffic to deal with different types of vehicles on multi-lane highways. Suppose $`g_i^a(x,v,v_{des};t)`$ is the distribution for vehicles of type $`a`$ on the $`i`$-th lane of the highway. Obviously, the Boltzmann-like equations for the different lanes are coupled to each other. However, one needs additional postulates to model the lane-changing rules .
Very little work has been done so far on developing kinetic theories of two-dimensional traffic flow which would represent, for example, traffic in cities. Suppose, for simplicity, that the network of the streets consists of north-south and east-west streets and that east-west streets allow only east-bound traffic while only north-bound traffic flow takes place along the north-south streets. Let $`P_{xx}`$ ($`P_{xy}`$) denote the probabilities of an east-bound vehicle passing another east-bound (north-bound) vehicle. Similarly, suppose, $`P_{yx}`$ ($`P_{yy}`$) denote the probabilities of a north-bound vehicle passing an east-bound (north-bound) vehicle. Let $`g_x(x,y,u,u_{des};t)`$ and $`g_y(x,y,v,v_{des};t)`$ represent the distributions for the east-bound and north-bound vehicles, respectively, where $`u`$ and $`v`$ refer to the actual their actual velocities whereas $`u_{des}`$ and $`v_{des}`$ refer to the corresponding desired velocities. The Boltzmann-like equations governing the time evolutions of these distributions are given by
$$\left[\frac{}{t}+u\frac{}{x}+v\frac{}{y}\right]g_x+\frac{}{u}\left[\frac{u_{des}u}{\tau }g_x\right]=\left(\frac{g_x}{t}\right)_{x,coll}+\left(\frac{g_x}{t}\right)_{y,coll}$$
(29)
where
$`\left({\displaystyle \frac{g_x}{t}}\right)_{x,coll}`$ $`=`$ $`f_x(x,y,u;t){\displaystyle _u^{\mathrm{}}}𝑑u^{}(1P_{xx})(u^{}u)g_x(x,y,u^{},u_{des};t)`$ (30)
$``$ $`g_x(x,y,u,u_{des};t){\displaystyle _0^u}𝑑u^{}(1P_{xx})(uu^{})f_x(x,y,u^{};t)`$
and
$`\left({\displaystyle \frac{g_x}{t}}\right)_{y,coll}`$ $`=`$ $`\delta _{u,0}{\displaystyle _0^{\mathrm{}}}𝑑v^{}f_y(x,y,v^{};t){\displaystyle _0^{\mathrm{}}}𝑑u^{}(1P_{xy})u^{}g_x(x,y,u^{},u_{des};t)`$ (31)
$``$ $`g_x(x,y,u,u_{des};t){\displaystyle _0^{\mathrm{}}}𝑑v^{}(1P_{xy})uf_y(x,y,v^{};t)`$
The first term on the right hand side of the equation (30) describes gain of population of east-bound vehicles with velocity $`u`$ because of interaction with other east-bound vehicles with velocity $`u^{}u`$ while the second term describes the loss of population of east-bound vehicles because of interaction with east-bound vehicles of velocity $`u^{}<u`$. The right hand side of the equation (31) is based on the assumption that when an east-bound vehicle interacts with a north-bound vehicle at a crossing, it either passes or stops.
### 5.3 Derivation of the phenomenological equations of the macroscopic fluid-dynamical theories from the microscopic gas-kinetic models
In this section we discuss the results of the attempts to derive the phenomenological equations of traffic flow in the macroscopic fluid-dynamical theories from the microscopic gas-kinetic models. Several attempts have been made so far to derive the equation of continuity and the Navier-Stokes-like equation for traffic from the corresponding Boltzmann-like equation in the same spirit in which the derivations of the equation of continuity and Navier-Stokes equation for viscous fluids from the Boltzmann-equation have been carried out. However, because of the postulate of ”vehicular chaos”, the equation (28) is expected to be valid only at very low densities where the correlations between the vehicles is negligibly small whereas traffic is better approximated as a continuum fluid at higher densities!
Let us define the moments
$$m_{k,\mathrm{}}(x;t)=𝑑v𝑑v_{des}v^kv_{des}^{\mathrm{}}g(x,v,v_{des};t)$$
(32)
Note that $`c=m_{0,0}`$, $`v=m_{1,0}`$. Integrating the Boltzmann-like equation (28) over the actual velocities we get
$$\frac{}{t}f_0(x,v_{des};t)+\frac{}{x}[\overline{v}(x,v_{des};t)f_0(x,v_{des};t)]=0$$
(33)
where $`\overline{v}(x,v_{des};t)`$ is defined as
$$\overline{v}(x,v_{des};t)=\frac{𝑑vvg(x,v,v_{des};t)}{f_0(x,v_{des};t)}$$
(34)
The equation (33) is an equation of continuity for each desired speed $`v_{des}`$ separately; it is consequence of the assumption that $`dv_{des}/dt=0`$, i.e., no driver changes the desired speed.
Using the Boltzmann-like equation (28) and the definition (32) we can get separate partial differential equations for the moments of $`v`$, moments of $`v_{des}`$ and the mixed moments of $`v`$ and $`v_{des}`$. Unfortunately, these lead to a hierarchy of moment equations where each evolution equation for moment of a given order involves also moments of the next higher order. In order to close this system of equations, one needs to make appropriate justifiable assumptions.
## 6 Car-following theories of vehicular traffic
In the car-following theories one writes, for each individual vehicle, an equation of motion which is the analogue of the Newton’s equation for each individual particle in a system of interacting classical particles. In Newtonian mechanics, the acceleration may be regarded as the response of the particle to the stimulus it receives in the form of force which includes both the external force as well as those arising from its interaction with all the other particles in the system. Therefore, the basic philosophy of the car-following theories can be summarized by the equation
$$[Response]_n[Stimulus]_n$$
(35)
for the $`n`$-th vehicle ($`n=1,2,\mathrm{}`$). Each driver can respond to the surrounding traffic conditions only by accelerating or decelerating the vehicle. Different forms of the equations of motion of the vehicles in the different versions of the car-following models arise from the differences in their postulates regarding the nature of the stimulus (i.e., ”behavioural force” or a ”generalized force” ). The stimulus may be composed of the speed of the vehicle, the difference in the speeds of the vehicle under consideration and its lead vehicle, the distance-headway, etc., and, therefore, in general,
$$\ddot{x}_n=f_{sti}(v_n,\mathrm{\Delta }x_n,\mathrm{\Delta }v_n)$$
(36)
where the function $`f_{sti}`$ represents the stimulus received by the $`n`$-th vehicle. Different versions of the car-following models model the function $`f_{sti}`$ differently. In the next two subsections we discuss two different conceptual frameworks for modeling $`f_{sti}`$.
### 6.1 Follow-the-leader model
In the earliest car-following models the difference in the velocities of the $`n`$-th and $`(n+1)`$-th vehicles was assumed to be the stimulus for the $`n`$-th vehicle<sup>4</sup><sup>4</sup>4In the following we label the vehicles in driving direction such that the $`(n+1)`$-th vehicle is in front of the $`n`$-th vehicle.. In other words, it was assumed that every driver tends to move with the same speed as that of the corresponding leading vehicle so that
$$\ddot{x}_n(t)=\frac{1}{\tau }\left[\dot{x}_{n+1}(t)\dot{x}_n(t)\right]$$
(37)
where $`\tau `$ is a parameter that sets the time scale of the model. Note that $`1/\tau `$ in the equation (37) can be interpreted as a measure of the sensitivity coefficient $`𝒮`$ of the driver; it indicates how strongly the driver responds responds to unit stimulus. According to such models (and their generalizations proposed in the fifties and sixties) the driving strategy is to follow the leader and, therefore, such car-following models are collectively referred to as the follow-the-leader model.
Pipes derived the equation (37) by differentiating, with respect to time, both sides of the equation
$$\mathrm{\Delta }x_n(t)=x_{n+1}(t)x_n(t)=(\mathrm{\Delta }x)_{safe}+\tau \dot{x}_n(t)$$
(38)
which encapsulates his basic assumption that (a) the higher is the speed of the vehicle the larger should be the distance-headway, and, (b) in order to avoid collision with the leading vehicle, each driver must maintain a safe distance $`(\mathrm{\Delta }x)_{safe}`$ from the leading vehicle.
It has been argued that, for a more realistic description, the strength of the response of a driver at time $`t`$ should depend on the stimulus received from the other vehicles at time $`tT`$ where $`T`$ is a response time lag. Therefore, generalizing the equation (37) one would get
$$\ddot{x}_n(t+T)=𝒮[\dot{x}_{n+1}(t)\dot{x}_n(t)]$$
(39)
where the sensitivity coefficient $`𝒮`$ is a constant independent of $`n`$.
According to the equations (37) and (39), a vehicle would accelerate or decelerate to acquire the same speed as that of its leading vehicle. This implies that, as if, slower following vehicle are dragged by their faster leading vehicle. In these linear dynamical models the acceleration response of a driver is completely independent of the distance-headway. Therefore, this oversimplified equation, fails to account for the clustering of the vehicles observed in real traffic. Moreover, since there is no density-dependence in this dynamical equation, the fundamental relation cannot be derived from this dynamics. In order to make the model more realistic, we now assume that the closer is the $`n`$-th vehicle to the $`(n+1)`$-th the higher is the sensitivity of the driver of the $`n`$-th car. In this case, the dynamical equation (39) is further generalized to
$$\ddot{x}_n(t+T)=\frac{\kappa }{[x_{n+1}(t)x_n(t)]}[\dot{x}_{n+1}(t)\dot{x}_n(t)]$$
(40)
where $`\kappa `$ is a constant. An even further generalization of the the model can be achieved by expressing the sensitivity factor for the $`n`$-th driver as
$$𝒮_n=\frac{\kappa [v_n(t+\tau )]^m}{[x_{n+1}(t)x_n(t)]^{\mathrm{}}}$$
(41)
where $`\mathrm{}`$ and $`m`$ are phenomenological parameters to be fixed by comparison with empirical data. These generalized follow-the-leader models lead to coupled non-linear differential equations for $`x_n`$. Thus, in this ”microscopic” theoretical approach, the problem of traffic flow reduces to problems of nonlinear dynamics.
So far as the stability analysis is concerned, there are two types of analysis that are usually carried out. The local stability analysis gives information on the nature of the response offered by the following vehicle to a fluctuation in the motion of its leading vehicle. On the other hand, the manner in which a fluctuation in the motion of any vehicle is propagated over a long distance through a sequence of vehicles can be obtained from an asymptotic stability analysis.
From experience with real traffic we know that drivers often observe not only the leading vehicle but also a few other vehicles ahead of the leading vehicle. For example, the effect of the leading vehicle of the leading vehicle can be incorporated in the same spirit as the effect of ”next-nearest-neighbour” in various lattice models in statistical mechanics. A linear dynamical equation, which takes into account this ”next-nearest-neighbour” within the framework of the follow-the-leader model, can be written as
$$\ddot{x}_n(t+T)=𝒮^{(1)}[\dot{x}_{n+1}(t)\dot{x}_n(t)]+𝒮^{(2)}[\dot{x}_{n+2}(t)\dot{x}_n(t)]$$
(42)
where $`𝒮^{(1)}`$ and $`𝒮^{(2)}`$ are two phenomenological response coefficients.
The weakest point of these theories is that these involve several phenomenological parameters which are determined through ”calibration”, i.e., by fitting some predictions of the theory with the corresponding empirical data. Besides, extension of these models to multi-lane traffic is difficult since every driver is satisfied if he/she can attain the desired speed!
### 6.2 Optimal velocity models
We can express the driving strategy of the driver in the car-following models in terms of mathematical symbols by writing
$$\ddot{x}_n(t)=\frac{1}{\tau }\left[V_n^{desired}(t)v_n(t)\right]$$
(43)
where $`V_n^{desired}(t)`$ is the desired speed of the $`n`$-th driver at time $`t`$. In all follow-the-leader models mentioned above the driver maintains a safe distance from the leading vehicle by choosing the speed of the leading vehicle as his/her own desired speed, i.e., $`V_n^{desired}(t)=\dot{x}_{n+1}`$.
An alternative possibility has been explored in recent works based on the car-following approach . This formulation is based on the assumption that $`V_n^{desired}`$ depends on the distance-headway of the $`n`$-th vehicle, i.e., $`V_n^{desired}(t)=V^{opt}(\mathrm{\Delta }x_n(t))`$ so that
$$\ddot{x}_n(t)=\frac{1}{\tau }\left[V^{opt}(\mathrm{\Delta }x_n(t))v_n(t)\right]$$
(44)
where the so-called optimal velocity function $`V^{opt}(\mathrm{\Delta }x_n)`$ depends on the corresponding instantaneous distance-headway $`\mathrm{\Delta }x_n(t)=x_{n+1}(t)x_n(t)`$. In other words, according to this alternative driving strategy, the $`n`$-th vehicle tends to maintain a safe speed that depends on the relative position, rather than relative velocity, of the $`n`$-th vehicle. In general, $`V^{opt}(\mathrm{\Delta }x)0`$ as $`\mathrm{\Delta }x0`$ and must be bounded for $`\mathrm{\Delta }x\mathrm{}`$. For explicit calculations, one has to postulate a specific functional form of $`V^{opt}(\mathrm{\Delta }x)`$. Car-following models along this line of approach have been introduced by Bando et al. . For obvious reasons, these models are usually referred to as optimal velocity (OV) models.
Since the equations of motion in the follow-the-leader models involve only the velocities, and not positions, of the vehicles these can be formulated as essentially first order differential equations (for velocities) with respect to time. In contrast, since the equations of motion in the OV model involve the positions of the vehicles explicitly, the theoretical problems of this model are formulated mathematically in terms of second order differential equations (for the positions of the vehicles) with respect to time .
The simplest choice for $`V^{opt}(\mathrm{\Delta }x)`$ is
$$V^{opt}(\mathrm{\Delta }x)=v_{max}\mathrm{\Theta }(\mathrm{\Delta }xd),$$
(45)
where $`d`$ is a constant and $`\mathrm{\Theta }`$ is the Heavyside step function. According this form of $`V^{opt}(\mathrm{\Delta }x)`$, a vehicle should stop if the corresponding distance-headway is less than $`d`$; otherwise, it can accelerate so as to reach the maximum allowed velocity $`v_{max}`$. A somewhat more realistic choice is
$$V^{opt}(\mathrm{\Delta }x)=\{\begin{array}{cc}0\hfill & \mathrm{for}\mathrm{\Delta }x<\mathrm{\Delta }x_A,\hfill \\ f\mathrm{\Delta }x\hfill & \mathrm{for}\mathrm{\Delta }x_A\mathrm{\Delta }x\mathrm{\Delta }x_B,\hfill \\ v_{max}\hfill & \mathrm{for}\mathrm{\Delta }x_B<\mathrm{\Delta }x.\hfill \end{array}$$
(46)
The main advantage of the forms (45) and (46) of the OV function is that exact analytical calculations, e.g. in the jammed region, are possible . Although (45) and (46) may not appear very realistic, they capture several key features of more realistic forms of OV function , e.g.,
$$V^{opt}(\mathrm{\Delta }x)=\mathrm{tanh}[\mathrm{\Delta }x\mathrm{\Delta }x_c]+\mathrm{tanh}[\mathrm{\Delta }x_c].$$
(47)
for which analytical calculations are very difficult. For the convenience of numerical investigation, the dynamical equation (44) for the vehicles in the OV model has been discretized and, then, rewritten as a difference equation .
The main question addressed by the OV model is the following: what is the condition for the stability of the homogeneous solution
$$x_n^h=bn+ct$$
(48)
where $`b=(\mathrm{\Delta }x)_{av}=L/N`$ is the constant average spacing between the vehicles and $`c`$ is the constant velocity. It is not difficult to argue that, in general, in the OV models the homogeneous flow becomes unstable when $`\frac{V^{opt}}{\mathrm{\Delta }x}|_{\mathrm{\Delta }x=b}>\frac{2}{\tau }`$ .
One can distinguish five different density regimes with respect to the stability of microscopic states (see Fig. 10). At low and at high densities the homogeneous states are stable. For intermediate densities one can distinguish three regimes where jammed states exist. In region III the jammed state is stable whereas in regions II and IV both homogeneous and jammed states form stable structures. Beyond the formation of jams also hysteresis effects have been observed. Thus the OV model is able to reproduce many aspects of experimental findings.
Modified Korteweg-de Vries (KdV) equation has been derived from the equation (44) in a special regime of the parameters and the relations between its kink solutions and traffic congestion have been elucidated . A generalization to two-lane traffic can be found in . In order to account for traffic consisting of two different types of vehicles, say, cars and trucks, Mason et al. generalized the formulation of Bando et al. by replacing the constant $`\tau `$ by $`\tau _n`$ so that
$$\ddot{x}_n(t)=\frac{1}{\tau _n}[V_n^{opt}(t)v_n(t)]$$
(49)
where $`\tau _n`$ now depends on whether the $`n`$-th vehicle is a car or a truck. Since a truck is expected to take longer to respond than a car we should assign larger $`\tau `$ to trucks and smaller $`\tau `$ to cars. Some other mathematically motivated generalizations of the OV model have also been considered .
As mentioned earlier, drivers often receive stimulus not only from the leading vehicle but also a few other vehicles ahead of the leading vehicle. One possible way to generalize the OV models for taking into account such multi-vehicle interactions is to write the dynamical equations as
$$\ddot{x}_n=\underset{j=1}{\overset{m}{}}𝒮_j\left[V^{opt}\left(\frac{x_{n+j}x_n}{j}\right)v_n\right]$$
(50)
where $`𝒮_j`$ are sensitivity coefficients; one of the commonly used explicit forms of the optimal velocity function, for example (47) can be chosen for that of the function $`V^{opt}`$ in (50).
Before concluding our discussions on the OV models let us mention a very simple model where one assumes that
$$\ddot{x}_n(t)=\{\begin{array}{cc}a\hfill & \mathrm{for}\mathrm{\Delta }x_n\mathrm{\Delta }x_c\hfill \\ a\hfill & \mathrm{for}\mathrm{\Delta }x_n<\mathrm{\Delta }x_c.\hfill \end{array}$$
(51)
with $`a>0`$. Obviously, $`\mathrm{\Delta }x_c`$ may be interpreted as the safety distance. Moreover, a restriction $`v_{min}v_nv_{max}`$ is imposed on the allowed velocities of the vehicles by introducing the allowed minimum and maximum velocities $`v_{min}`$ and $`v_{max}`$, respectively. Note that in this oversimplified model $`\ddot{x}_n`$ depends on the corresponding distance-headway and, therefore, has some apparent similarities with the OV models. But, unlike the more general OV models, in this model $`\ddot{x}_n`$ does not depend on its instantaneous velocity. The main reason for considering such a oversimplified scenario is that velocity of the propagation of jams can be calculated analytically.
Some ideas of the OV model has been utilized by Mahnke and Pieret in their master equation approach to the study of jam dynamics. They assumed that, at a time, only one vehicle can go into or come out of a jam; this, naturally, does not take into the merging or splitting of jams. Under this assumption, the master equation (see e.g. App. F, equation (176)) governing the probability distribution $`P(n;t)`$ of the jam sizes $`n`$ is given by
$`{\displaystyle \frac{dP(n;t)}{dt}}`$ $`=`$ $`W_+(n1)P(n1;t)+W_{}(n+1)P(n+1;t)`$ (52)
$``$ $`[W_+(n)+W_{}(n)]P(n;t)`$
where $`W_+`$ and $`W_{}`$, are the ”growth” and ”decay” transition rates, respectively. It has been argued that $`W_{}(n)=1/\tau =constant`$ since a vehicle would require a constant average time $`\tau `$ to come out of a jam. However, $`W_+(n)`$ would depend on $`V^{opt}(\mathrm{\Delta }x)`$ since the time taken by a free-flowing vehicle immediately behind a jam to get into the jam would depend on $`\mathrm{\Delta }x`$ as well as on $`V^{opt}(\mathrm{\Delta }x)`$ although the actual expression of $`W_+(n)`$ may be complicated in a reasonable ansatz . Moreover, several assumptions of the model will have to be relaxed before the results of this approach can be compared with those from real traffic.
Before concluding this section we would like to emphasize that while formulating the dynamical equations for updating the velocities and positions of the vehicles in any ”microscopic” theory the following points should be considered:
(i) in the absence of any disturbance from the road conditions and interactions with other vehicles, a driver tends to drive with a desired velocity $`v_{des}`$; if the actual current velocity of the vehicle $`v(t)`$ is smaller (larger) than $`v_{des}`$, the vehicle accelerates (decelerates) so as to approach $`v_{des}`$.
(ii) In a freely-flowing traffic, even when a driver succeeds in attaining the desired velocity $`v_{des}`$, the velocity of the vehicle fluctuates around $`v_{des}`$ rather than remaining constant in time.
(iii) The interactions between a pair of successive vehicles in a lane cannot be neglected if the gap between them is shorter than $`v_{des}`$; in such situations the following vehicle must decelerate so as to avoid collision with the leading vehicle.
Clearly, the reliability of the predictions of the OV model depends on the appropriate choice of the optimal velocity function.
## 7 Coupled-map lattice models of vehicular traffic
Recall that, in the car-following models, space is assumed to be a continuum and time is represented by a continuous variable $`t`$. Besides, velocity and acceleration of the individual vehicles are also real variables. However, most often, for numerical manipulations of the differential equations of the car-following models, one needs to discretize the continuous variables with appropriately chosen grids. In contrast, in the coupled-map lattice approach , one starts with a discrete time variable and, the dynamical equations for the individual vehicles are formulated as discrete dynamical maps that relate the state variables at time $`t`$ with those at time $`t+1`$, although position, velocity and acceleration are not restricted to discrete integer values. The unit of time in this scheme (i.e., one time step) may be interpreted as the reaction time of the individual drivers as the velocity of of a vehicle at the time step $`t`$ depends on the traffic conditions at the preceding time step $`t1`$.
Keeping in mind the general points raised at the end of the preceding section regarding the formulation of the dynamical equations for updating the velocities and positions of the vehicles, the general forms of the dynamical maps in the coupled-map lattice models can be expressed as:
$`v_n(t+1)`$ $`=`$ $`Map_n[v_n(t),v_{des},\mathrm{\Delta }x_n(t)],`$ (53)
$`x_n(t+1)`$ $`=`$ $`v_n(t)+x_n(t)`$ (54)
where $`v_{des}`$ is a desired velocity. In general, the dynamical map $`Map[v_n(t),v_{des},\mathrm{\Delta }x_n(t)]`$ takes into account the velocity $`v_n(t)`$ and the distance-headway $`\mathrm{\Delta }x_n(t)`$ of the $`n`$-th vehicle at time $`t`$ for deciding the velocity $`v_n(t+1)`$ at time $`t+1`$. The effects of the interactions among the vehicles enter into the dynamical updating rules (53), (54) only through the distance-headway $`\mathrm{\Delta }x_n`$.
### 7.1 The model of Yukawa and Kikuchi
Yukawa and Kikuchi have studied coupled-map models based on the map
$$v(t+1)=F(v(t)):=\gamma v(t)+\beta \mathrm{tanh}\left(\frac{v^Fv(t)}{\gamma }\right)+ϵ$$
(55)
for the uninfluenced motion of a single vehicle. $`v^F`$ is the preferred velocity of the vehicle and $`\beta `$, $`\gamma `$, $`\delta `$ and $`ϵ`$ are parameters. For $`\gamma `$ close to 1 the map becomes chaotic, but acceleration and deceleration are approximately constant far from $`v^F`$. Their magnitude is determined by the parameter $`\beta `$. $`ϵ`$ controls the difference of the acceleration and deceleration capabilities. Although the model is deterministic, fluctuations in the velocity are introduced through deterministic chaos. These fluctuations around $`v^F`$ are determined by the parameter $`\delta `$.
If there is more than one vehicle on the road one needs an additional deceleration mechanism to avoid collisions. This can be achieved by introducing a deceleration map. Assuming that deceleration is dominated by the headway, two models have been studied in . In model A, the deceleration map describes a sudden braking process. If the front-bumper to front-bumper distance $`\mathrm{\Delta }x_n`$ to the next vehicle ahead is less than the current velocity $`v_n(t)`$ of the following vehicle, then the velocity is reduced to $`\mathrm{\Delta }x_nl`$ where $`l`$ is the length of the vehicles. The corresponding map is $`B(\mathrm{\Delta }x_n(t))=\mathrm{\Delta }x_n(t)l`$.
Model B has a more complex deceleration map:
$`v_n(t+1)`$ $`=`$ $`G(\mathrm{\Delta }x_n(t),v_n(t)):={\displaystyle \frac{F(v_n(t))v_n(t)}{(\alpha 1)v_n(t)}}\left[\mathrm{\Delta }x_n(t)lv_n(t)\right]`$ (56)
$`(\mathrm{for}v_n(t)\mathrm{\Delta }x_n(t)l\alpha v_n(t)).`$
The parameter $`\alpha `$ determines the range within which the deceleration map $`G(\mathrm{\Delta }x,v)`$ is used. For headways less than $`\alpha v_n(t)`$ the map $`G`$ is used instead of $`F(v)`$. Note that $`G(\mathrm{\Delta }x,v=\mathrm{\Delta }xl)=\mathrm{\Delta }xl`$ and $`G(\mathrm{\Delta }x,v=(\mathrm{\Delta }xl)/\alpha )=F(v=(\mathrm{\Delta }xl)/\alpha )`$, i.e. $`G`$ interpolates between the free-motion map $`F`$ and the sudden braking map of model A. The full velocity map of Model B is thus given by
$$Map_n(v_n(t),\mathrm{\Delta }x_n(t))=\{\begin{array}{cc}F(v_n(t),v_n^F)\hfill & \text{for }\alpha v_n(t)\mathrm{\Delta }x_n(t),\hfill \\ G(\mathrm{\Delta }x_n(t),v_n(t))\hfill & \text{for }v_n(t)\mathrm{\Delta }x_n(t)\alpha v_n(t)),\hfill \\ B(\mathrm{\Delta }x_n(t))\hfill & \text{for }\mathrm{\Delta }x_n(t)v_n(t)).\hfill \end{array}$$
(57)
Local measurements of the flow for a system of vehicles with different preferred velocities $`v_n^F`$ produce a fundamental diagram of inverse-$`\lambda `$ shape (see Fig. 4A) . Here the non-uniqueness of the flow has a simple explanation. Due to the different $`v_n^F`$, platoons form behind the slowest vehicles. Whenever such a platoon passes the measurement region, a flow value on the lower branch is recorded. Otherwise, the flow corresponds to the upper branch.
Measurements of the power spectral density of temporal density fluctuations, i.e. the Fourier transform of the time-series of local densities, show a $`1/f^\alpha `$-behaviour with $`\alpha 1.8`$ in the free-flow regime. Due to the deterministic dynamics, the system evolves into a state with power-law fluctuations. In it has been suggested that the origin of the $`1/f^\alpha `$-fluctuations is the power-law distribution ($`1/(\mathrm{\Delta }x)^{3.0}`$) of the headways $`\mathrm{\Delta }x`$, since these are related to density waves. The occurence of jams destroys long-time correlations since vehicles loose their memory of current fluctuations when they are forced to stop in a jam . Therefore no $`1/f^\alpha `$-behaviour can be observed in the jammed regime.
In a coupled-map model based on optimal-velocity functions has been introduced by discretizing the time variable of the OV model (see Sec. 6.2). This allows to study systems with open boundaries and multilane systems. Furthermore a multiplicative random noise can be imposed in the velocity update so that the velocity map is given by
$$v(t+1)=\left[v(t)+\alpha \left(V_{opt}(\mathrm{\Delta }x)v(t)\right)\right](1+f_{noise}\xi )$$
(58)
where $`\xi [1/2,1/2]`$ is a uniform random variable and $`f_{noise}`$ the noise level.
### 7.2 The model of Nagel and Herrmann
Nagel and Herrmann (NH) have introduced a coupled-map model which is related to the continuum limit of the Nagel-Schreckenberg cellular automata model (see Sec. 8). A generalization of the NH model has later been presented in . Vehicles are characterized by a maximal velocity $`v_{max}`$ and a safety distance $`\alpha `$. The velocity map for the NH model is given by
$`v_n(t+1)=\{\begin{array}{cc}\mathrm{max}(\mathrm{\Delta }x_n(t)\delta ,0)\hfill & \text{for }v_n(t)>\mathrm{\Delta }x_n(t)\alpha ,\hfill \\ \mathrm{min}(v_n(t)+a,v_{max})\hfill & \text{for }v_n(t)<\mathrm{\Delta }x_n(t)\beta .\hfill \end{array}`$ (59)
In the velocity update step, vehicles which have a headway $`\mathrm{\Delta }x`$ smaller than the safety distance $`\alpha `$ decelerate. The headway distance after deceleration is determined by the parameter $`\delta `$. Vehicles which have a large enough headway, on the other hand, accelerate. The acceleration coefficient $`a`$ is determined by $`a=a_{max}\mathrm{max}(1,\mathrm{\Delta }x_n(t)/\gamma )`$.
Since the dynamics of the model is deterministic, the behaviour depends strongly on fluctuations of the initial state . For equidistant starting positions of the vehicles the fundamental diagram consists of two linear branches with maximum flow $`f(c_m)=v_{max}c_m`$ at density $`c_m=1/(v_{max}+\beta )`$. For homogeneous starting positions the system is free-flowing up to a critical density $`c_{crit}`$. Beyond this density free-flowing and congested areas coexist.
### 7.3 The model of Krauss, Wagner and Gawron
Krauss et al. introduced a whole class of stochastic models by considering necessery conditions for the collision-free motion of vehicles. The models are continuous in space and discrete in time.
The vehicles are characterized by a maximum velocity $`v_{max}`$, their acceleration and deceleration capabilities $`a(v)`$ and $`b(v)`$, respectively, and their length $`l`$ which will be taken to be $`l=1`$ in the following. Then the update rules for the velocity $`v`$ and the space coordinate $`x`$ of each vehicle are as follows:
Step 1: Determine desired velocity.
$$v_{des}=\mathrm{min}[v_{max},v+a(v),v_{safe}]$$
Step 2: Randomization.
$$v=\mathrm{max}[0,rand(v_{des}a,v_{des})]$$
Step 3: Vehicle movement.
$$xx+v$$
Here $`rand(v_1,v_2)`$ denotes a random number uniformly distributed in the interval $`[v_1,v_2)`$ and $`v_{safe}`$ is a velocity which guarantees collision motion of the vehicles. It is given explicitly by
$$v_{safe}=v_p+b(\widehat{v})\frac{gv_p}{\widehat{v}+b(\widehat{v})}.$$
(60)
where $`v_p`$ is the velocity of the preceding vehicle located at $`x_p`$ and $`g=x_px1`$ is the headway, i.e. the distance to the preceding vehicle.
In the simplest case the acceleration and deceleration capabilities do not depend on the velocity, i.e. $`a(v)=a=const`$ and $`b(v)=b=const`$. The behaviour of the model can be classified in three different families (see Fig. 11).
The three families of models sketched schematically in Fig. 11 can be characterized as follows:
Class I: High acceleration
Here no spontaneous jamming exists. For $`av_{max}`$ and $`b1`$ the behaviour is similar to that of a cellular automata model without velocity memory introduced in which is closely related to the Kasteleyn model of statistical physics. It can also be interpreted as 5-vertex model . Another model belonging to this class has been introduced by Fukui and Ishibashi (see Sec. 11.1).
Class II: High acceleration – low deceleration
The outflow from a jam is identical to the maximal possible flow. The jamming transition is not a true phase transition, but rather a crossover. The limit $`b\mathrm{}`$, $`a=1`$ corresponds to a continuum version of the Nagel-Schreckenberg cellular automata model which will be introduced in the next section.
Class III: Low acceleration – low deceleration
These models exhibit phase separation and metastability. The jamming transition is of first order. The outflow from a jam is not maximal. For $`av_{max}`$ and $`bv_{max}`$ the model is closely related to the Gipps model discussed in Sec. 6.1. Other models belonging to this class are the Kerner-Konhäuser model (see Sec. 4.4), the optimal-velocity model (Sec. 6.2) and the models with slow-to-start rules which will be introduced in Sec. 9.1.2.
On a macroscopic level, classes I, II and III can be distinguished by the ordering of the densities $`c_f`$ and $`c_c`$, where $`c_f`$ is the density of the outflow from a jam and $`c_c`$ is the density where homogeneous flow becomes unstable . For $`c_c>c_f`$ the outflow from a jam is stable and the system phase-separates into free-flow and jammed regions. Furthermore, metastable states can be found. This is the type of the behaviour found in class III. For $`c_c<c_f`$, on the other hand, the outflow from a jam is unstable and no metastable states or phase separation can be found. This is the typical behaviour of classes I and II. These classes can further be distinguished since in class I one does not find any structure formation, like spontaneous jamming, in contrast to class II.
A related model has been studied before by Migowsky et al. . In this model vehicles are also characterized by a maximum velocity $`v_{max}`$ and a bounded acceleration capability ($`b_{max}<\ddot{x}<a_{max}`$) which determines the safety distance $`d_s`$ necessary to avoid accidents. The investigations in focussed on the effect of so-called driving strategies. These strategies are characterized by a vector $`(f_v,f_a,f_s)`$, where $`f_v`$, $`f_a`$ and $`f_s`$ are the fraction of the vehicle’s maximal velocity, acceleration and safety distance actually used<sup>5</sup><sup>5</sup>5In other words, the maximum velocity is given by $`v_{max}^{(n)}=f_v^{(n)}v_{max}`$ etc., respectively. This can lead to the possibility of accidents and allows to study the number of crashes as a function of the driving strategies.
In different strategies have been compared. Furthermore, dynamical changes of strategies can be introduced which allow the drivers to adapt to the local traffic conditions. In general this leads to a decrease in the number of accidents, jams and fuel consumption, but at high densities the flow is reduced compared to the case of fixed strategies.
## 8 Nagel-Schreckenberg cellular automata model of vehicular traffic on highways
In general, CA are idealization of physical systems in which both space and time are assumed to be discrete and each of the interacting units can have only a finite number of discrete states. Note that for a discretization of differential equations, e.g. those of the hydrodynamic approach, space and time variables are discrete, but the state variable still is continuous. The concept of CA was introduced in nineteen fifties by von Neumann while formulating an abstract theory of self-replicating computing machines . However, it received the attention of a wider audience in the nineteen seventies through Conway’s game of life . The family of one-dimensional CA was studied systematically, in the nineteen eighties, from the point of view of dynamical systems and popularized by Wolfram . Since then the concept of CA has been applied to model a wide variety of systems . To our knowledge, the first CA model for vehicular traffic was introduced by Cremer and Ludwig .
In the CA models of traffic the position, speed, acceleration as well as time are treated as discrete variables. In this approach, a lane is represented by a one-dimensional lattice. Each of the lattice sites represents a ”cell” which can be either empty or occupied by at most one ”vehicle” at a given instant of time (see Fig. 12). At each discrete time step $`tt+1`$, the state of the system is updated following a well defined prescription (a summary of various possible different schemes of updating is given in Appendix A). The computational efficiency of the discrete CA models is the main advantage of this approach over the car-following and coupled-map lattice approaches.
In the NaSch model, the speed $`v`$ of each vehicle can take one of the $`v_{max}+1`$ allowed integer values $`v=0,1,\mathrm{},v_{max}`$. Suppose, $`x_n`$ and $`v_n`$ denote the position and speed, respectively, of the $`n`$-th vehicle. Then, $`d_n=x_{n+1}x_n`$, is the gap in between the $`n`$-th vehicle and the vehicle in front of it at time $`t`$. At each time step $`tt+1`$, the arrangement of the $`N`$ vehicles on a finite lattice of length $`L`$ is updated in parallel according to the following ”rules”:
Step 1: Acceleration. If $`v_n<v_{max}`$, the speed of the $`n`$-th vehicle is increased by one, but $`v_n`$ remains unaltered if $`v_n=v_{max}`$, i.e.
$$v_n\mathrm{min}(v_n+1,v_{max})$$
$`(\mathrm{U1}).`$
Step 2: Deceleration (due to other vehicles). If $`d_nv_n`$, the speed of the $`n`$-th vehicle is reduced to $`d_n1`$, i.e.,
$$v_n\mathrm{min}(v_n,d_n1)$$
$`(\mathrm{U2}).`$
Step 3: Randomization. If $`v_n>0`$, the speed of the $`n`$-th vehicle is decreased randomly by unity with probability $`p`$ but $`v_n`$ does not change if $`v_n=0`$, i.e.,
$$v_n\mathrm{max}(v_n1,0)\mathrm{with}\mathrm{probability}p$$
$`(\mathrm{U3}).`$
Step 4: Vehicle movement. Each vehicle is moved forward according to its new velocity determined in Steps 1–3, i.e.
$$x_nx_n+v_n$$
$`(\mathrm{U4}).`$
The NaSch model is a minimal model in the sense that all the four steps are necessary to reproduce the basic features of real traffic; however, additional rules need to be formulated to capture more complex situations. The step 1 reflects the general tendency of the drivers to drive as fast as possible, if allowed to do so, without crossing the maximum speed limit. The step 2 is intended to avoid collision between the vehicles. The randomization in step 3 takes into account the different behavioural patterns of the individual drivers, especially, nondeterministic acceleration as well as overreaction while slowing down; this is crucially important for the spontaneous formation of traffic jams. Even changing the precise order of the steps of the update rules stated above would change the properties of the model. E.g. after changing the order of steps 2 and 3 there will be no overreactions at braking and thus no spontaneous formation of jams. The NaSch model may be regarded as stochastic CA . In the special case $`v_{max}=1`$ the deterministic limit of the NaSch model is equivalent to the CA rule $`184`$ in Wolfram’s notation and some abstract extensions of this CA-184 rules have been studied in the more general context of complex dynamics and particle flow.
Why should the updating be done in parallel, rather than in random sequential manner, in traffic models like the NaSch model? In contrast to a random sequential update, parallel update can lead to a chain of overreactions. Suppose, a vehicle slows down due the randomization step. If the density of vehicles is large enough this might force the following vehicle also to brake in the deceleration step. In addition, if $`p`$ is larger than zero, it might brake even further in Step 3. Eventually this can lead to the stopping of a vehicle, thus creating a jam. This mechanism of spontaneous jam formation is rather realistic and cannot be modeled by the random sequential update.
The update scheme of the NaSch model is illustrated with a simple example in Fig. 13.
Space-time diagrams showing the time evolutions of the NaSch model demonstrate that no jam is present at sufficiently low densities, but spontaneous fluctuations give rise to traffic jams at higher densities (Fig. 14(a)). From the Fig. 14(b) it should be obvious that the intrinsic stochasticity of the dynamics , arising from non-zero $`p`$, is essential for triggering the jams . For a realistic description of highway traffic , the typical length of each cell should be about $`7.5`$m which is the space occupied by a vehicle in a dense jam. When $`v_{max}=5`$ each time step should correspond to approximately $`1`$ sec of real time which is of the order of the shortest relevant timescale in real traffic, namely the reaction time of the drivers.
Almost all the models of traffic considered in this review, including the NaSch model, have been formulated in such a way that no accident between successive vehicles is possible. However, accident of the vehicles is possible if the condition for safe driving is relaxed. For example, Boccara et al. replaced the update rule of the NaSch model by the rule
$$\mathrm{if}v_{n+1}(t)>0\mathrm{then}x_n(t+1)=x_n(t)+v_n(t+1)+\mathrm{\Delta }v$$
(61)
where $`\mathrm{\Delta }v`$ is a Bernoulli random variable which takes the value $`1`$ with probability $`p_{careless}`$ and zero with the probability $`1p_{careless}`$. The probability $`P_{ac}`$ of accident per vehicle per time step is a non-monotonic function of the vehicle density $`c`$ .
### 8.1 Relation with other models
#### 8.1.1 Relation with totally asymmetric simple exclusion process
Now we point out the similarities and differences between the $`v_{max}=1`$ limit of the NaSch model and the totally asymmetric simple exclusion process (TASEP) which is the simplest prototype model of interacting systems driven far from equilibrium . In the TASEP (Fig. 15) a randomly chosen particle can move forward, by one lattice spacing, with probability $`q`$ if the lattice site immediately in front of it is empty. It corresponds to a Kawasaki dynamics for exchange of a charged particle and hole on nearest-neighbour lattice sites at infinite temperature and in the presence of an infinite electric field (see Appendix B for some technical aspects of TASEP). Several different generalized variants of the TASEP have been considered. For example, in the $`k`$-hop model a particle can exchange its position with the nearest-hole on its right with probability $`q`$, provided the separation of the two sites under consideration is not more than $`k`$ lattice spacings. The $`k`$-hop model reduces to the TASEP in the special case $`k=1`$.
Note that in the NaSch model with $`v_{max}=1`$ every vehicle moves forward with probability $`q=1p`$ in the time step $`t+1`$ if the site immediately in front of it were empty at the time step $`t`$; this, is similar to TASEP. But, updating is done in parallel in the NaSch model whereas that in the TASEP is done in a random sequential manner. Nevertheless, the special case of $`v_{max}=1`$ for the NaSch model achieves special importance from the fact that so far it has been possible to derive exact analytical results for the NaSch model only in the special limits (a) $`p=0`$ and arbitrary $`v_{max}`$ (which we have already considered), and (b) $`v_{max}=1`$ and arbitrary $`p`$.
#### 8.1.2 Relation with surface growth models and the phenomenological fluid-dynamical theories of traffic
The NaSch model with $`v_{max}=1`$ can be mapped onto stochastic growth models of one-dimensional surfaces in a two-dimensional medium, the single-step model . Corresponding to each configuration $`\{\sigma _j\}`$ of the NaSch model in the site-oriented description, one can obtain a unique surface profile $`\{H_j\}`$ through the relation $`H_j=\frac{1}{2}_{jk}(12\sigma _k)`$ . Pictorially one can interpret this mapping as shown in Fig. 16. Particle (hole) movement to the right (left) correspond to local forward growth of the surface via particle deposition. In this scenario a particle evaporation would correspond to a particle (hole) movement to the left (right) which is not allowed in the NaSch model. It is worth pointing out that any quenched disorder in the rate of hopping between two adjacent sites would correspond to columnar quenched disorder in the growth rate for the surface .
The surface growth model described above is known to be a discrete counterpart of continuum models of growing surfaces whose dynamics are governed by the so-called Kardar-Parisi-Zhang equation . Since the Kardar-Parisi-Zhang equation can be mapped onto the Burgers equation using the Cole-Hopf transformation , it is not surprising that several features of vehicular traffic are described by the NaSch model at the microscopic level and by the noisy Burgers equation for the coarse-grained continuum of the fluid-dynamical theory .
### 8.2 Limiting cases of the NaSch model
In spite of the fact that the deterministic limits $`p=0`$ and $`p=1`$ of the NaSch model do not capture some of the most essential features of vehicular traffic it may be instructive to examine these limits to gain insight into the features of this simpler scenario. Another limiting case which exhibits a surprisingly complex behaviour is the case $`v_{max}=\mathrm{}`$.
#### 8.2.1 NaSch model in the deterministic limit $`p=0`$
The NaSch model, a stochastic CA, becomes a deterministic CA in the limit $`p=0`$. In this special case, the deterministic dynamical update rules of the model can be written as
$`v_n(t+1)`$ $`=`$ $`\mathrm{min}[v_{max},v_n(t)+1,d_n(t)1],`$
$`x_n(t+1)`$ $`=`$ $`x_n(t)+v_n(t+1)`$ (62)
which can lead to two types of steady states depending on the density $`c`$ . At low densities, the system can self-organize so that $`d_n>v_{max}`$ for all $`n`$ and, therefore, every vehicle can move with $`v_{max}`$, i.e., $`v_n(t)=v_{max}`$, giving rise to the corresponding flux $`cv_{max}`$. This steady-state is, however, possible only if enough empty cells are available in front of every vehicle, i.e., for $`cc_m^{det}=(v_{max}+1)^1`$ and the corresponding maximum flux is $`J_{max}^{det}=v_{max}/(v_{max}+1)`$. On the other hand, for $`c>c_m^{det}`$, $`d_n(t)1\mathrm{min}[v_n(t)+1,v_{max}]`$ and, therefore, the relevant steady-states are characterized by $`v_n(t)=d_n(t)1`$, i.e., the flow is limited by the density of holes. Since the average distance-headway is $`1/c1`$, the fundamental diagram in the deterministic limit $`p=0`$ of the NaSch model (for any arbitrary $`v_{max}`$) is given by the exact expression
$$J=\mathrm{min}(cv_{max},1c).$$
(63)
Note that the result $`v_n=1/c1`$ is identical with Greenshields ansatz $`v_n=1/c1/c_{jam}`$ if we identify $`c_{jam}=1`$.
#### 8.2.2 NaSch model in the deterministic limit $`p=1`$
Aren’t the properties of the NaSch model with maximum allowed speed $`v_{max}`$, in the deterministic limit $`p=1`$, exactly identical to those of the same model with maximum allowed speed $`v_{max}1`$? Although this expectation may seem to be consistent with the observation that $`J=0`$ for all $`c`$ in the special case $`v_{max}=1=p`$, the answer to the question posed above is: NO. To understand the subtle features of the deterministic limit $`p=1`$ one has to consider $`v_{max}>1`$. You can easily convince yourself that if, for example, $`v_{max}=2`$, then, for $`c>1/3`$, all stationary states correspond to $`J=0`$ because at least one vehicle will have only one empty cell in front (i.e. $`d_n=2`$) and it will never succeed in moving forward. For $`v_{max}=2`$ and $`p=1`$, although there are stationary states corresponding to $`J0`$ for all $`c1/3`$, such states are metastable in the sense that any local external perturbation leads to complete breakdown of the flow. If the initial state is random, such metastable states cannot lead to non-zero $`J`$ because they have a vanishing weight in the thermodynamic limit. Hence, if $`p=1`$, all random initial states lead to $`J=0`$ in the stationary state of the NaSch model irrespective of $`v_{max}`$ and $`c`$!
#### 8.2.3 NaSch model in the limit $`v_{max}=\mathrm{}`$
The limit $`v_{max}=\mathrm{}`$ has been introduced in . One has to be aware that there are several possible ways of performing this limit since only finite systems of length $`L`$ can be treated in computer simulations. In the case $`v_{max}=L`$ has been investigated<sup>6</sup><sup>6</sup>6See also , where the case $`p=0`$ was studied., but other limiting procedures are also possible, e.g. $`v_{max}L^\alpha `$ with $`\alpha >0`$ or even $`v_{max}=\mathrm{}`$ independent of the system size. In principle, these different limiting procedures could lead to different results, but up to now no systematic study has been performed. We therefore restrict ourselves to the case $`v_{max}=L`$ studied in .
Surprisingly one finds that the fundamental diagram has a form quite different from that of the case of a finite $`v_{max}`$. The flow does not vanish in the limit $`c0`$ since already one single car produces a finite value of the flow, $`J(c0)=1`$. Due to the hindrance effect of other cars, $`J(c)`$ is a monotonically decreasing function of the density $`c`$ (see Fig. 17). Another characteristic feature of the fundamental diagram is the existance of a plateau at flow $`J_P`$ where the value $`J_P`$ depends on the randomization $`p`$, but not on the system size $`L`$. The length of the plateau, on the other hand, increases with $`L`$.
What is the microscopic structure of the stationary state leading to such a fundamental diagram ? At low densities, where flow $`J`$ is larger than the plateau value $`J_P`$, the cars tend to be uniformly distributed just as in the deterministic case $`p=0`$ (see Sec. 8.2.1). For densities in the plateau regime, however, one jam exists in the system, whereas for higher densities there is more than one jam. In the thermodynamic limit, one expects a phase transition at $`c=0`$ between a jamless phase with $`J=1`$ and a phase with one jam and flow $`J_P`$ . Increasing the density further, more jams develop and the plateau ceases. Note that this behaviour is completely different from the prediction of mean-field theory in that limit (see Sec. 8.3.1) showing the importance of correlations.
### 8.3 Analytical theories of the NaSch model with periodic boundary conditions
CA are, by design, ideal for large scale computer simulations. However, proper interpretations of the numerical data obtained from computer simulations are not always quite straightforward because of the finite-size effects and ”numerical noise”. One cannot deny the importance of exact analytical results in providing a testing ground for the computer codes. On the other hand, the parallel updating makes exact analytical solution of CA models very difficult. Nevertheless, even in those situations where exact solutions are not possible, a combination of approximate analytical treatments and computer simulation often turns out out to be very powerful method of analysis of a problem. This approach has been quite successful in recent years in the studies of the NaSch model and its generalizations.
Before we proceed with the analytical theories in the non-deterministic NaSch model, we would like to point out that the fundamental diagram $`J(c)`$ is known exactly for arbitrary $`v_{max}`$ and $`p`$ in the two limits $`c0`$ and $`c1`$. In the former case, $`Jcv_F`$ where $`v_F=v_{max}p`$ is the free-flow velocity. On the other hand, in the latter case, $`J(1p)(1c)`$ as flow is determined by holes moving backwards at a speed $`1p`$.
#### 8.3.1 Site-oriented naive mean-field theory for the NaSch model
In the ”site-oriented” theories one describes the state of the finite system of length $`L`$ by completely specifying the state of each site, i.e., by the set $`(\sigma _1,\sigma _2,\mathrm{},\sigma _L)`$ where $`\sigma _j`$ ($`j=1,2,\mathrm{},L`$) can, in principle, take $`v_{max}+2`$ values one of which represents an empty site while the remaining $`v_{max}+1`$ correspond to the $`v_{max}+1`$ possible values of the speed of the vehicle occupying the site $`j`$. In some of the analytical calculations of steady-state properties of the NaSch model one follows, for convenience, the sequence $`2341`$, instead of $`1234`$ of the stages of updating as this merely shifts the starting step and, therefore, does not influence the steady-state properties of the model. The advantage of this new sequence is that, in a site-oriented theory, the variable $`\sigma _j`$ can now take $`v_{max}+1`$ values as none of the vehicles can have a speed $`v=0`$ at the end of the acceleration stage of the updating.
Let us introduce the lattice gas variables $`n(i;t)`$ through the following definition: $`n(i;t)=0`$ if the site labeled by $`i`$ is empty and $`n(i;t)=1`$ if it is occupied by a vehicle (irrespective of the speed). Obviously, the space-average of $`n(i;t)`$ is the density of the vehicles, i.e., $`_in(i;t)/L=c`$. Suppose, $`c_v(i;t)`$ is the probability that there is a vehicle with speed $`v`$ ($`v=0,1,2,\mathrm{},v_{max}`$) at the site $`i`$ at the time step $`t`$. Obviously, $`c(i;t)=_{v=0}^{v_{max}}c_v(i;t)`$ is the probability that the site $`i`$ is occupied by a vehicle at the time step $`t`$ and $`d(i;t)=1c(i;t)`$ is the corresponding probability that the site $`i`$ is empty at the time step $`t`$.
In the naive site-oriented mean-field (SOMF) approximation for the NaSch model one writes down the equations relating $`c_v(i;t+1)`$ ($`v=1,\mathrm{},v_{max}`$) with the corresponding probabilities at time $`t`$ and, then, solves the equations in the steady-state (see the Appendix C for a detailed derivation of these equations for arbitrary $`v_{max}`$). In the simplest case of $`v_{max}=1`$ and periodic boundary conditions one gets
$`c_0(i;t+1)`$ $`=`$ $`c(i;t)c(i+1;t)+pc(i;t)d(i+1;t),`$ (64)
$`c_1(i;t+1)`$ $`=`$ $`qc(i1;t)d(i;t).`$ (65)
The equation (64) expresses the simple fact that at the time step $`t+1`$ the speed of the vehicle at the $`i`$-th site can be zero either because the site $`i+1`$ was occupied at time $`t`$ or because of random deceleration (if the site $`i+1`$ was empty at time $`t`$). Similarly, the equation (65) implies that the speed of the vehicle at the site $`i`$ can be $`1`$ at time $`t+1`$ if at the time step $`t`$ the site $`i`$ was empty while the site $`i1`$ was occupied by a vehicle which did not decelerate during the random deceleration stage of updating.
In the steady state, $`c_v(i,t)`$ are independent of $`t`$. Besides, if periodic boundary conditions are imposed, the $`i`$-dependence of $`c_v(i)`$ also drops out in the translational-invariant steady-state. Therefore, in the steady-state
$$J=c_1=qc(1c)$$
(66)
It turns out that the naive SOMF underestimates the flux for all $`v_{max}`$. Curiously, if instead of parallel updating one uses the random sequential updating, the NaSch model with $`v_{max}=1`$ reduces to the TASEP for which the equation (66) is known to be the exact expression for the corresponding flux (see, e.g., )!
#### 8.3.2 Paradisical mean-field theory of the NaSch model
What are the reasons for these differences arising from parallel updating and random sequential updating? There are ”Garden of Eden” (GoE) states (dynamically forbidden states) of the NaSch model which cannot be reached by the parallel updating whereas no state is dynamically forbidden if the updating is done in a random sequential manner. For example, the configuration shown in Fig. 18 is a GoE state<sup>7</sup><sup>7</sup>7The configuration shown in Fig. 12 is also a GoE state! because it could occur at time $`t`$ only if the two vehicles occupied the same cell simultaneously at time $`t1`$.
The naive SOMF theory, discussed in the preceding subsection, does not exclude the GoE states. On the other hand, results of the paradisical mean-field (pMF) theory are derived by repeating the calculations of the naive SOMF theory excluding all the GoE states from consideration. The exact expression, given in the next subsection, for the flux in the steady-state of the NaSch model with $`v_{max}=1`$ is obtained in the pMF theory (see Appendix D for detailed calculations), thereby indicating that the only source of correlation in this case is the parallel updating . But, for $`v_{max}>1`$, there are other sources of correlation because of which exclusion of the GoE states merely improves the naive SOMF estimate of the flux (Fig. 19) but does not yield exact results .
#### 8.3.3 Site-oriented cluster-theoretic approach to the NaSch Model
The site-oriented cluster theoretic approach leads to a systematic improvement of the naive SOMF theory of the NaSch model. We define a $`n`$-cluster to be a collection of $`n`$ successive sites and denote the probability of finding an $`n`$-cluster in the state $`(\sigma _1,\sigma _2,\mathrm{},\sigma _n)`$ in the steady-state of the system by the symbol $`P_n(\sigma _1,\sigma _2,\mathrm{},\sigma _n)`$. In the general $`n`$-cluster approximation , one divides the lattice into ”clusters” of length $`n`$ such that two neighbouring clusters have $`n1`$ sites in common (see Fig. 20); an $`n`$-cluster is treated exactly and the cluster is coupled to the rest of the system in a self-consistent way, as we shall show in this subsection. Even without any calculation, one would expect that, for a given $`v_{max}`$, the $`n`$-cluster approximation should yield more accurate results with increasing $`n`$ and should give exact results in the limit $`n\mathrm{}`$. Fortunately, in the special case $`v_{max}=1`$ exact results are obtained already for $`n`$ as small as $`2`$, i.e., the results of 2-cluster calculations are exact for $`v_{max}=1`$ .
Let us first explain the key concepts involved in the cluster theory. It is straightforward to verify, for example, in the special case of $`v_{max}=1`$, that the state of the 2-cluster $`\sigma _i,\sigma _{i+1}`$ at time $`t+1`$ depends on the state of the 4-cluster $`(\tau _{i1},\tau _i,\tau _{i+1},\tau _{i+2})`$ at time $`t`$. In general, in the $`n`$-cluster approximation for an arbitrary $`v_{max}`$ one has to take into account the vehicles that can enter an $`n`$-cluster from one of the $`v_{max}`$ cells to its left and can leave it to occupy one of the $`v_{max}`$ cells to its right. Therefore, in general, the state $`\sigma _j,\sigma _{j+1},\mathrm{}\sigma _{j+n1}`$ of an $`n`$-cluster at time $`t+1`$ depends on the state of a $`n+2v_{max}`$ cluster $`\tau _{jv_{max}},\mathrm{}\tau _j,\tau _{j+1},\mathrm{},\tau _{j+n1},\mathrm{},\tau _{j+n1+v_{max}}`$ at time $`t`$. Therefore, in the special case of $`v_{max}=1`$, the master equations
$`P_2(\sigma _i,\sigma _{i+1};t+1)={\displaystyle \underset{\tau _j}{}}W(\sigma _i,\sigma _{i+1}|\tau _{i1},\tau _i,\tau _{i+1},\tau _{i+2})P_4(\tau _{i1},\tau _i,\tau _{i+1},\tau _{i+2};t)`$
(67)
governing the time evolution of the 2-cluster probabilities $`P_2(\sigma _i,\sigma _{i+1})`$ involve the 4-cluster probabilities for all those configurations $`(\tau _{i1},\tau _i,\tau _{i+1},\tau _{i+2};t)`$ which can lead to the 2-cluster configuration $`(\sigma _i,\sigma _{i+1};t+1)`$ under consideration as well as the corresponding transition probabilities $`W(\sigma _i,\sigma _{i+1}|\tau _{i1},\tau _i,\tau _{i+1},\tau _{i+2})`$. Similarly, the master equation governing the time evolution of the 4-cluster probabilities $`P_4(\tau _{i1},\tau _i,\tau _{i+1},\tau _{i+2})`$ involve 6-cluster probabilities, and so on. In order to obtain a closed set of equations one has to truncate this hierarchy in an appropriate manner; in the $`n`$-cluster approximation one expresses the $`(n+2v_{max})`$-cluster probabilities in terms of products of $`n`$-cluster probabilities.
The $`n`$-cluster approximation represented geometrically in Fig. 20 for $`n=1`$ can be expressed mathematically as
$$P(\tau _{j2},\tau _{j1},\tau _j,\tau _{j+1},\tau _{j+2})=\underset{i=\tau _{j2}}{\overset{\tau _{j+2}}{}}P_1(\tau _i)$$
(68)
Thus, 1-cluster approximation is equivalent to the naive SOMF approximation. The $`2`$-cluster approximation represented geometrically in Fig. 20 can be expressed mathematically as
$`P(\tau _{j2},\tau _{j1},\tau _j,\tau _{j+1},\tau _{j+2},\tau _{j+3})`$ $``$ $`P_2(\tau _{j2},\tau _{j1})P_2(\tau _{j1},\tau _j)P_2(\tau _j,\tau _{j+1})`$ (69)
$`P_2(\tau _{j+1},\tau _{j+2})P_2(\tau _{j+2},\tau _{j+3})`$
or, more precisely,
$`P(\tau _{j2},\tau _{j1},\tau _j,\tau _{j+1},\tau _{j+2},\tau _{j+3})`$ $`=`$ $`P_2(\tau _{j2}|\underset{¯}{\tau _{j1}})P_2(\tau _{j1}|\underset{¯}{\tau _j})P_2(\tau _j,\tau _{j+1})`$ (70)
$`P_2(\underset{¯}{\tau _{j+1}}|\tau _{j+2})P_2(\underset{¯}{\tau _{j+2}}|\tau _{j+3})`$
where
$$P_2(\tau _{j1}|\underset{¯}{\tau _j})=\frac{P_2(\tau _{j1},\tau _j)}{_{\tau _{j1}}P_2(\tau _{j1},\tau _j)}$$
(71)
are 2-cluster conditional probabilities. Similarly, the 3-cluster approximation consists of the approximate factorization
$`P(\tau _{j2},\tau _{j1},\tau _j,\tau _{j+1},\tau _{j+2},\tau _{j+3},\tau _{j+4})=P_3(\tau _{j2}|\underset{¯}{\tau _{j1},\tau _j})P_3(\tau _{j1}|\underset{¯}{\tau _j,\tau _{j+1}})`$
$`P_3(\tau _j,\tau _{j+1},\tau _{j+2})P_3(\underset{¯}{\tau _{j+1},\tau _{j+2}}|\tau _{j+3})P_3(\underset{¯}{\tau _{j+2},\tau _{j+3}}|\tau _{j+4}).`$ (72)
Analoguous factorizations hold for an arbitrary number of sites on the left-hand-side of (68),(70) and (72).
Let us now illustrate the scheme of the cluster calculations for the NaSch model by carrying out the calculation for the simplest case, namely, the 2-cluster calculations for $`v_{max}=1`$. For convenience, one follows the sequence $`2341`$, instead of $`1234`$ of the stages of updating so that, for $`v_{max}=1`$, a two-state variable $`\sigma `$ is adequate to describe the state of a lattice site; $`\sigma =0,1`$ correspond, respectively, to an empty site and a site occupied by a vehicle with speed $`1`$. Thus, in the special case of $`v_{max}=1`$ one would need only four $`2`$-cluster probabilities, namely, $`P_2(0,0),P_2(1,0),P_2(0,1),P_2(1,1)`$. Interestingly, the constraints
$$P_2(1,0)+P_2(1,1)=c$$
(73)
and
$$P_2(0,0)+P_2(0,1)=1c$$
(74)
together with the particle-hole symmetry
$$P_2(1,0)=P_2(0,1)$$
(75)
leave only one of the four $`2`$-cluster probabilities, say, $`P_2(1,0)`$, as an independent variable which one needs to calculate by solving the corresponding master equation. For general $`v_{max}`$, the $`n`$-cluster approximation on the right hand side of the master equation leads to $`(v_{max}+1)^n`$ nonlinear equations; the number of independent equations gets reduced by the so-called Kolmogorov consistency conditions .
Using (70) one factorizes the 4-cluster probabilities on the right-hand-side of (67) for $`P_2(1,0)`$ in terms of 2-cluster conditional probabilities. In the first column of the table in Fig. 21 we list all those configurations $`(\tau _{i1},\tau _i,\tau _{i+1},\tau _{i+2};t)`$ which can lead to the configurations, shown in the second column, which is the exhaustive list of the 4-cluster configurations each having $`\sigma _i=1,\sigma _{i+1}=0`$; the corresponding transition probabilities $`W(1,0|\tau _{i1},\tau _i,\tau _{i+1},\tau _{i+2})`$ are given in the third column.
Using the configurations at $`t`$ and $`t+1`$ as well as the corresponding transition probabilities given in table in Fig. 21 the master equation (67) for $`P_2(1,0)`$ reduces to the quadratic algebraic equation
$$qy^2y+c(1c)=0$$
(76)
where we have used the shorthand notation $`y=P_2(1,0)`$. Solving this quadratic equation we get (see also )
$$P_2(1,0)=\frac{1}{2q}\left[1\sqrt{14qc(1c)}\right]$$
(77)
and, hence, $`P_2(1,1),P_2(0,0),P_2(0,1)`$ from the equations (73), (74) and (75). Moreover, the expression (77) establishes that $`P_2(1,0)P_1(1)P_1(0)=c(1c)`$, which indicates an effective particle-hole attraction (particle-particle repulsion) in the NaSch model with $`v_{max}=1`$. Furthermore, from equation (77) one gets the expression
$$J(c,p)=qP_2(1,0)=\frac{1}{2}\left[1\sqrt{14qc(1c)}\right]$$
(78)
which can be proved to be the exact expression for the corresponding flux. It is not difficult to carry out 2-cluster calculations for higher values of $`v_{max}`$, but one gets only approximate results for $`v_{max}>1`$ .
An interesting feature of the expression (78) is that the flux is invariant under charge conjugation, i.e., under the operation $`c(1c)`$ which interchanges particles and holes. Therefore, the fundamental diagram is symmetric about $`c=1/2`$ when $`v_{max}=1`$ (see Fig. 22(a)). Although this symmetry breaks down for all $`v_{max}>1`$ (see Fig. 22(b)), the corresponding fundamental diagrams appear more realistic. Moreover, for given $`p`$, the magnitude of $`c_m`$ decreases with increasing $`v_{max}`$ as the higher is the $`v_{max}`$ the longer is the effective range of interaction of the vehicles (see Fig. 22). Furthermore, for $`v_{max}=1`$, flux merely decreases with increasing $`p`$ (see eqn. (78)), but remains symmetric about $`c=1/2=c_m`$. On the other hand, for all $`v_{max}>1`$, increasing $`p`$ not only leads to smaller flux but also lowers $`c_m`$ (Fig. 23).
For $`v_{max}>1`$, one needs to carry out higher order cluster calculations to get more accurate results than those obtained in the 2-cluster approximation. For $`v_{max}=2`$, the fundamental diagrams obtained from the $`n`$-cluster approximation ($`n=1,2,..,5`$) are compared in Fig. 24 with the Monte Carlo data. This comparison clearly establishes a rapid convergence with increasing $`n`$ and already for $`n=4`$ the difference between the cluster calculation and MC data is extremely small. In the cluster probabilities for $`v_{max}=2`$ have been obtained from computer simulations. The results suggest that the $`n`$-cluster approximation for $`n3`$ becomes asymptotically exact in the limit $`p0`$.
#### 8.3.4 Car-oriented mean-field theory of the NaSch model
In the ”car-oriented” theories the state of the traffic system is described by specifying the positions and speeds of all the $`N`$ vehicles in the system . Suppose, $`P_n(t)`$ is the probability to find at time $`t`$ exactly $`n`$ empty sites immediately in front of a vehicle. Another auxiliary quantity, which turns out to be very convenient to use in several different calculations, is $`g(t)`$, the probability at time $`t`$ that a vehicle moves in the next time step. These two sets of quantities, namely, $`P_n(t)`$ and $`g(t)`$ are related. For example, in the NaSch model with $`v_{max}=1`$, a vehicle will move in the next time step if there is at least one empty cell in front of it (probability $`_{n1}P_n(t)`$) and if it does not decelerate in the randomization step (probability $`q`$); therefore, $`g(t)=q[_{n1}P_n(t)]=q[1P_0(t)]`$. Hence the flux $`J(c,p)`$ can be obtained from $`J(c,p)=cg=cq[1P_0]`$.
The essence of the car-oriented mean-field (COMF) approximation is to neglect the correlations between the gaps in front of the successive cars<sup>8</sup><sup>8</sup>8A similar approach, the so-called interparticle distribution functions technique, is used for studying reaction-diffusion systems .. The equations describing the time evolution of the probabilities $`P_n(t)`$, under this approximation (see Appendix E for the derivation of these equations), can be solved in the steady-state using a generating function technique . Following this approach, in the special case $`v_{max}=1`$, one recovers the exact expression (78) for $`J(c,p)=cq[1P_0]`$.
For $`v_{max}=2`$ one has to distinguish between $`P_n(v=1)`$ and $`P_n(v=2)`$. Moreover, one has to generalize the quantity $`g`$ to $`g_\alpha `$, the probability that the vehicle moves $`\alpha `$ cells ($`\alpha =1,2`$) in the next time step. Applying the same generating function techniques as for $`v_{max}=1`$, one can also solve the coupled sets of steady-state equations for $`P_n(v=1)`$ and $`P_n(v=2)`$ for $`v_{max}=2`$ but gets only approximate results .
Interestingly, finite size of the system affects the equations for $`v_{max}=2`$ in a much more dramatic way than those for $`v_{max}=1`$ thereby revealing the intrinsic qualitative differences in the nature of correlations in the NaSch model for $`v_{max}=1`$ and $`v_{max}>1`$.
Comparisons with Monte Carlo simulations show that in contrast to the 3-cluster approximation for $`v_{max}=2`$ COMF does not become asymptotically exact in the limit $`p0`$. This implies that even in this limit correlations between the headways are not negligible. It is interesting, however, that for the fundamental diagram one finds an excellent agreement between MC simulations and the predictions of COMF for $`p0`$. The reason is that in the deterministic limit many configurations exist which produce the same flow. COMF is not able to identify the dominating structures correctly, but nevertheless can predict the correct current.
#### 8.3.5 Microscopic structure of the stationary state
As we have seen MFT underestimates the flow in the stationary state considerably. Deviations become larger for higher velocities $`v_{max}`$. This shows the importance of correlations. As described above a particle-hole attraction exist. Using the 2-cluster probabilities for $`v_{max}=1`$ this attraction can be expressed in mathematical form as $`P_2(1,0)>P_1(1)P_1(0)=c(1c)`$.
For $`v_{max}=1`$ all improvements of MFT (2-Cluster, COMF and pMFT) are exact. Here only correlations between neighbouring cells are important. The fact that pMFT is exact shows that no ’true’ correlations exist. All correlations have their origin in the existence of GoE states. This also helps to understand why for random-sequential dynamics already MFT is exact for $`v_{max}=1`$ and the stationary state is uncorrelated. The reason is simply that for random-sequential dynamics no GoE states exist!
The situation changes for higher velocities $`v_{max}>1`$. Here pMFT is no longer exact<sup>9</sup><sup>9</sup>9Note that for random-sequential dynamics also MFT is no longer exact!. Therefore ’true’ correlations exist. This corresponds to the observation made in that the NaSch model shows a qualitatively different behaviour for $`v_{max}=1`$ and $`v_{max}>1`$. Furthermore it explains why so far the exact determination of the stationary state for $`v_{max}>1`$ has not been possible.
It is interesting to investigate how the microscopic structure of the stationary state depends on the randomization $`p`$. For $`p=0`$ we have seen in Sec. 8.2.1 that for densities $`c1/(v_{max}+1)`$ the vehicles arrange themselves in such a way that all headways are at least $`v_{max}`$. This is no longer possible for larger densities, but still the vehicles have the tendency to maximize their headway. Furthermore, for $`p=0`$ no spontaneous formation of jams exists since overreactions are not possible. The behaviour in this limit can be interpreted as coming from a kind of ”repulsive interaction” between the vehicles.
The behaviour for $`p=1`$ is a little bit different. Here we have seen in Sec. 8.2.2 that metastable states with finite flow exist for $`c1/3`$ and $`v_{max}>1`$.
For $`0<p<1`$ the microscopic structure interpolates between these two limiting cases. This can be seen by analysing the 3-cluster probabilities obtained from Monte Carlo simulations . For small $`p`$ the microscopic structure of the stationary state is determined by the ’repulsive interactions’ between vehicles. With increasing $`p`$ one finds a tendency towards phase separation into jammed and free-flow regions. A standing vehicle is able to induce a jam even at low densities since the restart probability is small. The jams formed are typically not compact, but of the form ‘.0.0.0.’ since a vehicle approaching the jam slows down in the randomization step with a rather high probability.
Concluding one might say that the microscopic structure for $`0<p<1`$ is determined by the competition of the two ”fixed points” $`p=0`$ and $`p=1`$.
### 8.4 Spatio-temporal organization of vehicles; is there a phase transition?
The density $`c_m`$ corresponding to maximum flux is an obvious first candidate for a critical density separating the regimes of free-flow and congested flow in the NaSch model. We shall show in this subsection that this, indeed, is a critical point provided $`p=0`$. However, in spite of strong indications that, probably, a noise-induced smearing of the transition takes place when $`p0`$, rigorous proofs are still lacking.
$``$ Order parameter
For a proper description of a phase transition one should introduce an appropriate order parameter which can distinguish the two phases because of its different qualitative behavior within the two phases .
A first candidate for the NaSch model would be the average fraction of vehicles at rest, i.e., with instantaneous speed $`v=0`$. In the deterministic limit $`p=0`$ this, indeed, serves the purpose of the order parameter for the sharp transition at $`c_m^{det}`$ from the free-flowing dynamical phase to the congested dynamical phase. But, in the general case of non-zero $`p`$, there is a non-vanishing probability that a vehicle comes to an instantaneous rest merely because of random braking even at extremely low density $`c`$; this probability is $`p`$ for $`v_{max}=1`$ and decreases with increasing $`c`$.
The next obvious choice would be the density of nearest-neighbor pairs in the stationary state
$$m=\frac{1}{T}\frac{1}{L}\underset{t=1}{\overset{T}{}}\underset{j=1}{\overset{L}{}}n_jn_{j+1},$$
(79)
where, as defined earlier, $`n_j=0`$ for an empty cell and $`n_j=1`$ for a cell occupied by a vehicle (irrespective of its velocity). Because of the step 2 of the updating rule (deceleration due to other vehicles) $`m`$ gives the space-time-averaged density of those vehicles with velocity $`0`$ which had to brake due to the next vehicle ahead.
The Fig. 25(a) shows that $`m`$ vanishes at $`c_m^{det}`$ if $`p=0`$. On the other hand if $`p0`$, $`m`$ does not vanish even for $`c<c_m`$ although $`m`$ becomes rather small at small densities (see Fig. 25(b)).
We now present a heuristic argument to point out why any quantity related to the fraction of jammed vehicles is non-zero at any density $`c>0`$ and, hence, inadequate to serve as an order parameter . To slow down to $`v_{max}2`$ a vehicle must be hindered by one randomly braking vehicle in front. Similarly, to reach a speed $`v_{max}3`$ a vehicle must find two vehicles within the range of interaction, and so on. The probability for $`n`$ vehicles to be found in the close vicinity of a given vehicle is proportional to $`c^n`$. Therefore, the probability $`P_v(c)`$ of finding a vehicle with speed $`v<v_{max}1`$ is proportional to $`c^{v_{max}1v}`$ and, hence, even for $`v=0`$, $`P_v(c)`$ is, in general, non-zero for all $`c0`$.
$``$ Spatial correlations
A striking feature of second-order phase transitions is the occurrence of a diverging length scale at criticality and a corresponding algebraic decay of the correlation function . Using lattice gas variables $`n_j`$, the equal-time density-density correlation function is defined by
$$G(r)=\frac{1}{T}\frac{1}{L}\underset{t=1}{\overset{T}{}}\underset{j=1}{\overset{L}{}}n_jn_{j+r}c^2.$$
(80)
which measures the correlations in the density fluctuations that occur at the same time at two different points in space separated by a distance $`r`$.
Again it is very instructive to consider first the deterministic case $`p=0`$ (Fig. 26(a)). Since, as argued before, there are exactly $`v_{max}`$ empty sites in front of each vehicle at $`c=c_m^{det}`$ the correlation function at $`c=c_m^{det}`$ is given by
$$G(r)=\{\begin{array}{cc}c_m^{det}(1c_m^{det})\hfill & \text{for }r0\text{ mod }(v_{max}+1)\text{,}\hfill \\ (c_m^{det})^2\hfill & \text{else.}\hfill \end{array}$$
(81)
For all $`v_{max}`$ the correlation function for small non-zero $`p`$ (Fig. 26(b)) has essentially the same structure as that for $`p=0`$ (Fig. 26(a)) but the amplitude decays exponentially for all $`c`$. In the general case of non-vanishing $`p`$, the asymptotic behaviour ($`r\mathrm{}`$) of the correlation length $`\xi `$ can be obtained analytically only for $`v_{max}=1`$. It turns out that, for given $`p`$, $`\xi `$ is maximum at $`c=1/2=c_m`$ and that $`\xi (c=1/2)p^{1/2}`$. Thus, for $`v_{max}=1`$, $`\xi `$ diverges only for $`p=0`$ but remains finite for all non-zero $`p`$. For $`v_{max}>1`$ the trend of variation of $`\xi `$ with $`c`$ (Fig. 27(a)) in the vicinity of $`c_m`$ is the same as that for $`v_{max}=1`$ . Moreover, for $`v_{max}>1`$, the maximum value of the correlation length, $`\xi _{max}`$ plotted against $`p`$ (Fig. 27(b)), is also consistent with the corresponding trend of variation for $`v_{max}=1`$. Thus, the correlation function $`G(r)`$ gives a strong indication that the NaSch model exhibits a second order phase transition, at $`c=c_m^{det}`$, only for $`p=0`$ but this transition is smeared out if $`p0`$. This noise-induced smearing of the phase transition in the NaSch model is very similar to the smearing of critical phenomena by finite-size effects.
$``$ Distribution of lifetimes of jams
Another quantity which should be able to give information about the nature of the transition from free-flow to the jammed regime is the distribution of lifetimes of jams. Following Nagel each vehicle which has a velocity less than $`v_{max}`$ before the randomization step will be considered jammed. This definition is motivated by the cruise-control limit (see Sec. 9.1.1) where it is more natural than in the NaSch model. One expects, however, that the long-time behaviour of the lifetime distribution is independent of the exact definition of a jam. The short-time behaviour, on the other hand, might differ strongly, e.g. for “compact jams” where a jam is defined as a series of consecutive standing vehicles without any empty cells in between.
Fig. 28 shows the results of Monte Carlo simulations for the lifetime distribution in the NaSch model for different densities near the transition region, $`cc_m=0.085\pm 0.005`$ (for $`v_{max}=5`$, $`p=0.5`$), where $`c_m`$ is the density where the flow is maximal. The most interesting feature of the lifetime distribution is the existence of a cutoff near $`\tau _c=10000`$. It has been shown that this cutoff is neither a finite-size nor a finite-time effect. For times smaller than $`\tau _c`$ a scaling regime exists where the distribution decays algebraically.
$``$ Dissolution of a megajam
Gerwinski and Krug tried to find an intuitive criterion which allows the distinction of free-flow and jammed phases. It is based on the investigation of jam dissolution times. Starting from a megajam configuration, i.e. a block of $`N`$ consecutive cells occupied by vehicles with the remaining $`LN`$ cells being empty, they determined the time until the jam<sup>10</sup><sup>10</sup>10In the same definition of a jam as in the previous point ”Distribution of lifetimes of jams” (see ) has been used. has dissolved completely.
A simple estimate gives the density at which the lifetime is expected to become infinite. Suppose that the jam dissolves with velocity $`v_J`$. Since the first vehicle move freely with average velocity $`v_F=v_{max}p`$ it will reach the end of the jam at the same time as the dissolution wave if the condition $`(LN)/v_F=N/v_J`$ is satisfied. The corresponding density is then given by
$$c^{}=\frac{v_J}{v_J+v_F}=\frac{v_J}{v_J+v_{max}p}.$$
(82)
For $`v_{max}=1`$ vehicles accelerate immediately to $`v_{max}`$. In this case one has $`v_J=q=1p`$. For higher velocities, $`q=1p`$ is only an upper bound for $`v_J`$. Inserting $`v_J=1p`$ into (82) one therefore obtains an upper bound for the density $`c^{}`$. Taking into account interactions between vehicles in the outflow region of the jam, one can derive an effective acceleration rate $`q_{eff}`$, and thus the jam dissolution velocity $`v_J=q_{eff}`$, as a function of $`p`$ .
Computer simulations show a sharp increase of the lifetime near the density $`c^{}`$. It becomes “infinite”, i.e. the jam does not dissolve within the measurement time, only at a higher density $`c_1^{}`$ which is considerably larger than the density $`c_m`$ of maximum flow. At intermediate densities $`c^{}<c<c_1^{}`$ the jam does not dissolve during the first lap, but later due to fluctuations of the two ends of the jam. During this time other jams have usually formed. All results found in are consistent with the measurements of the lifetime distributions presented in the previous point.
$``$ Relaxation time
A characteristic feature of a second order phase transition is the divergence of the relaxation time at the transition point. For $`p=0`$ this has been studied first by Nagel and Herrmann . They found a maximum $`\tau _{max}`$ of the relaxation time at the density $`c=1/(v_{max}+1)`$ which diverges in the thermodynamic limit $`L\mathrm{}`$ as $`\tau _{max}L`$. For finite $`p`$ the behaviour of the relaxation time is more complicated. For technical reasons Csányi and Kertész made no direct measurements of the relaxation time, but used the following approach: Starting from a random configuration of cars with velocity $`v_j=0`$ the average velocity $`\overline{v}(t)`$ is measured at each time step $`t`$. For $`t\mathrm{}`$ the system reaches a stationary state with average velocity $`\overline{v}_{\mathrm{}}`$. The relaxation time is characterised by the parameter
$$\tau =\underset{0}{\overset{\mathrm{}}{}}\left[\mathrm{min}\{v^{}(t),\overline{v}_{\mathrm{}}\}\overline{v}(t)\right]𝑑t.$$
(83)
$`v^{}(t)`$ denotes the average velocity in the acceleration phase $`t0`$ for low vehicle density $`\rho 0`$. In this regime, due to the absence of interactions between the vehicles, one has $`v^{}(t)=(1p)t`$. Thus the relaxation time is obtained by summing up the deviations of the average velocity $`\overline{v}(t)`$ from the values of a system with one single vehicle which can move without interactions with other cars ($`\rho 0`$). Note that for a purely exponentially decaying quantity $`v(t)=v_{\mathrm{}}+C\mathrm{exp}(t/\tau ^{})`$ the definition (83) is proportional to $`\tau ^{}`$, i.e. the standard definition of the relaxation time. One finds a maximum of the relaxation parameter near, but below, the density of maximum flow for $`p=0.25`$ and $`v_{\mathrm{𝑚𝑎𝑥}}=5`$ (see ).
The relaxation time (83) shows interesting behavior which is difficult to interpret in terms of critical point phenomena. One finds a maximum of the relaxation parameter near, but below, the density of maximal flow. This maximum value increases with system size, but the width of the transition region does not seem to shrink . Two scenarios are possible: 1) The relaxation time converges to a large but finite value for large system sizes beyond the present computer power; 2) The relaxation time diverges for $`L\mathrm{}`$. Scenario 1) appears to be more plausible in view of the finite lifetimes of jams discussed above. Complicated interactions between jams could in principle lead to a divergence. Keeping in mind the unusual scaling behaviour of the width, this should occur in a finite interval, not at a special (critical) point.
The interpretation of the parameter $`\tau `$ as a relaxation time can be problematic. This can be seen clearly for $`c>c_{transition}`$, where $`\tau `$ can become negative . Here it is possible that during relaxation the system can temporarily reach states with a higher average velocity than in the stationary state. This overreaction can be divided into two phases for $`p>0`$. Within the first few time steps small clusters which occur in the initial configuration vanish. The second phase is characterized by the growth of surviving jams. More and more cars get trapped into large jams and therefore the average flow decreases to its stationary value. This decrease causes negative values of $`\tau `$ at large densities.
$``$ Distribution of distance-headways
In order to get information on the spatial organization of the vehicles, one can calculate the distance-headway distribution $`𝒫_{dh}(\mathrm{\Delta }x)`$ by following either a site-oriented approach or a car-oriented approach if $`\mathrm{\Delta }x_j=x_{j+1}x_j`$, i.e., if the number $`\mathrm{\Delta }x_j1`$ of empty lattice sites in front of the $`j`$-th vehicle is identified as the corresponding distance-headway.
Stated precisely, $`𝒫_{dh}(k)`$ is the conditional probability of finding a string of $`k`$ empty sites in front of a site which is given to be occupied by a vehicle. A comparison between the naive mean-field expression
$$𝒫_{dh}^{mfa}(j)=c(1c)^j$$
(84)
for the distance-headway distribution in the NaSch model with $`v_{max}=1`$ and the corresponding Monte Carlo data reveals the inadequacy of equation (84) at very short distances which indicates the existence of strong short-range correlations in the NaSch model that are neglected by the mean-field treatment. This is consistent with our earlier observation that there are particle-hole effective short-range attraction in the NaSch model with $`v_{max}=1`$. Again, this correlation disappears when a random sequential updating is carried out! The exact distance-headway distribution in the NaSch model with $`v_{max}=1`$ is found to be
$$𝒫_{dh}^{2c}(j)=\frac{y^2}{c(1c)}\left[1\frac{y}{(1c)}\right]^{j1}(j=1,2,\mathrm{})$$
(85)
where $`y=P_2(1,0)`$ is given by the equation (76).
For all $`v_{max}>1`$, at moderately high densities, $`𝒫_{dh}(\mathrm{\Delta }x)`$ exhibits two peaks, in contrast to a single peak in the distance-headway distributions for $`v_{max}=1`$ at all densities (Fig. 30); the peak at $`\mathrm{\Delta }x=1`$ is caused by the jammed vehicles while that at a larger $`\mathrm{\Delta }x`$ corresponds to the most probable distance-headway in the free-flowing regions. At first sight, the simultaneous existence of free-flowing and jammed regions may appear analogous to the coexistence of gaseous and liquid phases of matter in equilibrium. In fact, when the two-peaked structure of the distance-headway distribution was first observed , it was erroneously interpreted as a manifestation of the coexistence of two dynamical phases, namely, the free-flowing phase and the jammed phase. But, later works established that the analogy between the coexistence of free-flowing and jammed regions in the NaSch model and the coexistence of the gas and liquid phases of matter cannot be pushed too far because the analogue of the gas-liquid interfacial tension is zero in the NaSch model. Thus, one can not conclude that the NaSch model exhibits a first order dynamical phase transition.
$``$ Distributions of jam sizes and gaps between jams
One can identify a string of $`k`$ successive stopped vehicles as a jam of length $`k`$ (by definition, such jams are compact). Similarly, when there are $`k`$ lattice sites between two successive jams, each occupied by a moving vehicle or is vacant then we say that there is a gap of length $`k`$ between the two successive jams. Analytical expressions for the distributions of the jam sizes as well as of the gaps between jams can be calculated for the NaSch model (and some of its extensions) using the 2-cluster approximation or or COMF . The expressions are exact in the case $`v_{max}=1`$ with periodic boundary conditions. For higher velocities the results are only approximative. In COMF the probability $`C_k`$ to find a jam of length $`k`$ is given by
$$C_k^{(COMF)}=(1P_0)P_0^{k1}$$
(86)
whereas in the 2-cluster approach one finds
$`C_k^{(2cl)}`$ $`=`$ $`{\displaystyle \frac{1}{𝒩_J}}P(\underset{¯}{0}|1)P(\underset{¯}{1}|1)^{k2}P(\underset{¯}{1}|1)P(\underset{¯}{1}|0)(k2),`$
$`C_1^{(2cl)}`$ $`=`$ $`{\displaystyle \frac{1}{𝒩_J}}{\displaystyle \underset{v=1}{\overset{v_{max}}{}}}P(\underset{¯}{0}|v)P(\underset{¯}{v}|0).`$ (87)
For the $`n`$cluster approximation similar expressions can be derived.
Both distribitions (86) and (87) decay exponentially for large jam sizes. COMF always predicts a monotonous distribution with $`C_k^{(COMF)}C_{k+1}^{(COMF)}`$. In contrast, the jam size distribution in the $`n`$cluster approximation can in principle exhibit a maximum at small jam sizes $`1kn`$.
$``$ Distribution of time-headways
Since flux is equal to the inverse of the average time-headway, much more detailed information is contained in the full distribution of the time-headway than in the fundamental diagram. The time-headway distribution contains information on the temporal organization of the vehicles.
Suppose, $`𝒫_m(t_1)`$ is the probability that the following vehicle takes time $`t_1`$ to reach the detector, moving from its initial position where it was located when the leading vehicle just left the detector site. Suppose, after reaching the detector site, the following vehicle waits there for $`\tau t_1`$ time steps, either because of the presence of another vehicle in front of it or because of its own random braking; the probability for this event is denoted by $`Q(\tau t_1|t_1)`$. The distribution $`𝒫_{th}(\tau )`$, of the time-headway $`\tau `$, can be obtained from
$$𝒫_{th}(\tau )=\underset{t_1=1}{\overset{\tau 1}{}}𝒫_m(t_1)Q(\tau t_1|t_1)$$
(88)
Substituting the expressions for $`𝒫_m(t_1)`$ and $`Q(\tau t_1|t_1)`$ for $`v_{max}=1`$ in (88) we, finally, get
$`𝒫^{th}(t)={\displaystyle \frac{qy}{cy}}\left(1{\displaystyle \frac{qy}{c}}\right)^{t1}+{\displaystyle \frac{qy}{dy}}\left(1{\displaystyle \frac{qy}{d}}\right)^{t1}`$
$`\left[{\displaystyle \frac{qy}{cy}}+{\displaystyle \frac{qy}{dy}}\right]p^{t1}q^2(t1)p^{t2}.`$ (89)
where, $`q=1p`$, $`d=1c`$ and, for the given $`c`$ and $`p`$, $`y`$ can be obtained from equation (77). The expression is plotted in Fig. 31(a) for a few typical values of $`c`$ for a given $`p`$. A few typical time-headway distributions in the NaSch model for $`v_{max}>1`$, obtained through computer simulation, are shown in Fig. 31(b).
$``$ Temporal correlations
In order to probe the spatio-temporal correlations in the fluctuations of the occupation of the cells, one can study the space-time correlation function
$$C(r,\tau )=\frac{1}{T}\frac{1}{L}\underset{j=1}{\overset{L}{}}\underset{t=1}{\overset{T}{}}n_i(t)n_{j+r}(t+\tau )c^2$$
(90)
which, by definition, vanishes in the absence of any correlation. In three different regimes have been distinguished.
Free-flow ($`0<\rho \rho _1`$): “Minijams” occur which resolve immediately. The correlation function shows anticorrelations around propagating peak.
Jamming ($`\rho _1<\rho \rho _2`$): Free-flow and jamming coexist, i.e. jams with a finite lifetime and vehicles moving with $`v_{max}`$ occur. This coexistence is reflected in the behaviour of the correlation function which exhibits a double-peak structure.
Superjamming ($`\rho _2<\rho 1`$): The whole system is congested. Jamming waves are connected and form an infinite wave. As a consequence, the propagating peak in the correlation function has disappeared.
Neubert et al. have introduced a special autocorrelation function of the density in order to study the velocity of jams. They have determined jam velocities for several variants of the NaSch model which will be introducted in later sections.
$``$ Structure factor
Structure factors are known to give valuable information about driven systems . For the NaSch model the static structure factor
$$S(k)=\frac{1}{L}\left|\underset{j=1}{\overset{L}{}}n_je^{ikj}\right|$$
(91)
has been investigated in . Again $`n_j`$ denotes the occupation number of cell $`j`$. Note that $`S(k)`$ is related to the Fourier transform of the density-density correlation function (80).
For all densities, $`S(k)`$ exhibits a maximum at $`k_00.72`$ which corresponds to the characteristic wavelength $`\lambda _0=2\pi /k_0`$ of the density fluctuations in the free-flow regime. For $`v_{max}>1`$ one finds $`k_0(v_{max}+1)=const`$.
In these investigations have been extended to the dynamical structure factor in velocity-space,
$$S_v(k,\omega )=\frac{1}{NT}\left|\underset{n,t}{}v_n(t)e^{i(kn\omega t)}\right|^2,$$
(92)
with $`k=2\pi m_k/N`$, $`\omega =2\pi m_\omega /T`$, where $`N`$ is the number of vehicles and $`m_k`$ and $`m_\omega `$ are integers. $`v_n(t)`$ is the velocity of the $`n`$-th vehicle at time $`t`$.
Compared to the dynamical structure factor in real space, (92) has the advantage that the free-flow regime only contributes white noise, $`S_v(k,\omega )|_{freeflow}=const`$. Therefore it is easier to study jamming properties. It is found in that $`S_v(k,\omega )`$ exhibits one ridge with negative slope, corresponding to backward moving jams. One finds that the velocity of the jams is a function of the randomization parameter $`p`$ only. It is independent of the density $`c`$ and the maximal velocity $`v_{max}`$ . This is consistent with results from a direct study of the autocorrelation function <sup>11</sup><sup>11</sup>11Measurements of the jam dissolution speed in , however, show a decrease with increasing $`v_{max}`$ and saturation for large $`v_{max}`$.. Above a transition density, an algebraic behavior $`S_v(k,\omega )|_{\omega /k=v_\mathrm{j}}k^\gamma `$ of the structure factor is found. This has been interpreted as an indication of critical behavior in . However, due to the difficulties involved in the calculation of (92) only relatively short times $`T2048`$ have been considered in . This is much smaller than the cutoff found in the lifetime of jams (see the discussion above) and lies well in the region where an algebraic decay is found. Therefore the results for the dynamical structure factor (92) and the lifetime measurements are consistent, but the algebraic decay is not to be interpreted as an evidence for the existence of a critical point in the NaSch model. In order to see the cutoff, times $`T>10^4`$ would have to be considered.
### 8.5 Exact solution of the NaSch model with $`v_{max}=1`$ and open boundary conditions
The analytical methods presented in Section 8.3 are well suited for the investigation of translationally-invariant stationary states which are achieved by imposing periodic boundary conditions. For both practical and theoretical reasons sometimes different boundary conditions are preferable. Imagine a situation where a multilane road is reduced to one lane, e.g. due to road construction. Such a bottleneck can be modeled by using a NaSch model with open boundaries. The multilane part of the road acts as a particle reservoir. If the first cell of the one-lane part is empty a car is inserted here with probability $`\alpha `$. At the other end a car is removed from the last cell with probability $`\beta `$ (see Fig. 32). These boundary conditions break the translational invariance and in general one can expect stationary states with a non-trivial density profile $`\tau _j`$. From a more theoretical point of view such models have been studied intensively as prototypes of systems exhibiting so-called boundary-induced phase transitions . In contrast to what one expects from experience with equilibrium systems, one-dimensional driven nonequilibrium systems can exhibit phase transitions, even when the interactions are short-ranged, just by ’slightly’ changing the boundary conditions.
The probability distribution characterizing the steady state of the TASEP with parallel dynamics (i.e., NaSch model with $`v_{max}=1`$) and open boundary conditions has been obtained recently in and using generalizations of techniques based on the Matrix Product Ansatz (MPA) (see Appendix F for a more technical introduction). By varying the boundary rates $`\alpha `$ and $`\beta `$ one obtains a surprisingly rich phase diagram (see Fig. 33) which is qualitatively the same for all types of dynamics. Three phases can be distinguished by the functional dependence of the current through the system on the system parameters. In the low-density phase A ($`\alpha <\beta ,\alpha _c(p)`$) the current is independent of $`\beta `$. Here the current is limited by the rate $`\alpha `$ which then dominates the behaviour of the system. In the high-density phase B ($`\beta <\alpha ,\beta _c(p)`$) the behaviour is dominated by the output rate $`\beta `$ and the current is independent of $`\alpha `$. In the maximum current phase C ($`\alpha >\alpha _c(p)`$ and $`\beta >\beta _c(p)`$) the limiting factor for the current is the bulk rate<sup>12</sup><sup>12</sup>12Note that conventionally the hopping rate in the ASEP is denoted as $`p`$. Since in the NaSch model $`p`$ is the braking probability the hopping rate in the ASEP (for $`v_{max}=1`$) becomes $`q=1p`$. $`q=1p`$. Here the current becomes independent of both $`\alpha `$ and $`\beta `$.
High- and low-density phase can be subdivided into two phases AI, AII and BI and BII, respectively. These subphases can be distinguished by the asymptotic behaviour of the density profiles at the boundaries.
The transitions between the phases can be characterized by the behaviour of two correlation lengths $`\xi _\alpha `$ and $`\xi _\beta `$ which only depend on $`p`$ and $`\alpha `$ or $`\beta `$. These lengths can be obtained explicitly from the exact solution. Apart from $`\xi _\alpha `$ and $`\xi _\beta `$ also a third length $`\xi ^1=|\xi _\alpha ^1\xi _\beta ^1|`$ plays an important role.
The transition from AII (BII) to C is continuous with diverging correlation length $`\xi _\alpha `$ ($`\xi _\beta `$). The transition from the high- to the low-density phase is of first order. Here both $`\xi _\alpha `$ and $`\xi _\beta `$ are finite, but $`\xi `$ diverges. On the transition line one finds a linear density profile created by the diffusion of a domain wall between a low-density region at the left end of the chain and a high-density region at the right end.
For the case of parallel dynamics, i.e. the NaSch model with $`v_{max}=1`$, the currents in the three phases are given by
$$J_A=\frac{\alpha (q\alpha )}{q\alpha ^2},J_B=\frac{\beta (q\beta )}{q\beta ^2},J_C=\frac{1}{2}\left(1\sqrt{p}\right).$$
(93)
The corresponding bulk densities<sup>13</sup><sup>13</sup>13In the maximum current phase no real bulk density can be defined due to the algebraic behaviour of the density profile. $`c_C`$ is therefore just $`\tau _{L/2}`$. are
$$c_A=\frac{\alpha (1\alpha )}{q\alpha ^2},c_B=\frac{q\beta }{q\beta ^2},c_C=\frac{1}{2}.$$
(94)
The phase boundaries are determined by the critical rates
$$\alpha _c(p)=\beta _c(p)=1\sqrt{p}.$$
(95)
In Fig. 33 also the special line $`(1\alpha )(1\beta )=p`$ is indicated. Here the density profile is flat (i.e. constant). On this line the exact solution can be obtained by the 2-cluster approach of Sec. 8.3.3 . Since it goes through all three phase these results are sufficient to obtain exact analytic expressions, e.g. for the currents, once the structure of the phase diagram is established (e.g. by Monte Carlo simulations).
The stationary state of the ASEP can also be obtained for other types of updates (see App. A), e.g. random-sequential , ordered-sequential and sublattice-parallel update . One finds that the phase diagram has the same basic structure for all updates . The functional dependence of the currents, density profiles etc. on the model parameters differs, however. For the important case of random-sequential updating (93)-(95) have to be replaced by
$`J_A`$ $`=q\alpha (1\alpha ),`$ $`J_B`$ $`=q\beta (1\beta ),`$ $`J_C`$ $`={\displaystyle \frac{q}{4}},`$ (96)
$`\rho _A`$ $`=\alpha ,`$ $`\rho _B`$ $`=1\beta ,`$ $`\rho _C`$ $`={\displaystyle \frac{1}{2}},`$ (97)
$`\alpha _c(p)`$ $`=\beta _c(p)={\displaystyle \frac{q}{2}}.`$ (98)
Results for other updates can be found in . For a discussion of the calculation of diffision constants and shock profiles we refer to the reviews and references therein.
In the behaviour of the ASEP for $`\beta =1`$ was explained by postulating a maximal-current principle. According to this principle, independent of the details of the dynamics, the system tries to maximize the stationary current $`J`$:
$$J=\underset{c[0,c_{}]}{\mathrm{max}}J(c)$$
(99)
Here $`J(c)`$ is the fundamental diagram (for periodic boundary conditions) and $`c_{}`$ is the density at the left (input) boundary, i.e. $`c_{}=\alpha `$ in the case described above.
In a nice physical picture has been developed which explains the structure of the phase diagram not only qualitatively, but also (at least partially) quantitatively. It is determined by the dynamics of a domain walls<sup>14</sup><sup>14</sup>14A somewhat related approach has been used to obtain an approximate solution for the special case of parallel dynamics with deterministic bulk dynamics ($`p=0`$) .. In nonequilibrium systems, a domain wall is an object connecting two possible stationary states. The notion of domain walls in the ASEP can be illustrated in the limit $`\alpha L1`$ and $`\beta L1`$ of small boundary rates. At late times there will be a low-density region at the left end of the chain and a high-density region at the right end, with a domain wall in between. Schematically this state can be depicted as $`000011111`$. For general values of the rates the wall not be sharp in general, but spread over a few lattice sites.
For late times the dynamics of the system can then be interpreted in terms of the motion of the domain wall. A particle entering the system leads moves the wall one cell to left, and a particle leaving the system moves it one cell to the right. Therefore the domain wall performs a biased random walk with drift velocity $`v_D=\beta \alpha `$ and diffusion coefficient $`D=(\alpha +\beta )/2`$. For $`\alpha <\beta `$ the domain wall moves to the right until it reaches the end of the system which is thereafter in the low-density stationary state. For $`\alpha >\beta `$ the wall moves to the left until it reaches the left end and the systems goes into the high-density stationary state. In the case $`\alpha =\beta `$ there is no net drift in the position of the wall. It fluctuates with its rms displacement increasing with time as $`(Dt)^{1/2}`$, i.e. it can be anywhere in the system resulting in a linear density profile.
In order to understand the case of general $`\alpha `$ and $`\beta `$ one has to introduce a second kind of domain wall separating a maximum current phase from the high-density phase. Since the maximal possible flow for periodic boundary conditions is $`J_{max}=1/4`$ (for $`p=0`$ and random-sequential update) the dynamics for $`\alpha =1/2`$ is dominated by the overfeeding at the left boundary. The injection rate could support a current larger than $`1/4`$, but in the bulk it can not exceed this value. Therefore at the left boundary a maximum current state is formed. If the particles are not extracted fast enough at the right boundary a high-density region will develop there. These two regions are separated by a new kind of domain wall, the maximum current/high-density domain wall. Schematically it can be represented as mmmm1111. Again this domain wall performs a biased random walk.
In order to obtain more quantitative predictions one goes to a coarse-grained picture. Then it is useful to replace the boundary rates $`\alpha `$ and $`\beta `$ by particle reservoirs with densities $`c_{}`$ and $`c_+`$. The continuity equation $`c/t+J/x=0`$ in the continuum limit has traveling wave solutions of the form $`c(xv_Dt)`$ with the domain wall velocity
$$v_D=\frac{J_+J_{}}{c_+c_{}}$$
(100)
which can be obtained by integration over the chain.
For the low-density/high-density domain wall one has<sup>15</sup><sup>15</sup>15We assume $`p=0`$ and random-sequential dynamics. $`c_+=1\beta `$, $`J_+=J(c_+)=\beta (1\beta )`$ and $`c_{}=\alpha `$, $`J_{}=J(c_{})=\alpha (1\alpha )`$. This gives indeed $`v_D=\beta \alpha `$ which should be valid for $`\alpha ,\beta <1/2`$. For the maximum current/high-density domain wall $`c_{}`$ takes the value $`c_{}=1/2`$ so that $`J_{}=1/4`$ and thus $`v_D=\beta 1/2`$.
The arguments described above can be generalized to any process where the fundamental diagram $`J(c)`$ of the periodic system has only one maximum at a density $`c^{}`$. For all currents $`J<J(c^{})`$ there exist two corresponding densities $`c_1`$ and $`c_2`$ with $`J(c_1)=J=J(c_2)`$. For fundamental diagrams with more than one maximum, more than two densities might exist for a given current $`J`$. This implies the existence of a larger number of domain wall types. The phase diagram of the open system than exhibits a larger number of phases . The maximal-current principle (99) for the TASEP with $`\beta =1`$ is generalized to the extremal-current principle
$$J=\{\begin{array}{cc}\mathrm{max}_{c[c_+,c_{}]}J(c)\hfill & \text{for }c_{}>c_+,\hfill \\ \mathrm{min}_{c[c_{},c_+]}J(c)\hfill & \text{for }c_{}<c_+.\hfill \end{array}$$
(101)
Since the above phenomenological picture does not depend on the microscopic details of the dynamics, it plausible that the phase diagrams for different updates are qualitatively the same. Boundary-induced phase transitions have recently been observed in measurements on a German motorway . One finds a first-order nonequilibrium phase transition between a free-flow and a congested phase. This transition is induced by the interplay between density waves induced by an on-ramp and a shock wave moving on the motorway .
## 9 Generalizations and extensions of the NaSch model
As stated earlier, the NaSch model is a minimal model. The first obvious possible generalization would be to replace the acceleration stage of updating rule (U1) to
$$v_n\mathrm{min}(v_n+a_n,v_{max})$$
$`(\mathrm{U1}^{})`$
where, $`a_n`$, acceleration assigned to the $`n`$-th vehicle, need not be unity and, in general, may depend on $`n`$. In the following subsections we consider more non-trivial generalizations and extensions of the NaSch model.
### 9.1 Single-lane highways
In the next few subsubsections we shall demonstrate the rich variety of traffic phenomena that can be observed by appropriate modifications of the random braking. We have earlier mentioned in the context of empirical results that traffic flow exhibits metastability and the related hysteresis effects. Such phenomena have been observed in continuum formulations of ”microscopic” models, i.e., in coupled-map lattice models . However, the NaSch model is too simple to account for these phenomena. We now briefly describe a few generalizations of the NaSch model, each of which is based on modifications of the braking rules of the original NaSch model; one common feature of all of these generalized models is that they show metastability and hysteresis.
Before we begin our discussions on specific generalized versions of the NaSch model, which exhibit metastability, we make some general remarks. In the schematic stationary fundamental diagram of Fig. 34, the low density branch corresponds to homogeneous free-flow states, while the high density branch corresponds to configurations, where jammed states are present. Obviously, at densities $`c_1<c<c_2`$, the flow depends non-uniquely on the global density.
In order to establish the existence of meta-stable states one can follow two basic strategies. In the first method, the density of vehicles is changed adiabatically by adding or removing vehicles from the stationary state at a certain density. Starting in the jamming phase (large densities) and removing vehicles, one obtains the lower branch of the hysteresis curve. On the other hand, by adding vehicles to a free flowing state (low densities), one obtains the upper branch. This method is closely related to the experimental situation, where the occupancy of the road varies continuously.
The second method does not require changing the density. Instead one starts from two different initial conditions, the mega-jam and the homogeneous state. The mega-jam consists of one large, compact cluster of standing vehicles. In the homogeneous state, vehicles are distributed periodically with same constant gap between successive vehicles (with one larger gap for incommensurate densities). Then, for $`c>c_1`$ the homogeneous initialization leads to a free-flow state, while the mega-jam initialization leads to the jammed high-density states.
#### 9.1.1 Cruise-control limit and self-organized criticality
In the cruise-control limit of the NaSch model vehicles moving with their desired velocity $`v_{max}`$ are not subject to noise. This is exactly the effect of a cruise-control which automatically keeps the velocity constant at a desired value. In this model the acceleration, deceleration (due to other vehicles) and movement stages of updating are identical to those in the general case of the NaSch model; however, the randomization step is applied only to vehicles which have a velocity $`v<v_{max}`$ after step 2 of the update rule. We can express this more formally by recasting the randomization stage of the update rules in the NaSch model as follows:
$$v_n\mathrm{max}(0,v1)$$
with probability
$$p=\{\begin{array}{cc}p_{v_{max}}\hfill & \text{if }v=v_{max}\text{,}\hfill \\ p\hfill & \text{if }v<v_{max}\text{,}\hfill \end{array}$$
(102)
where $`v`$ is the velocity of the vehicle at the end of the step 2 of the update rule, i.e., after deceleration due to blocking by other vehicles. In the original formulation of the NaSch model $`p_{v_{max}}=p`$. On the other hand, the cruise-control limit corresponds to $`p_{v_{max}}0,p0`$.
For $`p_{v_{max}}1`$, at sufficiently low densities, all the vehicles move deterministically with the velocity $`v_{max}`$; this deterministic motion is, however, interrupted by small perturbations at a vanishingly small rate. Consequently, the system gets enough time to relax back to the state corresponding to the deterministic algorithm before it is perturbed again. This effectively separates completely the time scales for perturbing the system and the response of the system.
First, let us consider the periodic boundary conditions which is easier to treat than than the open boundary conditions. In this model, a sharp transition from the ”free-flowing” dynamical phase to the ”congested” phase takes place at a critical concentration $`c_{}(p,v_{max})`$ which depends on $`p`$ as well as $`v_{max}`$ and, for all $`p0`$, $`c_{}(p,v_{max})`$ is smaller than $`c_m^{det}=1/(v_{max}+1)`$. For a given $`v_{max}`$, $`c_{}(p)`$ increases with decreasing $`p`$ and, in the deterministic limit $`p0`$, $`c_{}(0,v_{max})c_m^{det}`$.
In this model a jam is defined to consist of vehicles all of which have their instantaneous velocities smaller than $`v_{max}`$. For all $`c<c_{}`$, jams present in the initial configuration eventually disappear and in the jam-free stationary state every vehicle moves with the velocity $`v_{max}`$. Therefore, in the density regime $`c<c_{}`$ the flux increases linearly with density following $`J=cv_{max}`$, just like that in the deterministic limit $`p=0`$ of the NaSch model (Fig. 35). But, unlike the deterministic limit $`p=0`$, the cruise-control limit of the NaSch model exhibits metastability in the interval $`c_{}<c<c_m^{det}`$. In this context, the metastability means that, in the interval $`c_{}<c<c_m^{det}`$, on appropriate initialization, the system can reach apparently steady states where no jam appears and where the fluxes are higher than $`J(c_{})`$; but, perturbations of such a ”metastable” state creates long-lived jams thereby reducing the flux to a level consistent with the stable branch of the fundamental diagram. At all $`c>c_{}`$ jams present in the initial configuration never disappear completely and, in this density regime, the stable steady-state is a mixture of laminar flow regions and jams. The long-lived jams lower the flux beyond $`c_{}`$ and the flux decreases linearly with density (Fig. 35).
Let us assume that at densities slightly above $`c_{}`$, only one jam of length $`L_{jam}`$ containing $`N_{jam}`$ vehicles exists in the system. Then, because of the periodic boundary conditions, the total number of vehicles $`N`$ is conserved and, hence,
$$N=c_{jam}L_{jam}+c_{out}(LL_{jam})$$
(103)
where $`c_{jam}=N_{jam}/L_{jam}`$ and $`c_{out}=(NN_{jam})/(LL_{jam})`$ are the densities of the vehicles in the jam and in the outflow region, respectively. Dividing both sides of (103) by $`L`$ we get
$$c=c_{jam}\frac{L_{jam}}{L}+c_{out}\left(1\frac{L_{jam}}{L}\right)$$
(104)
Since in the cruise-control limit of the NaSch model $`L_{jam}`$ must vanish as $`cc_{}`$, we conclude that we must have $`c_{out}=c_{}`$, i.e., the average density in the outflow region of a jam is equal to the critical density $`c_{}`$.
In order to study the traffic at the critical point of the cruise-control limit, Nagel and Paczuski used a a special boundary condition which enables the system to select automatically the state of maximum throughput, i.e, the system exhibits self-organized criticality. This special boundary condition consists of an infinite jam from $`\mathrm{}`$ to $`0`$ (i.e., at the left boundary) while the right boundary is open. Vehicles emerge from the infinite jam in a jerky way, before attaining the velocity $`v_{max}`$. Far away from the infinite jam all vehicles move with the same velocity $`v_{max}`$. In order to show that the state selected this way is ”critical” we perturb a vehicle, far downstream from the infinite jam, slightly by reducing its velocity from $`v_{max}`$ to $`v_{max}1`$. This particular vehicle initiates a chain reaction and gives rise to a jam if the following vehicle is sufficiently close to it although it itself accelerates and, eventually, attains $`v_{max}`$. This phantom jam has a time-dependent size $`n(t)`$, measured by the number of vehicles $`n`$ in this jam at time $`t`$ and it has a lifetime $`\tau _{life}`$. The statistics of this features of the phantom jam can be obtained by repeating the computer experiment sufficiently large number of times; sometimes the phantom jam created is small and has a short lifetime and sometimes it is large and has quite long lifetime. Interestingly, the characteristic quantities like, for example, the distributions of the sizes of the jams, lifetimes of the jams, etc. do, indeed, exhibit power-laws which are hall mark of the self-organized criticality . E.g. the branching behaviour of the jams gives rise to intermittent dynamics with a $`1/f`$ power law spectrum . $`1/f`$ noise in real traffic has been discovered by Musha and Higuchi . They recorded transit times of vehicles passing underneath a bridge. The corresponding power spectral density of the flow fluctuations shows $`1/f`$ behaviour at low frequencies.
The exponents associated with the various power laws in the cruise-control limit of the NaSch model can be calculated analytically, at least for $`v_{max}=1`$, by utilizing a formal relation with one-dimensional unbiased random walk. If $`v_{max}=1`$, all the vehicles in the jams have velocity $`v=0`$. Moreover, the jams are compact so that the number of vehicles in a jam is identical to its spatial extent. The probability distribution $`P(n;t)`$ for the number of vehicles $`n`$ in such a jam, at time $`t`$, is determined by the following equation:
$$P(n;t+1)=(1r_{in}r_{out})P(n;t)+r_{in}P(n1;t)+r_{out}P(n+1;t)$$
(105)
where the phenomenological parameters $`r_{in}`$ and $`r_{out}`$ are the rates of incoming and outgoing vehicles. Of course, $`r_{in}`$ depends on the density of the vehicles behind the jam. For large $`n`$ and $`t`$, taking the continuum limit of the equation (105) and expanding we get
$$\frac{P}{t}=(r_{out}r_{in})\frac{P}{n}+\frac{r_{in}+r_{out}}{2}\frac{^2P}{n^2}$$
(106)
If $`r_{in}>r_{out}`$ the jams would grow for ever. On the other hand, the jams would shrink, and eventually disappear, if $`r_{in}<r_{out}`$. If $`r_{in}=r_{out}`$, the first term on the right hand side of the equation (106) vanishes and the resulting equation governing the time evolution of $`P(n;t)`$ is identical to that of the probability of finding, at time $`t`$, an unbiased one-dimensional random walker at a distance $`n`$ which was initially at the origin. Thus, when $`r_{in}=r_{out}`$, the jams exhibit large (”critical”) fluctuations which can be characterized by critical exponents. Using this formal mapping onto unbiased random walk, we find (a) that the mean size of jam at time $`t`$ corresponds to the mean displacement of the random walker from the origin after time interval $`t`$ and (b) that the lifetime of a jam corresponds to the time taken by the random walker to return to the origin for the first time. Hence, using the well known results from the theory of random walks , we get
$$n(t)t^{1/2},\text{and}P(\tau _{life})\tau _{life}^{3/2}.$$
(107)
It turns out that the power-law exhibited by the size of the jams, the distributions of the lifetimes, etc. are not restricted merely to the special case $`v_{max}=1`$ of the cruise-control limit but is also shown by the corresponding computer simulation data also for arbitrary $`v_{max}`$. The power-law distributions of $`P(\tau _{life})`$ in the cruise-control limit of the NaSch model is in sharp contrast with the exponential distribution observed in the NaSch model . Thus, in the cruise-control limit of the NaSch model, the large jams are fractal in the sense that there are smaller sub-jams inside larger jams, ad infinitum. In other words, in between sub-jams, there are holes of all sizes.
#### 9.1.2 Slow-to-start rules, metastability and hysteresis
The slow-to-start rules can lead not only to metastability and, consequently, hysteresis, but also to phase separated states at high densities, as we now show.
$``$ Takayasu-Takayasu slow-to-start rule
Takayasu and Takayasu (TT) were the first to suggest a CA model with a slow-to-start rule. Here, we investigate the generalization, as suggested in , of the original slow-to-start rule. According to this generalized version, a standing vehicle (i.e., a vehicle with the instantaneous velocity $`v=0`$) with exactly one empty cell in front accelerates with probability $`q_t=1p_t`$, while all other vehicles accelerate deterministically. The other steps of the update rule (U2-U4) of the NaSch model are kept unchanged.
As in the case of the NaSch model, it is instructive to consider first the deterministic limits of the TT model . The TT model reduces to the NaSch model in the limit $`p_t=0`$. What happens in the other deterministic limit, namely, $`p_t=1`$? In the latter deterministic limit, a stopped vehicle can move only if there are at least two empty cells in front . Obviously, completely blocked states exist for densities $`c0.5`$, where the number of empty cells in front of each vehicle is smaller than two. However, in the region $`0.5c0.66`$ the number of blocked configurations is very small compared to the total number of configurations and states with a finite flow exist. Precisely at $`c=0.5`$, there are only two blocked states and the time to reach these states diverges exponentially with the system size.
The fundamental diagram for the TT model with $`v_{max}=1`$ has been derived analytically by carrying out (approximate) 2-cluster calculations in the site-oriented approach . But, the fundamental diagrams of the TT model for all $`v_{max}>1`$ have been obtained so far only numerically by carrying out computer simulations (see Fig. 36).
Comparing these fundamental diagrams with the corresponding ones for the NaSch model ($`p_t=0`$), we find the following effects of the TT slow-to-start rule: (i) for a given density $`c`$, the flux $`J(c)`$ is smaller in the TT model as compared to that in the NaSch model; (ii) the particle-hole symmetry is not exhibited by the TT model for any $`v_{max}`$ (not even for $`v_{max}=1`$) and (iii) the TT model exhibits metastability and hysteresis which are absent in the NaSch model. Note that the mechanism for meta-stability in the case $`p_t=1`$ is different from that for the metastability observed for $`0<p_t<1`$.
Because of the slow-to-start rules, the separations between the vehicles coming out of a jam are larger than those between the vehicles coming out a jam in the NaSch model. Since the density far downstream is smaller than the density of maximum flow, the vehicles can propagate freely in the low density regions of the lattice where spontaneous formation of jams is highly unlikely, if the parameter $`p`$ is sufficiently small. Therefore, the phase-separated steady-states at high global densities consist of a macroscopic jam and a macroscopic free-flow regime both of which coexist simultaneously (Fig. 37).
$``$ The BJH model of slow-to-start rule
Benjamin, Johnson and Hui (BJH) modified the updating rules of the NaSch model by introducing an extra step where their ”slow-to-start” rule is implemented; this slow-to-start rule is different from that introduced by TT . According the BJH ”slow-to-start” rule , the vehicles which had to brake due to the next vehicle ahead will move on the next opportunity only with probability $`1p_s`$. The steps of the update rules can be stated as follows:
Step 1: Acceleration. $`v_n\mathrm{min}(v_n+1,v_{max})`$.
Step 2: Slow-to-start rule: If $`flag=1`$, then $`v_n0`$ with probability $`p_s`$.
Step 3: Blockage (due to other vehicles). $`v_n\mathrm{min}(v_n,d_n1)`$ and, then,
$`flag=1`$ if $`v_n=0`$, else $`flag=0`$.
Step 4: Randomization. $`v_n\mathrm{max}(v_n1,0)`$ with probability $`p`$.
Step 5: Vehicle movement. $`x_nx_n+v_n`$.
Here $`flag`$ is a label distinguishing vehicles which have to obey the slow-to-start rule ($`flag=1`$) from those which do not have to ($`flag=0`$).
Obviously, for $`p_s=0`$ the above rules reduce to those of the NaSch model. The slow-to-start rule of the TT model is a ‘spatial’ rule. In contrast, the BJH slow-to-start rule requires ‘memory’, i.e. it is a ‘temporal’ rule depending on the number of trials and not on the free space available in front of the vehicle. The fundamental diagram of the BJH model (Fig. 38) clearly shows the existence of metastable states and, therefore, expected to exhibit hysteresis effects . But, in the special case of $`v_{max}=1`$, for which approximate analytical calculations can be carried out , no meta-stable states exist.
Since for all $`v_{max}>1`$ in the BJH model, just as in the TT model, the outflow from a jam is smaller than the maximal flow, the phase-separated steady-states at high global densities consist of a macroscopic jam and a macroscopic free-flow regime both of which coexist simultaneously (Fig. 39) .
However, the macroscopic jam is not compact. The typical size of the macroscopic free-flow regime can be estimated by measuring the distribution of the gaps between the successive jams . A peak occurs in this distribution for headways of the order of the system size (see the inset of the right part of Fig. 40). The position of the peak indicates the typical size of the macroscopic free-flow regime.
$``$ The NaSch model with a velocity-dependent slow-to-start rule
Although the NaSch model does not exhibit metastable states and hysteresis, a simple generalization exists which is able to reproduce these effects. It is the so-called Velocity-Dependent-Randomization (VDR) model . Here, in contrast to the original NaSch model, the randomization parameter depends on the velocity of the vehicle, $`p=p(v)`$. The rules (see Sec. 8) are supplemented by a new rule,
Step 0: Determination of the randomization parameter. The randomization parameter used in step 3 for the $`n`$-th vehicle is given by $`p=p(v_n(t))`$.
This new step has to be carried out before the acceleration step 1. The randomization parameter used in step 3 depends on the velocity $`v_n(t)`$ of the $`n`$th vehicle after the previous timestep. In order to implement a simple slow-to-start rule one chooses
$$p(v)=\{\begin{array}{cc}p_0\hfill & \text{for }v=0\text{,}\hfill \\ p\hfill & \text{for }v>0\text{,}\hfill \end{array}$$
(108)
with $`p_0>p`$. This means that vehicles which have been standing in the previous timestep have a higher probability $`p_0`$ of braking in the randomization step than moving vehicles.
The rules of the VDR model can be recast in a form similar to those of the BJH model. We define a label $`flag`$ which distinguishes between vehicles which have to obey the slow-to-start rule ($`flag=1`$) from those which do not have to ($`flag=0`$). $`flag=1`$ if $`v_n=0`$ at the beginning of a time step, else $`flag=0`$. Explicitly, the update rules are as follows:
Step 1: Acceleration. $`v_n\mathrm{min}(v_n+1,v_{max})`$.
Step 2: Blockage (due to other vehicles). $`v_n\mathrm{min}(v_n,d_n1)`$,
Step 3: Randomization. $`v_n\mathrm{max}(v_n1,0)`$ with probability $`p_0`$ if $`flag=1`$. else, $`v_n\mathrm{max}(v_n1,0)`$ with probability $`p`$.
Step 4: Vehicle movement. $`x_nx_n+v_n`$.
Let us compare this VDR model with the cruise-control limit of the NaSch model. The vehicles with velocity $`v=v_{max}`$ (at the end of the step 2) are treated deterministically in the cruise-control limit whereas in the VDR model velocities of all those with the velocity $`v>0`$ (just before the step 3) are updated stochastically, but using different values of the braking parameter.
Typical fundamental diagrams look like the one shown in Fig. 41 where, over a certain interval of $`c`$, $`J(c)`$ can take one of the two values depending on the initial state and, therefore, exhibit metastability.
Moreover, typical space-time diagrams of the VDR model (see Fig.42) clearly demonstrate that metastable homogeneous states have a lifetime after which their decay leads to a phase separated steady state. The microscopic structure of these phase-separated high-density states is qualitative similar to those observed in the high-density regimes of the TT and BJH models but differs drastically from those found in the NaSch model.
It is instructive to compare the fundamental diagram of the VDR model with those of the corresponding NaSch models. We now present a simple derivation of the fundamental diagram of the VDR model on the basis of heuristic arguments utilizing the observed structures of the steady-states. For small densities $`c1`$ there are no slow vehicles in the VDR model since interactions between vehicles are extremely rare. In this regime every vehicle can move with the free-flow velocity $`v_f=(1p)v_{max}+p(v_{max}1)=v_{max}p`$ and, therefore, the flux is given by
$$J_{hom}(c)=c(v_{max}p)$$
(109)
which is identical to the NaSch model with randomization $`p`$. On the other hand, for densities close to $`c=1`$, the vehicles are likely to have velocities $`v=0`$ or $`v=1`$ only and, therefore, the random braking is dominated by $`p_0`$, rather than $`p`$, while the flow is determined by the movement of the holes. Hence, for large densities, i.e., $`1c1`$, the flow is given by $`J(c)(1p_0)(1c)`$ which corresponds to the NaSch model with randomization $`p_0`$. This expression for flux in the high-density regime can also be derived as follows. In the phase-separated state the vehicles are expected to move with the velocity $`v_f=v_{max}p`$ in the free-flow region. Neglecting interactions between vehicles in the free-flowing region (which is justified because of the corresponding low density), the average distance of two consecutive vehicles in the free-flow region is given by $`\mathrm{\Delta }x=c_f^1=T_wv_f+1`$ where the average waiting time $`T_w`$ of the first vehicle at the head of the megajam is given by $`T_w=\frac{1}{1p_0}`$. In other words, the density in the free-flow regime $`c_f`$ is determined by the average waiting time $`T_w`$ and $`v_f`$. Now suppose that $`N_J`$ and $`N_F`$ are the number of vehicles in the megajam and free-flowing regions, respectively. Using the normalization $`L=N_J+N_F\mathrm{\Delta }x`$ we find that for the density $`c=\frac{N_F+N_J}{L}`$, the flux $`J_{sep}(c)`$ is given by $`J_{sep}(c)=\frac{N_F}{L}(v_{max}p)`$ and, hence,
$$J_{sep}(c)=(1p_0)(1c).$$
(110)
Obviously, $`c_f`$ is precisely the lower branching density $`c_1`$, because for densities below $`c_f`$ the jam-length is zero. It should be noted that the heuristic arguments presented above remain valid for $`p_0p`$ and $`v_{max}>1`$. The condition $`p1`$ guarantees that the jams are compact in that limit. In the case $`v_{max}=1`$, vehicles can stop spontaneously, even in the free-flow regime and these vehicles might initiate a jam. This is the basic reason why hysteresis is usually not observed for $`v_{max}=1`$.
Analogous to the BJH model phase separation can be directly identified using the results of the jam-gap distribution. Fig. 43 shows that the size of the free-flow regime is proportional to the system size.
The results for the slow-to-start models discussed above have been obtained by computer simulations of periodic systems of finite length. It was shown that the fundamental diagram which is sketched in Fig. 34 is generic for all models under consideration. Now it is self-evident to ask what kind of stationary states are realised in the thermodynamic limit $`L\mathrm{}`$. The simulation results indicate that $`\mathrm{\Delta }c=c_1c_2`$ decreases with larger system sizes and is expected to vanish for $`L\mathrm{}`$, i.e. the jammed branch is stable in that limit. This is readily understood if one analyses the typical configurations which lead to an emerging jam or, vice versa, the mechanism of the dissolution of a jam. Jams emerge if overreactions of drivers lead to a chain reaction. This is possible in dense regions of the free-flow state where the gap between the vehicles is not larger than $`v_{max}`$. Obviously the probability to find large platoons of vehicles driving with small spatial headways increases with the system size (for fixed density). In addition to that the jammed states are phase separated, i.e. the size of the jam is of the order of the system size. During a simulation run the size of the jam fluctuates due to the stochastic movement and acceleration of the vehicles. Jams can dissolve if the amplitude of these fluctuations are of the order of the length of the jam, which is impossible in the thermodynamic limit.
Therefore the non-unique behaviour of the fundamental diagram is only observable if finite system sizes are considered or if the vehicles move deterministically in the free flow regime. Nevertheless the results discussed above are highly relevant for practical purposes, because the hysteresis effects have been observed at realistic system sizes (e.g. $`L=10000`$ corresponds to a highway of length $`75km`$).
#### 9.1.3 Flow-optimization and meta-stable states
Hysteresis effects and meta-stable states are not only of theoretical interest, but also have interesting applications. From the previous discussion of the slow-to-start models it is evident, that one can optimize the maximum flow, if the homogeneous state is stabilized by controlling the density so that it never exceeds $`c_2`$. This strategy was followed in minimizing frequent jams in the Lincoln- and the Holland-Tunnels in New York. Before the traffic lights were installed at the entrance of the tunnels jams used to form spontaneously within the tunnel because (a) the vehicle density used to be sufficiently high and (b) the drivers used to drive more carefully inside the tunnel thereby giving rise to stronger fluctuations which caused the jams. But, the traffic lights installed at the entrance of the tunnels do not allow the density to exceed $`c_2`$ and, consequently, jams are not formed spontaneously by the decay of any metastable high-density state.
One can mimic the situation of Lincoln and Holland-tunnels within the framework of the CA models in the following way . The tunnel is considered as part of the road, where the braking probabilities $`p^t,p_0^t`$ are higher compared to the remaining part of the lattice ($`p,p_0`$, see Fig. 44). Therefore, if one allows for an uncontrolled inflow of the vehicles, jams typically appear inside the “tunnel” and the system capacity is governed by $`p_0^t`$.
The situation differs drastically if traffic lights are implemented . As shown in Fig. 45, a considerable increase of the maximum capacity can be achieved for an optimal combination of the red-/green-signal periods. The gain of capacity obtained for the optimal intervals of the signal is of the same order as for the realistic examples .
### 9.2 Multi-lane highways
For a realistic description of traffic on highways the idealized single-lane models must be generalized to develop CA models of multi-lane traffic; the main ingredient required for this generalization being the lane-changing rules. Several attempts have been made so far in this direction . The lane changing rules for two-lane traffic can be symmetric or asymmetric with respect to the lanes. Similarly, if there are two (or more) different types of vehicles (say, cars and trucks) with two different $`v_{max}`$, the lane-changing rule can be symmetric or asymmetric with respect to the vehicles.
In general, the update in the two-lane models is divided into two sub-steps: in one sub-step, the vehicles may change lanes in parallel following the lane-changing rules and in the other sub-step each vehicle may move forward effectively as in the single-lane NaSch model. Drivers must find some incentive in changing the lane. Two obvious incentives are (a) the situation on the other lane is more convenient for driving, and (b) the need to make a turn in near future. Two general prerequisites have to be fulfilled in order to initiate a lane change: first, there must be an incentive and second, the safety rules must be fulfilled . Lane changing rules according to this scheme have been introduced by Rickert et al. . They suggested that vehicles are allowed to change the lane if the following four criteria are satisfied:
(C1) $`gap(i)<l`$,
(C2) $`gap_o(i)>l_o`$,
(C3) $`gap_{o,back}(i)>l_{o,back}`$,
(C4) $`rand()<p_c`$.
Here $`gap(i)`$ and $`gap_o(i)`$ are the gaps in front of vehicle $`i`$ on the own lane and the other lane<sup>16</sup><sup>16</sup>16The gap on the other lane is defined in the same way as the gap on the own lane by imagining that the vehicle occupies the site parallel to its current position., respectively. $`gap_{o,back}(i)`$ is the gap on the other lane to the next vehicle behind. $`l`$, $`l_o`$, $`l_{o,back}`$ and $`p_c`$ are parameters specifying the rule and $`rand()`$ is a random number in the interval $`[0,1]`$.
The first rule C1 represents the incentive criterion, i.e. if the gap $`gap(i)`$ in front of the vehicle is not sufficiently large vehicles want to change the lane. Typical choices of the parameter $`l`$ are given by $`l=\mathrm{min}(v+1,v_{max})`$. This choice of the minimal headway ensures that vehicles driving in a slow platoon try to change the lane if possible. In the next rule C2 it is checked if the situation on the other lane is indeed more convenient. This motivates the choice $`l=l_o`$. The third rule C3 avoids too small distances to following vehicles on the other lane. Rickert and coworkers suggested $`l_{o,back}=v_{max}`$. It is also important to perform lane changing stochastically. Even if the incentive and safety criteria are fulfilled a lane change is performed only with probability $`p_c`$ (C4). This avoids, at least partially, so called ping-pong lane changes, i.e. multiple lane-changes of vehicles in consecutive timesteps<sup>17</sup><sup>17</sup>17This artifact of the parallel update was already pointed out by Nagatani , who simulated a two-lane system with $`v_{max}=1`$.. Already implementations of the NaSch model using the basic lane-changing rules revealed quite realistic results. Nevertheless several variants of the basic rules have been developed in order to improve the realism. A large number of lane changing rules considered in the literature have been tabulated and compared by Nagel et al. (see Fig. 46).
The lane changing rules for two-lane traffic can be symmetric or asymmetric with respect to the lanes . If symmetric lane changing rules are applied the rules do not depend on the direction of the lane changing maneuver. In contrast also asymmetric lane changing rules have been considered. Lane changing rules can be asymmetric in two ways. First it is possible that it is preferred to drive on the right lane at low densities. This behavior can be implemented simply by leaving out the first rule for a change from the left to the right lane. Second it is also possible that it is even forbidden to overtake a vehicle on the right lane, e.g. on german highways. Then the single lane dynamics on the right lane depends on the configuration on the left lane. These examples show the flexibility of the CA approaches. Moreover the simulations also show that the details of the lane changing rules may lead to considerable changes of the model results .
In multi-lane traffic it is of particular interest to investigate systems with different types of vehicles. For CA models this has been done first by Chowdhury et al. who simulated a periodic two-lane system with slow and fast vehicles, i.e. vehicles with different $`v_{max}`$. The simulation results have been shown that already for small densities the fast vehicles take on the average the free-flow velocity of the slow vehicles, even if only a small fraction of slow vehicles have been considered. Analogous results have been obtained by Helbing and Huberman who used a different CA model for the in-lane update (see Sec. 11.3 for the definition of the model). In addition to that Nagel et al. have been shown that for a suitable choice of the lane changing rules and different types of vehicles even the phenomenon of ”lane inversion” which has been observed at german highways can be reproduced<sup>18</sup><sup>18</sup>18On german highways the left lane is considered for overtaking vehicles only. Therefore, at low densities, the right lane is used more often. Surprisingly, at higher densities not simply a balancing of the lane usage has been observed, but for densities close to the optimum flow the left lane is even higher frequented.. The results discussed so far show the strong influence of slow vehicles in multi-lane systems. They fit fairly well the empirical results, which show an alignment of the speeds on different lanes and and of different types of vehicles. Nevertheless recent simulation results of Knospe et al. indicate that the influence of slow vehicles seems to be overestimated by the multi-lane variants of the NaSch model. In particular for symmetric lane changing rules even a single slow vehicle can dominate the dynamics close to the optimal value of the flow. In order to weaken the effect of slow vehicles they suggested to consider anticipation effects, i.e. the driver estimates the velocity of the vehicle in the next timestep .
Another interesting quantity to look at is the frequency of lane changes at different densities. Here the simulation results show that close to the density of maximal flow the number of lane changing maneuvers drastically decreases if the small values of the braking noise are considered in CA models where the velocities of vehicles are solely determined by the distance to the vehicle ahead. This is due to the fact that for homogeneous states at high densities no sufficiently large gaps exist. For larger values of the braking noise large density fluctuations are observable. Therefore the local minimum of the lane-changing frequency is not found for larger values of $`p`$.
In general the simulation results show that some generic multi-lane effects can be pointed out. First of all the maximal performance of multi-lane systems is slightly increased compared to corresponding single-lane network. In addition, slow vehicles lead to an alignment of velocities of different type of vehicles already at low densities which is confirmed by empirical observation. This effect is quite robust for different choices of the CA model as well as for different lane changing rules. It can be weakened most efficiently if anticipation effects are applied. The details of the lane changing rules, however, may have strong influence on the lane usage characteristics.
### 9.3 Bidirectional traffic
Simon and Gutowitz have introduced a two-lane CA model where the vehicles move in opposite directions. Passing may be allowed on one or on both lanes. It is only attempted if there is a chance to complete the pass. Therefore drivers measure the local density, i.e. the density of vehicles in front that have to be passed. If it is sufficiently low, a pass will be attempted. This means that at high global densities the lanes are effectively decoupled since only very few passes will occur.
In principle, three types of jams can occur on a bidirectional road: 1) Spontaneous jamming and start-and-stop waves on one of the lanes; 2) jams caused by drivers who try to pass but can not return to their home lane since there is not enough space and 3) “super jams” when an adjacent pair of drivers tries to pass simultaneously. These super jams halt traffic on both lanes and can be prevented by breaking the symmetry between the lanes.
The precise rules of the CA are in the same spirit as the rules for multilane traffic described in the previous subsection 9.2. First, the situation on the own lane is examined. If the motion is hindered by another vehicle (moving in the same direction), a pass is attempted. This will only be initiated if the safety criteria are satisfied: 1) The gap on the other lane has to be sufficiently large, and 2) the number of vehicles to be passed has to be small. Even if these criteria are satisfied a lane change occurs only with probability $`p_{change}`$. After this lane changing step the vehicles move forward similar to the dynamics of the NaSch model. There are, however, important differences. Passing vehicles never decelerate randomly. In order to break the symmetry between the two lanes moving vehicles which are on their home lane and see oncoming traffic decelerate deterministically by one unit. This rule prevents the occurance of a super jam.
The results of show the expected behaviour, namely that passing makes traffic more fluid. Start-stop waves are surpressed if the density is not too large. The improvement of the flow on one lane compared to the the one-lane model depends on the density of vehicles on the other lane. It is maximal for very small densities ($`c0`$) on the passing lane. If the density on the other lane is small ($`c<0.25`$) the flow may be lower than in the one-lane model since passing oncoming vehicles create an additional hindrance. For large densities on at least on of the lanes there is little difference between the one- and two-lane models.
Lee et al. have proposed a toy model for bidirectional traffic based on a multispecies generalization of the ASEP. Here no passing is allowed. Instead oncoming traffic on the opposite lane reduces the hopping rates of the vehicles. The dynamics on each lane is given by that of the ASEP with random-sequential update and $`v_{max}=1`$, but the hopping rate from an occupied cell $`j`$ to an empty cell $`j+1`$ on lane 1 depends on the occupancy of cell $`j+1`$ on the opposite lane (lane 2). When this cell is empty, vehicles hop with rate 1, otherwise with rate $`1/\beta `$. On lane 2 vehicles move in the opposite direction and the hopping rate from cell $`j+1`$ to cell $`j`$ depends on the occupancy of cell $`j`$ on lane 1. It is given by $`\gamma `$ when this cell is empty and by $`\gamma /\beta `$ if it is occupied.
For $`\gamma <1`$ the uninfluenced hopping rate on lane 2 is smaller than that of lane 1. The vehicles on lane 2 might therefore be interpreted as trucks. The interlane interaction parameter $`\beta `$ can be interpreted as a kind of road narrowness. For $`\beta =1`$ vehicles are not slowed down by oncoming traffic. This corresponds to a highway with divider. The case $`\beta 0`$ corresponds to a narrow road being completely blocked by the oncoming traffic.
The behaviour of the model with only one truck is rather similar to that of the NaSch model with quenched disorder (see Sec. 10). For $`\beta >\beta _c`$ the system segregates into two phases, a high-density phase in front of the truck and a low-density phase behind it.
By forbidding trucks and cars to occupy parallel cell $`j`$ simultaneously the model can be mapped onto an exactly solvable 2-species variant of the ASEP. Using the matrix-product Ansatz (see App. F) many steady-state properties for the single-truck case can be obtained exactly. Two phase can be distinguished: A low-density phase for $`c\beta <1`$ and a jammed phase for $`c\beta >1`$ where $`c`$ is the density of vehicles on lane 1. In contrast to the case of a fixed defect site (see Sec. 10.3) only one critical density $`c_{crit}=1/\beta `$ exists since the particle-hole symmetry is broken.
Generalizations of this model to other updates and higher velocities can be found in .
## 10 Effects of quenched disorder on traffic
### 10.1 Randomness in the braking probability of drivers and Bose-Einstein-like condensation
We have seen how modifications of the random braking probability or the rule(s) for random braking in the NaSch model can give rise to a rich variety of physical phenomena, e.g., self-organized criticality, metastability and hysteresis, etc. Now we consider the effects of quenched randomness in the random braking probability $`p`$, i.e., we study the effects of assigning randomly different time-independent braking probabilities $`p_i`$ to different drivers $`i`$ in the NaSch model. Such ”quenched” (i.e., time-independent) randomness in the random braking in the NaSch model can lead to exotic phenomena which are reminiscent of ”Bose-Einstein-like condensation” in the TASEP where particle-hopping rates are quenched random variables . Various aspects of these phenomena have been thoroughly reviewed by Krug and, therefore, we’ll restrict our discussion to only the essential points.
Let us first consider the special case of the NaSch model with $`v_{max}=1`$. As explained earlier, this model reduces to the TASEP if the parallel updating is replaced by random sequential updating scheme. If the same hopping probability $`q`$ is assigned to every particle except one for which the the hopping probability is $`q^{}<q`$, then the single ”impurity” particle is the slowest moving one. The faster particles can be allowed to overtake the slow one at a non-zero rate ; however, if this rate of overtaking vanishes the slow particle will give rise to a platoon of particles behind it. This phenomenon is very similar to the formation of platoons of vehicles in a traffic behind the slow vehicles (e.g., trucks).
Here we are interested in a more general situation of quenched ”disorder” in the form of a distribution of intrinsic hopping probabilities of the vehicles in the system rather than that of the single ”defect” particle. In such situations random initial conditions can lead to the formation of platoons if (a) slow particles are sufficiently rare and (b) if the density of vehicles is sufficiently low. Following their formation, starting from a random initial condition, the platoons grow through coalescence. The coarsening of the platoons has been investigated in the same spirit in which coarsening of domains (the so-called Oswald ripening) is monitored while studying spinodal decomposition in, for example, binary alloys . Suppose, $`\xi (t)`$ is the typical platoon size at time $`t`$. Starting from a homogeneous spatial distribution of the vehicles, $`\xi (t)`$ can be monitored as a function of time $`t`$ to find out the law of ”growth” of the size of the platoons.
Before describing the effects of the quenched randomness in the hopping probabilities on the steady-states of the TASEP and the NaSch model, we consider an even simpler model of platoon formation which was developed using the language of aggregation phenomena. In this model an initial velocity $`v_j`$ is assigned to each vehicle $`j`$, drawn randomly from a continuous probability density $`f(v)`$. The particles then move ballistically along a line and coalesce whenever a faster vehicle catches up with a slower one in front. It has been found that $`\xi (t)`$ increases indefinitely according to the power law
$$\xi (t)t^{(n+1)/(n+2)}$$
(111)
where the exponent $`n`$ characterizes the behaviour of $`f(v)`$ in the vicinity of the minimal velocity $`v_{min}`$, i.e., $`f(v)A(vv_{min})^n`$ as $`vv_{min}`$ with some positive constant $`A`$. An attempt has been made to develop a coarse-grained description of this phenomenon .
It has been shown that if quenched random hopping probabilities are assigned to each particle in the TASEP, there are small gaps between particles in the high-density congested phase but in the inhomogeneous low-density phase there is a macroscopically large empty region in front of the slowest particle (i.e., the particle with smallest hopping probability) behind which a platoon is formed. The phase transition from the low-density inhomogeneous phase (which consists of a macroscopic free region and a platoon) to the high-density congested phase is, in many respects, analogous to the Bose-Einstein transition.
In order to see this analogy, let us imagine that the empty sites are bosons and the state of a boson is determined by which particle it is immediately in front of. In the language of the ideal Bose gas, in the high-density phase the bosons are thinly spread over all the states. On the other hand, in the low-density phase there is a finite fraction of the empty sites are condensed in front of the slowest particle in such ”Bose-Einstein-like condensed” state. The steady-state velocity of the particles is the analogue of the fugacity of the ideal Bose gas. What makes the system interesting is the fact that the platoon appears at low-density rather than at high density of the vehicles.
The Bose-Einstein-like-condensation in the TASEP with quenched random hopping probabilities of the individual particles survives when the random sequential updating is replaced by parallel updating . Finally, it is worth emphasizing that, the analogy with the ideal Bose gas is only formal as the empty sites in the TASEP are not non-interacting quantum particles.
The qualitative features of the dynamical phases and phase transitions observed in the NaSch model with random braking probabilities, for $`v_{max}=1`$ as well as for larger $`v_{max}`$, are very similar to those described above for the TASEP with random hopping probabilities . Typical snapshots of the system at three different stages of evolution from a random initial state are shown in Fig. 47.
The typical size of the platoons $`\xi (t)`$ can be computed directly by computing the correlation function (80) and identifying the separation $`r=R_0`$ of the first zero-crossing of this correlation as $`\xi (t)`$. Following this procedure, it has been observed that $`\xi (t)`$ follows the power law (111) when the distribution of the random braking probabilities is given by $`P(p)=2^n(n+1)(\frac{1}{2}p)^n`$.
### 10.2 Random $`v_{max}`$
The two important parameters of the NaSch model are $`p`$ and $`v_{max}`$. In the preceding subsection we have seen the effects of randomizing $`p`$ assigning the same $`v_{max}`$ to all the vehicles. In this subsection, on the other hand, we investigate the effects of randomizing $`v_{max}`$, assigning a non-random constant $`p`$ to every driver.
The simplest possible model to investigate the effects of quenched randomness in $`v_{max}`$ is that considered by Ben-Naim et al. which was discussed to motivate the phenomenon of platoon formation in the preceding subsection.
In order to model traffic consisting of two different types of vehicles, say, for example, cars and trucks, of which a fraction $`f_{fast}`$ are intrinsically fast (say, cars) while the remaining fraction $`1f_{fast}`$ are intrinsically slow (say, trucks), Chowdhury et al. assigned a higher $`v_{max}`$ (e.g., $`v_{max}=5`$) to a fraction $`f_{fast}`$ of vehicles chosen randomly while the remaining fraction $`1f_{fast}`$ were assigned a lower $`v_{max}`$ (e.g., $`v_{max}=3`$). As the density of the vehicles increases, the vehicles with higher $`v_{max}`$ find it more difficult to change lane in order pass a vehicle with lower $`v_{max}`$ ahead of it in the same lane. This leads to the formation of ”coherent moving blocks” of vehicles each of which is led by a vehicle of lower $`v_{max}`$ . Two main causes of traffic accidents, namely, differences in vehicles speeds and lane changes, are reduced considerably in this state thereby making this state of traffic much safer. It is worth mentioning that even a small number of slow vehicles in 2-lane models, where overtaking is possible, can have a drastic effect. For details we refer to and the discussion in Sec. 9.2.
### 10.3 Randomly placed bottlenecks on the roads and the maximum flux principle
So far we have investigated the effects of two different types of quenched randomness both of which were associated with the vehicles (i.e., particles). We now consider the effects of yet another type of quenched randomness which is associated with the road (i.e., lattice).
In order to anticipate the effects of such randomness associated with the highway, let us begin with the simplest possible caricature of traffic with a ”point defect” : a single ”impurity” (or, ”defect”) site in the deterministic limit $`p=0`$ of the NaSch model with $`v_{max}=1`$. In this model, vehicles move forward, in parallel, by one lattice spacing if the corresponding site in front is empty; each vehicle takes $`T_{imp}`$ ($`>1`$) timesteps to cross the ”impurity” site but only one time step to cross a normal site when the next site is empty. The impurity sites acts like a blockage for all $`T_{imp}>1`$. As explained in section 8.2.1, in the absence of the impurity, $`J=c`$ for $`0<c1/2`$ and $`J=1c`$ for $`1/2<c1`$. Note that, if the impurity is present, $`1/T_{imp}`$ vehicle passes through the impurity site per unit time. Therefore, the bottleneck created by the impurity introduces an upper cut-off of the flux, viz., $`1/T_{imp}`$. Obviously, $`J=c<1/T_{imp}`$ so long as $`c<c_1=1/T_{imp}`$. Similarly, $`J=1c<1/T_{imp}`$ for $`c>c_2`$ where $`c_2=1c_1`$. In the density interval $`c_1<c<1c_1`$, the bottleneck at the impurity is the flow-limiting factor and, hence, in this regime, $`J=1/T_{imp}`$ is independent of $`c`$. Thus, in the simple caricature of traffic under consideration one would expect the flux to vary with density following the relation
$$J=\{\begin{array}{cc}c\hfill & \mathrm{if}0<cc_1,\hfill \\ 1/T_{imp}\hfill & \mathrm{if}c_1<cc_2,\hfill \\ 1c\hfill & \mathrm{if}c_2<c1,\hfill \end{array}$$
(112)
where
$$c_1=\frac{1}{1+(\mathrm{\Delta }t)_{imp}}\text{and}c_2=\frac{(\mathrm{\Delta }t)_{imp}}{1+(\mathrm{\Delta }t)_{imp}},$$
(113)
and $`T_{imp}=1+(\mathrm{\Delta }t)_{imp}`$ such that $`(\mathrm{\Delta }t)_{imp}=0`$ for the normal sites but $`(\mathrm{\Delta }t)_{imp}>0`$ for the impurity site.
The fundamental diagram, obtained numerically through computer simulation of the NaSch model with a single defect and non-zero $`p`$ (Fig. 48) is in qualitative agreement with those of the fundamental diagram (112). The qualitative features of the fundamental diagram in Fig. 48 are also similar to those of the TASEP with a single defect where the hopping probability $`q`$ is smaller than that at all the normal sites. Equation (112) also indicates that the larger is $`(\mathrm{\Delta }t)_{imp}`$ the lower is the maximum flux $`1/[1+(\mathrm{\Delta }t)_{imp}]`$ and the wider is the interval $`c_1cc_2`$ over which the flux remains constant.
What makes the problem of a single ”point defect” nontrivial is that, over the interval $`c_1cc_2`$ of the density of the vehicles, where the $`J`$ is maximum and independent of $`c`$, the localized blockage has global effects whereby the traffic exhibits macroscopic phase segregation into high-density and low-density regions. Evidence for such macroscopic phase segregation can be obtained directly from the density profiles (see Fig. 49). Fig. 49 implies that so long as $`c<c_1`$ the particles will not pile up but a local increase of density will compensate for the reduced local velocity at the blockage so that the flux around the blockage is identical to that far from it. However, if the global density exceeds $`c_1`$, the particles pile up during the transient period leading to the phase-segregated steady state. Because of the particle-hole symmetry the phase-segregation does not take place if the particle density exceeds $`c_2`$.
We now develop a semi-phenomenological theory<sup>19</sup><sup>19</sup>19A microscopic approach for deterministic dynamics can be found in . for the NaSch model with $`v_{max}=1`$, non-zero $`p`$ and a single ”impurity” site assuming the steady-state to be phase-segregated, as demonstrated by computer simulation (Fig. 49). Naturally, this theory cannot explain the underlying mechanism that gives rise to the phase-segregated structure of the steady-state. But, as we shall see soon, it provides a good estimate of the flux in the phase-segregated regime. Our calculations are based on arguments similar to those suggested originally by Janowsky and Lebowitz in the context of TASEP with a single defect.
Using equation (78), the flux in the high-density and low-density regions, far from their interface, are given by $`J_h=\frac{1}{2}\left(1\sqrt{14qc_h(1c_h)}\right)`$ and $`J_{\mathrm{}}=\frac{1}{2}\left(1\sqrt{14qc_{\mathrm{}}(1c_{\mathrm{}})}\right)`$ and that across the defect bond is given by $`J_{def}\frac{1}{2}\left(1\sqrt{14q_dc_h(1c_{\mathrm{}})}\right)`$. Since, in the steady state, the flux is same across the entire system, we must have $`qc_h(1c_h)=qc_{\mathrm{}}(1c_{\mathrm{}})`$ and, hence,
$$c_h=c_{\mathrm{}}\text{or}c_h=1c_{\mathrm{}}.$$
(114)
The condition $`c_h=c_{\mathrm{}}`$ is satisfied by the uniform density profiles whereas the condition $`c_h=1c_{\mathrm{}}`$ is satisfied by the phase-segregated density profile (see Fig. 49). Moreover, using the condition $`J_h=J_{def}=J_{\mathrm{}}`$ we get
$$c_h(1c_h)=c_{\mathrm{}}(1c_{\mathrm{}})rc_h(1c_{\mathrm{}})$$
(115)
where $`r=q_d/q<1`$ may be interpreted as the ”transmission probability” or ”permeability” of the blockage. From the (approximate) equations (115) we get
$$c_h\frac{1}{r+1}=\frac{p}{p+p_d}\text{and}c_{\mathrm{}}\frac{r}{r+1}=\frac{p_d}{p+p_d}$$
(116)
and, hence,
$$J=\frac{1}{2}\left[1\sqrt{1\frac{4qr}{(1+r)^2}}\right]$$
(117)
The estimate (117) is in good agreement with the numerical data (Fig. 48) obtained from computer simulation . Moreover, the estimates $`c_{\mathrm{}}`$ and $`c_h`$ are also in good agreement with $`c_{\mathrm{}}`$ and $`c_h`$, respectively, in Fig. 49.
Note that $`c_{\mathrm{}}`$ and $`c_h`$ depend only on $`r`$ and are independent of $`c`$. Moreover, the estimates (116) of $`c_{\mathrm{}}`$ and $`c_h`$ are in excellent agreement with $`c_1`$ and $`c_2`$, respectively, in Fig. 48. At first sight, these two results may appear surprising and counter-intuitive. But, we’ll now show that these are related to the mechanism of the phase-segregation. Conservation of the vehicles demand that
$$cL=c_hh+c_{\mathrm{}}=c_hh+c_{\mathrm{}}(Lh)$$
(118)
where $`h`$ and $`\mathrm{}=Lh`$ are the lengths of the high-density and low-density regions, respectively. Thus,
$$\frac{h}{L}=\frac{cc_{\mathrm{}}}{c_hc_{\mathrm{}}}=\frac{c(1+r)r}{1r}$$
(119)
The equation (119) shows that $`h/L0`$ as $`cc_{\mathrm{}}`$ and $`h/L1`$ as $`cc_h`$. Therefore, keeping $`r`$ fixed as the density is increased beyond $`c_1=c_{\mathrm{}}`$, the densities of the two regions remain fixed but the high-density region grows thicker at the cost of the length of the low-density region as more and more vehicles pile up and, eventually, at $`c=c_2=c_h`$ the low-density region occupies a vanishingly small fraction of the total length of the system signaling the disappearance of the phase segregation. Interestingly, recasting the expressions for $`c_h`$ and $`c_{\mathrm{}}`$ as $`c_{\mathrm{}}=1/[1+(\mathrm{\Delta }t)_{imp}]`$ and $`c_h=(\mathrm{\Delta }t)_{imp}/[1+(\mathrm{\Delta }t)_{imp}]`$, where $`(\mathrm{\Delta }t)_{imp}=1/r`$, we find close formal analogies with $`c_1`$ and $`c_2`$, respectively, in equation (113) .
Schütz considered a TASEP with sublattice-parallel update (see App. A) where the motion of the particles is deterministic (i.e., $`q=1`$) everywhere except at a defect site where they move with the probability $`q_d<1`$ (i.e., $`r=q_d<1`$). Exact solution is possible through a mapping on a 6-vertex model. Later, a solution using the matrix-product Ansatz was presented in . Except for minor differences, the qualitative features of the results do not differ from the corresponding approximate results obtained for $`q_d1`$ . Qualitatively similar phase segregation phenomena have also been observed in a related model .
The qualitative features of the fundamental diagram do not change significantly if the ”point-like” defect (or, impurity) is replaced by an ”extended” defect , i.e. a few consecutive defect sites. However, with increasing length of the defect, the maximum value of the flux decreases monotonically and approaches the maximum flow of the homogeneous system where the hopping probability associated with each of the bonds is identical to that associated with the defects in our model (Fig. 50). From Figs. 50 and 51 we conclude that the monotonic decrease of the flow with increasing length of the extended defects, leads to a larger difference $`c_hc_{\mathrm{}}`$ between the densities of the high-density and the low-density regions of the phase-segregated steady-state.
Next, instead of a single point-like or extended ”defect”, let us consider the more general case of quenched ”disorder” in the NaSch model with $`v_{max}=1`$ where the quenched random hopping probabilities $`q_{j,j+1}=1p_{j,j+1}`$ are chosen independently from some probability distribution $`P(q)`$, for the hopping from the cell $`j`$ to the cell $`j+1`$ ($`j=1,2,\mathrm{},L`$). For a given realization of the disordered system, every vehicle hopping from a given cell $`i`$ to the next cell $`i+1`$ must hop with the same probability $`q_{i,i+1}`$ and a given vehicle hops across different bonds, in general, with different probabilities assigned to these bonds as it moves forward with time. A similar generalization of the TASEP has also been studied . We shall refer to this model as disordered TASEP (i.e., DTASEP).
Suppose, $`q`$ are chosen from the binary distribution
$$P(q_{j,j+1}=q_d)=f,P(q_{j,j+1}=q)=1f$$
(120)
i.e., a fraction $`f`$ of the bonds have a permeability $`r<1`$ while the remaining fraction $`1f`$ have unit permeability. A mean-field theory has been developed (see Appendix G for details) for computing the fundamental diagram of the DTASEP. The flux in this model has interesting symmetry properties under the operations of ”charge conjugation” (which interchanges particles and holes), ”parity” (which interchanges forward and backward hopping rates on each bond and ”time reversal” (which reverses the direction of the current) .
The quenched disorder in these ”disordered” models can be viewed as ”point-like impurities” distributed randomly over the lattice. But, the qualitative features of the fundamental diagram of DTASEP are similar to those observed for a single point-like defect and those for a single extended defect. Although the random distribution of the point-like impurities leads to a ”rough” density profile for all densities, in an intermediate regime of density, phase-segregated steady-sates with macroscopic high- and low-density regions have been identified.
What is the underlying mechanism for the ”macroscopic” phase segregation in all the models DTASEP ? Let us denote the stretches of bonds with permeability $`1`$ by $`X`$ and the stretches of bonds with permeability $`r`$ by $`Y`$. The two parabolas in Fig. 52 are the two steady-state fundamental diagrams for the two pure reference systems consisting of all $`X`$ and all $`Y`$, respectively. Since the flux must be spatially constant in the steady-state, the possible densities are given by the four intersections of the line $`J=J_0`$ with the two parabolas. If the average density is less (greater) than $`1/2`$ then the two possible densities are $`c_1`$ and $`c_2`$ ($`c_3`$ and $`c_4`$). The variation of density between $`c_1`$ and $`c_2`$ (or $`c_3`$ and $`c_4`$) in the $`X`$ and $`Y`$ stretches is merely micro phase-segregation while, on a macroscopic scale, the density remains uniform. For simplicity, we assume that the density in each stretch of like bonds is uniform. The global density of the system is approximately $`c(1f)c_{1,4}(J_0)+fc_{2,3}(J_0)`$ where $`f`$ is given by equation (120). However, as the density increases the flux also increases till it attains the maximum allowed flux of the pure system consisting of all $`Y`$ (this happens at a global density smaller than $`1/2`$). What happens when the density increases further? According to the ”maximum current principle” , no further increase of the flux is possible and the excess density is taken care of by increasing the density in some of the $`X`$ stretches from $`c_1`$ to $`c_4`$ (or, vice versa if $`c>1/2`$). This conversion takes place adjacent to the largest stretch of $`Y`$ bonds where the density also changes from $`c_2`$ to $`c_3`$ (or, vice versa if $`c>1/2`$) to accommodate the additional particles added to the system. This leads to the macroscopic phase segregation as the system consists of two macroscopic regions of two different mean densities- one with lower densities $`c_1,c_2`$ in the $`X`$ and $`Y`$ stretches and the other with the higher densities $`c_3,c_4`$ in the $`X`$ and $`Y`$ stretches.
It is not difficult to generalize the DTASEP to disordered NaSch (DNaSch) model by replacing the random sequential updating by parallel updating. However, we face subtle conceptual difficulties in extending DNaSch model further to arbitrary $`v_{max}`$; if the position of a vehicle at time $`t+1`$ is decided by $`q`$ at its current position at time $`t`$, it may be forced to move $`v`$ sites downstream by hopping over sites of even smaller $`q`$ if $`v>1`$.
Some alternative parameterization of the defect (or, disordered) sites in the NaSch model for arbitrary $`v_{max}`$ have also been suggested. In localized defects have been investigated where the randomization parameter $`p_d`$ is larger than in the rest of the system. Csahok and Vicsek have considered the blockages as sites with a ”permeability” smaller than unity whereas the permeability of all the other sites is unity. This effectively reduces $`v_{max}`$ while the vehicle is at a blockage. On the other hand, Emmerich and Rank considered a model of where the velocity of every vehicle in the region occupied by the blockage (or, more appropriately, hindrance) at the time step $`t+1`$ is half of that at time $`t`$, i.e., $`v_n(t+1)=v_n(t)/2`$ is the $`n`$-th vehicle is located within the hindrance region. Some effects of static hindrances on vehicular traffic have also been investigated following alternative approaches, e.g., car-following theory . From a practical point of view, ramps have effects very similar to those of a static defect. For the NaSch model this has been investigated in .
## 11 Other CA models of highway traffic
All the CA models of highway traffic described so far are basically generalizations of the minimal CA model proposed originally by Nagel and Schreckenberg . We now describe a few other alternative minimal CA models and the interesting features of the corresponding results.
### 11.1 Fukui-Ishibashi model
The update rules of the Fukui and Ishibashi (FI) model of single-lane highway traffic are as follows:
If $`v_{max}`$ or more sites in front of the $`n`$-th vehicle is empty at the time step $`t`$, then it has a probability $`1p`$ to move forward by $`v_{max}`$ sites and a probability $`p`$ to move forward by $`v_{max}1`$ sites in the time step $`t+1`$. However, if only $`d`$ sites ($`d<v_{max}`$) in front of the $`n`$-th vehicle are empty at time $`t`$ then it moves by $`d`$ sites in the next time step. Since no site can be occupied simultaneously by more than one vehicle, a vehicle must not move forward in the time step $`t+1`$ if the site immediately in front of it is occupied by a vehicle at the time step $`t`$. The model becomes deterministic in both the limits $`p=0`$ and $`p=1`$.
The FI model differs from the NaSch model in two respects: (a) the increase of speed of the vehicles is not necessarily gradual and (b) the stochastic delay applies only to high-speed vehicles. The FI model, obviously, reduces to the NaSch model if $`v_{max}=1`$. A site-oriented mean-field theory and a car-oriented mean-field theory for the FI model have been developed for arbitrary $`v_{max}`$ and $`p`$. Note that the FI model is equivalent to the deterministic CA rule $`184`$ (in the notation of Wolfram ) in the limit $`v_{max}=1,p=0`$. Generalizations of the deterministic limit $`p=0`$ of the FI model have also been proposed .
According to the classification of Sec. 7.3, the FI model belongs to class I, i.e. the high-acceleration limit where no spontaneous jamming exists. In an alternative high-acceleration variant has been proposed. Here only the acceleration step $`(U1)`$ of the NaSch model is changed to $`vv_{max}`$, i.e. all vehicles accelerate immediately to the maximal possible velocity. The other update steps of the NaSch model are left unchanged. In contrast to the FI model, all vehicles are subject to the randomization step. The behaviour of this variant is therefore similar to that of the NaSch model, e.g. one finds spontaneous jam formation.
### 11.2 Galilei-invariant vehicle-vehicle interaction and metastability
In the NaSch model it is postulated that the gap between a pair of successive vehicles is adjusted according to the velocity of the leading vehicle alone. In contrast, often in real traffic, drivers tend to adjust the gap in front taking into account the difference between the velocity of their own vehicle and that of the leading vehicle. The latter aspect of real traffic is captured by a recent model developed by Werth, Froese and Wolf (WFW) .
If both the following vehicle and the leading vehicle move with constant acceleration $`b`$, then a collision between the two can be avoided provided
$$gap+\mathrm{}(v_{LV})v_{FV}\tau _r+\mathrm{}(v_{FV})$$
(121)
where $`v_{LV}`$ and $`v_{FV}`$ are the velocities of the leading vehicle and the following vehicle, respectively, $`\tau _r`$ is the reaction time of the following vehicle and
$$\mathrm{}(v)=\frac{v^2}{2b}$$
(122)
is the distance covered by a vehicle with initial velocity $`v`$ before it comes to a stop by moving with a constant deceleration $`b`$. Using equation (122) the condition (121) for avoiding collision can be written as
$$gapv_{FV}\tau _r+\frac{\overline{v}}{b}(v_{FV}v_{LV})$$
(123)
where $`\overline{v}=(v_{FV}+v_{LV})/2`$ if the average velocity of the pair of vehicles under consideration. Therefore, a sufficient condition for avoiding a collision is
$$gap\{\begin{array}{cc}v_{FV}\tau _r\hfill & \text{for }v_{FV}v_{LV}\text{,}\hfill \\ v_{FV}\tau _r+\frac{v_{max}}{b}(v_{FV}v_{LV})\hfill & \text{for }v_{FV}>v_{LV}\text{.}\hfill \end{array}$$
(124)
In the limiting case $`b\mathrm{}`$ the sufficient condition (124) reduces to $`gapv_{FV}\tau _r`$ which is identical to the form of vehicle-vehicle interaction in the NaSch model if one chooses $`\tau _r`$ as the unit of time. In the opposite limit $`\tau =0`$, the sufficient condition (124) reduces to
$$gap\{\begin{array}{cc}0\hfill & \text{for }v_{FV}v_{LV}\text{,}\hfill \\ \frac{v_{max}}{b}(v_{FV}v_{LV})\hfill & \text{for }v_{FV}>v_{LV}\text{.}\hfill \end{array}$$
(125)
Since this type of vehicle-vehicle interaction involves the difference of the velocities $`v_{FV}v_{LV}`$, it is clearly invariant under a Galilean transformation and, hence, the name. The vehicle-vehicle interactions in real traffic may be somewhere in between the two limiting cases of NaSch model and the Galilei-invariant model.
Suppose the indices $`n1`$ and $`n`$ label the leading vehicle and the following vehicle, respectively, of a pair. The update rules suggested by WFW for implementing the Galilei-invariant vehicle-vehicle interaction are as follows:
Step 1: Acceleration.
$$v_n^{(1)}=\mathrm{min}(v_n+1,v_{max})$$
Step 2: Deceleration (due to other vehicles).
$$v_n^{(2)}=\mathrm{min}(v_n^{(1)},d_n1+v_{n1})$$
Step 3: Randomization.
$$v_n^{(3)}\stackrel{p}{=}\mathrm{max}(v_n^{(2)}1,0)\text{with probability }p$$
Step 4: Deceleration (due to other vehicles).
$$v_n^{(4)}=\mathrm{min}(v_n^{(3)},d_n1+v_{n1}^{(4)})$$
Step 5: Vehicle movement.
$`x_n`$ $`=`$ $`x_n+v_n^{(4)},`$
$`v_n`$ $`=`$ $`v_n^{(4)}.`$
Thus, the rule for deceleration (due to other vehicles) is applied twice. Step 4 makes sure collisions are avoided. Since also the new velocity $`v_{n1}^{(4)}`$ of the preceding car enters, this step can not be performed in parallel for all cars. Instead it is performed sequentially, but the final configuration is independent of the starting point of this sequential updating. Step 4 has then to be applied twice in order to determine all velocities $`v_n^{(4)}`$ consistently.
The rules as given above define the retarded version of the Galilei-invariant model. In the non-retarded version, in step 2 $`v_{n1}`$ is replaced by the new velocity $`v_{n1}^{(2)}`$. To determine $`v_n^{(2)}`$ consistently for all cars, step 2 has then to be iterated $`v_{max}1`$ times.
The most interesting feature of the Galilei-invariant model is that its fundamental diagram has a metastable branch although its update scheme involve neither cruise-control nor slow-to-start rules. The mechanism leading to the existence of metastable states is different from the models with slow-to-start rules (see Sec. 9.1.2). The outflow from jams is the same as in the NaSch model since it is independent of the interaction between vehicles. However, due to the inclusion of anticipation effects (i.e. the driver knows the velocity of the preceding vehicle) the free-flow is less sensitive to fluctuations.
### 11.3 CA versions of the optimal-velocity model
The traffic jams appear spontaneously in both the OV models and the CA models. However, in the OV models spontaneous formation of the jams are caused by the non-linearity of the dynamical equations whereas in the CA models it is triggered primarily by the stochasticity of the update ”rules”. The mechanism for the spontaneous formation of jams in real traffic may be a combination of these two.
In the OV model, the control of velocity is given by the control of acceleration through the OV function which gives the optimal velocity for the current distance-headway. Thus, unlike the CA models like the NaSch model, the vehicles in the OV models get an opportunity to avoid crash without any need to exert unphysically large deceleration. In fact, collision of vehicles may take place in naive discretization of the dynamical equations for the OV models unless special care is taken in the discretization process (see below).
In the following we present several CA model analogues of the OV model that have been proposed by different authors. In principle, the NaSch model is also an OV model, but with a linear OV function, $`v(d)=\mathrm{min}[d1,v_{max}]`$. The first attempt to generalize this relation is due to Emmerich and Rank . The update rules of their model are as follows:
Step 1: Find largest gap.
Find the vehicle with the largest gap to the next vehicle ahead.
Step 2: Acceleration.
$$v_n\mathrm{min}(v_{max},v_n+1)$$
Step 3: Deceleration due to other vehicles.
$$v_nM_{d_n1,v_n}\text{if }d_n1v_{max},$$
i.e. a vehicle with velocity $`j`$ and $`i`$ empty cells in of it (i.e. a gap $`d_n=i+1`$) reduces its velocity to $`M_{i,j}`$ ($`0i,jv_{max}`$).
Step 4: Randomization.
$$v_n\mathrm{max}(v_n1,0)\text{with probability }p$$
Step 5: Vehicle movement.
$$x_nx_n+v_n$$
Step 6: Next vehicle.
Repeat steps 2-5 for the next vehicle behind, i.e. proceed in the direction opposite to the motion of the vehicles.
For the NaSch model with $`v_{max}=5`$ the matrix $`M_{i,j}`$ is given by
$$M_{i,j}^{(NaSch)}=\left(\begin{array}{cccccc}0& 0& 0& 0& 0& 0\\ 0& 1& 1& 1& 1& 1\\ 0& 1& 2& 2& 2& 2\\ 0& 1& 2& 3& 3& 3\\ 0& 1& 2& 3& 4& 4\\ 0& 1& 2& 3& 4& 5\end{array}\right)$$
(126)
A general matrix $`M_{i,j}`$ has to satisfy certain conditions (e.g. $`M_{i,j}\mathrm{min}(i,j)`$ and $`M_{i,j}M_{i,k}`$ for $`jk`$) to guarantee e.g. the absence of collisions in the model. In order to model the fact that faster vehicles keep a relatively larger headway to the preceding vehicle, Emmerich and Rank suggested the following matrix:
$$M_{i,j}^{(ER)}=\left(\begin{array}{cccccc}0& 0& 0& 0& 0& 0\\ 0& 1& 1& 1& 1& 1\\ 0& 1& 2& 2& 2& 2\\ 0& 1& 2& 2& 3& 3\\ 0& 1& 2& 3& 3& 4\\ 0& 1& 2& 3& 4& 4\end{array}\right)$$
(127)
Using a parallel update scheme, the model shows unrealistic behaviour in the free-flow regime, especially in the deterministic limit $`p1`$. Here the fundamental diagram is non-monotonic , as can be seen from a simple example for $`p=0`$. At density $`c=1/7`$, the stationary state is of the form $`5\mathrm{}\mathrm{}5\mathrm{}\mathrm{}5\mathrm{}\mathrm{}`$ where numbers denote the velocity of vehicles and ’.’ an empty cell. At density $`c=1/6`$, on the other hand, the stationary state is $`4\mathrm{}\mathrm{..4}\mathrm{}\mathrm{..4}\mathrm{}..`$. Comparing the corresponding flows, one finds $`J(c=1/7)=5/7>J(c=1/6)=2/3`$. At density $`c=1/5`$, the stationary state is $`4\mathrm{}.4\mathrm{}.4\mathrm{}.`$ with flow $`J(c=1/5)=4/5`$ which is again larger than the flow at $`c=1/6`$ and corresponds to the maximal possible flow. This kind of behaviour persists even in the presence of randomness ($`p>0`$) . In order to circumvent this problem, Emmerich and Rank had to introduce a special kind of ordered-sequential update, where first the vehicle with the largest gap ahead is updated. Then, the position of the next vehicle upstream is updated, and so on, using periodic boundary conditions.
Emmerich and Rank also investigated more general rules where even for gaps larger than $`v_{max}`$ the velocity of the vehicles is reduced to a value $`v<v_{max}`$.
Later, a similar model has been proposed by Helbing and Schreckenberg . It is closer to the spirit of the original optimal-velocity model (see Sec. 6.2).
Step 1: Vehicle movement.
$$x_nx_n+v_n(t)$$
Step 2: Acceleration.
$$v_n^{}(t+1)=v_n(t)+\lambda [V_{opt}(d_n(t))v_n(t)],$$
Step 3: Randomization.
$$v_n(t+1)=v_n^{}(t+1)\{\begin{array}{cc}1\hfill & \text{with probability }p\text{ (if }v_n^{}(t+1)>0\text{)}\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
Here $`y`$ denotes the floor function, i.e. the largest integer $`iy`$. In various optimal velocity functions $`V_{opt}(d)`$ have been used in order to fit experimental data. The simplest, but unrealistic, choice was $`V_{opt}(d)=\mathrm{min}(d,v_{max})`$ where $`d`$ is the distance-headway. The parameter $`\lambda `$ corresponds formally to the sensitivity parameter in the OV model where it determines the timescale of relaxation towards the stationary fundamental diagram. Such an interpretation is not possible for discrete time models. Here the main effect of the parameter $`\lambda `$ is a rescaling of the OV function.
The naive discretization of the OV function produces some undesirable features of the model, e.g. the flow corresponding to the OV function is non-monotonic in the free-flow region. Furthermore one finds – for certain initial conditions – a breakdown of the flow at a finite density $`c_{}<1`$. This indicates a lack of robustness of model against small modifications of the rules.
For $`\lambda <1`$ the model is not intrinsically collision-free , in contrast to most other models discussed in this review. Problems occur e.g. when fast vehicles approach the end of a jam. For the $`V_{opt}(d)=\mathrm{min}(d,v_{max})`$ and the ’realistic’ choice $`\lambda =0.77`$ collisions can occur for $`v_{max}5`$ . For $`\lambda >1`$, on the other hand, a backward motion of vehicles is possible. For a given OV function it is possible to derive conditions for the parameter $`\lambda `$ which ensure the absence of collisions and backward motion .
Nagatani has suggested a CA model which combines the OV idea with the TASEP. It is the discrete analogue of the simplified OV model (51) presented in Sec. 6.2. Here the $`n`$-th vehicle moves ahead with probability $`v_n(t)`$ where $`v_n(t)`$ is interpreted as velocity. This velocity is obtained by integration of the OV function
$$\ddot{x}_n(t)=\{\begin{array}{cc}a\hfill & \text{ for }\mathrm{\Delta }x_n(t)\mathrm{\Delta }x_c\hfill \\ a\hfill & \text{ for }\mathrm{\Delta }x_n(t)<\mathrm{\Delta }x_c\hfill \end{array}$$
(128)
where $`a>0`$, $`\mathrm{\Delta }x_n(t)=x_{n+1}(t)x_n(t)`$ and $`\mathrm{\Delta }x_c`$ is a safety distance. Furthermore the velocity is restricted to the interval $`0v_n(t)v_{max}1`$.
### 11.4 CA from ultra-discretization
Several CA which can be interpreted as traffic models have been derived using the so-called ultra-discretization method (UDM) . This approach allows to establish a direct connection between certain differential equations and CA. The problem in the derivation of CA from differential equations lies in the discretization of the ’state’ (or dependent) variable. E.g. in a numerical treatment, only space and time variables are discretized.
The basic procedure of the UDM is as follows: 1) Start from a nonlinear wave equation, e.g. the KdV equation or Burgers’ equation; 2) Discretize space- and time-variables in a standard way to obtain the discrete analogue of the wave equation which is still continuous in the state variable $`u_j(t)`$; 3) The discrete analogue is now ultra-discretized. Defining $`U_j(t)=ϵ\mathrm{ln}(u_j(t))`$ (where $`ϵ`$ depends on the discretization $`\mathrm{\Delta }x`$ and $`\mathrm{\Delta }t`$ of space and time) one can use the identity $`lim_{ϵ0}ϵ\mathrm{log}(e^{A/ϵ}+e^{B/ϵ}+\mathrm{})=\mathrm{max}[A,B,\mathrm{}]`$ to derive the CA analogue fo the nonlinear wave equation.
By applying the UDM to Burgers’ equation $`v_t=2vv_x+v_{xx}`$ one obtains the $`(L+1)`$-state, deterministic CA .
$$n_j(t+1)=n_j(t)+\mathrm{min}[M,n_{j1}(t),Ln_j(t)]\mathrm{min}[M,n_j(t),Ln_{j+1}(t)].$$
(129)
$`n_j(t)`$ is the occupation number of cell $`j`$ at time $`t`$. In contrast to most other CA discussed in this review multiple occupations of cells are allowed. The maximum number of particles which can occupy the same cell is given by $`L`$, i.e. $`0n_j(t)L`$. The model defined by (129) might therefore by interpreted as a simple model for a highway with $`L`$ lanes where the effects of lane changes are completely neglected. In it has also been suggested to interprete $`n_j(t)/L`$ for large $`L`$ as a coarse-grained density. The parameter $`M`$ denotes the maximum number of vehicles that can move from cell $`j`$ to cell $`j+1`$.
For $`ML=1`$ the model reduces to the rule-184 CA, i.e. the NaSch model with $`v_{max}=1`$ and $`p=0`$. In the case $`L2M`$ the fundamental diagram looks similar to that of rule-184 CA, but it is ’degenerate’ in the sense that qualitative different stationary states with the same flow exist. For $`L>2M`$ the fundamental diagram the fundamental diagram resembles that of rule-184 CA with a blockage site (see Sec. 10.3, especially Fig. 48). In the region $`M/Lc(LM)/L`$ the flow takes the constant value $`M/L`$, i.e. the parameter $`M`$ can be interpreted as a flow limiter.
In also a generalization of (129) to higher velocities $`v_{max}>1`$ has been suggested. For $`v_{max}=2`$ the generalized update rule is given by
$`n_j(t+1)`$ $`=`$ $`n_j(t)+a_{j2}(t)a_j(t)`$ (130)
$`+`$ $`\mathrm{min}[b_{j1}(t)a_{j1}(t),Ln_j(t)a_{j2}(t)]`$
$``$ $`\mathrm{min}[b_j(t)a_j(t),Ln_{j+1}(t)a_{j1}(t)].`$
Here $`b_j(t)=\mathrm{min}[n_j(t),Ln_{j+1}(t)]`$ is the maximum number of vehicles at site $`j`$ that can move and $`a_j(t)=\mathrm{min}[n_j(t),Ln_{j+1}(t),Ln_{j+2}(t)]`$ is the number of vehicles that move two cells forward. The idea behind this dynamics is that first the vehicles try to move two cells forward. This is only possible if the two cells ahead are not fully occupied, i.e. $`Ln_{j+1}(t)>0`$ and $`Ln_{j+2}(t)>0`$. Then $`\mathrm{min}[b_j(t)a_j(t),Ln_{j+1}(t)a_{j1}(t)]`$ vehicles move forward one cell.
The model defined by (130) can be considered as a generalization of the Fukui-Ishibashi model (see Sec. 11.1) to which it reduces for $`L=1`$. Note that the flow limiter $`M`$ has been dropped in the extended model.
The fundamental diagram of the model (130) has a structure similar to that shown in Fig. 34, i.e. states of high flow exist. Due to the higher velocity the degeneracy found in the model (129) is lifted. In the case $`L=2`$, the $`c=1/2`$ states $`\mathrm{}1111\mathrm{}`$ and $`\mathrm{}2020\mathrm{}`$ are degenerate in the simple model. In the $`v_{max}=2`$ case, however, in the state $`\mathrm{}1111\mathrm{}`$ the vehicles can increase their velocity, while in the state $`\mathrm{}2020\mathrm{}`$ they can not due to the presence of the fully-occupied cells ’2’. The high-flow states are unstable against local perturbations.
Another CA model obtained by ultra-discretization of the modified KdV-equation has been suggested in . It is a second-order CA since the configuration at time $`t+1`$ does not only depend on the configuration at time $`t`$, but also on the previous one at time $`t1`$. The update rule for the position $`x_j(t)`$ of vehicle $`j`$ is given by
$$x_j(t+2)=x_j(t)+\mathrm{\Delta }x_j(t)+\mathrm{max}[0,\mathrm{\Delta }x_j(t+1)M]\mathrm{max}[0,\mathrm{\Delta }x_j(t+1)L]$$
(131)
where $`\mathrm{\Delta }x_j(t)=x_{j+1}(t)x_j(t)`$ is the gap and $`M`$ and $`L`$ are constants.
Although this model is deterministic it exhibits start-stop waves similar to those found in the NaSch model . It seems that a second-order deterministic CA might produce effects similar to those of noise in a first-order stochastic CA. One should note, however, that the rules allow vehicles to move backwards. Other models obtained using the UDM have been discussed in .
## 12 Cellular automata models of city traffic and road networks
### 12.1 Biham-Middleton-Levine model and its generalizations
In the BML model , each of the sites of a square lattice represent the crossing of a east-west street and a north-south street. All the streets parallel to the $`\widehat{x}`$-direction of a Cartesian coordinate system are assumed to allow only single-lane east-bound traffic while all those parallel to the $`\widehat{y}`$-direction allow only single-lane north-bound traffic. Let us represent the east-bound (north-bound) vehicles by an arrow pointing towards east (north). In the initial state of the system, vehicles are randomly distributed among the streets. The states of east-bound vehicles are updated in parallel at every odd discrete time step whereas those of the north-bound vehicles are updated in parallel at every even discrete time step following a rule which is a simple extension of the TASEP: a vehicle moves forward by one lattice spacing if and only if the site in front is empty, otherwise the vehicle does not move at that time step.
Thus, the BML model is also a driven lattice gas model where each of the sites can be in one of the three possible states: either empty or occupied by an arrow $``$ or $``$. Note that the parallel update rules of the BML model is fully deterministic and, therefore, it may also be regarded as a deterministic CA. The randomness arises in this model only from the random initial conditions . Suppose, $`N_{}`$ and $`N_{}`$ are the numbers of the east-bound and north-bound vehicles, respectively, in the initial state of the system. If periodic boundary conditions are imposed in all directions, the number of vehicles in every street is conserved since no turning of the vehicles are allowed by the updating rules. In a finite $`L\times L`$ system the densities of the east-bound and north-bound vehicles are given by $`c_{}=N_{}/L^2`$ and $`c_{}=N_{}/L^2`$, respectively, while the global density of the vehicles is $`c=c_{}+c_{}`$.
Computer simulations of the BML model with periodic boundary conditions demonstrate that a first order phase transition takes place at a finite non-vanishing density $`c_{}`$, where the average velocity of the vehicles vanishes discontinuously signaling complete jamming; this jamming arises from the mutual blocking (”grid-locking”) of the flows of east-bound and north-bound traffic at various different crossings (see for the corresponding results of the BML model with open boundary conditions).
At concentrations just above $`c_{}`$, in the jammed phase, all the vehicles together form a single cluster which is stretched along the diagonal connecting the south-west to the north-east of the system. In other words, the lowest-density jammed configurations consist of a single diagonal band where the $``$ and $``$ occupy nearest-neighbour sites on the band in a zigzag manner. With further increase of density more and more vehicles get attached to the band in the form of off-diagonal branches and the infinite cluster of the jammed vehicles looks more and more random. Thus, in general, a typical infinite cluster of the jammed vehicles consists of a ”backbone” and ”dangling vehicles” which are the analogs of the ”backbone” and the ”dangling ends” of the infinite percolation clusters in the usual site/bond percolation . However, in contrast to the infinite percolation cluster in the usual random site/bond percolation, the infinite spanning cluster of vehicles in the BML model emerges from the self-organization of the system. Nevertheless, concepts borrowed from percolation theory have been used to characterize the structure of the infinite cluster of jammed vehicles in the BML model at $`c>c_{}`$ . The distribution of the waiting times of the vehicles at the signals (i.e., at the lattice sites) has also been investigated through computer simulations .
#### 12.1.1 Poor man’s mean-field estimates for the BML model
If one is not interested in detailed information on the ”structure” of the dynamical phases, one can get a mean-field estimate of $`c_{}`$ by carrying out a back-of-the-envelope calculation. Suppose, $`c_{}`$ and $`c_{}`$ denote the average densities while $`v_{}`$ and $`v_{}`$ denote the average speeds of the east-bound and north-bound vehicles, respectively. In order to estimate $`c_{}`$ one has to take into account interaction of the east-bound (north-bound) vehicles not only with the north-bound (east-bound) vehicles but also, in a self-consistent manner, with other east-bound (north-bound) vehicles . Following the arguments of Appendix I, one can show that in the symmetric case $`c_{}=c_{}=c`$, a self-consistency equation for the speed is
$$v_{}=v_{}=v=\frac{1}{2}\left[1+\frac{c}{2}+\sqrt{\left(1+\frac{c}{2}\right)^24c}\right]$$
(132)
for $`c<c_{}`$. The critical density $`c_{}`$ is determined by the condition that at $`cc_{}`$ the equation (132) does not give a real solution. Hence we get $`c=c_{}=6\sqrt{32}0.343`$ which, in spite of the approximations made, is surprisingly close to the corresponding numerical estimate obtained from computer simulation . Moreover, the mean-field estimate $`c_{}0.343`$ is also consistent with the more rigorous result<sup>20</sup><sup>20</sup>20Ishibashi and Fukui claimed that complete jamming can occur in the BML model only for $`c=1`$. However, a plausible flaw in their arguments was pointed out by Chau et al. . In it has been argued that $`c_{}L^{0.14}`$, i.e. $`c_{}=0`$ in the thermodynamic limit. that $`c_{}`$ is strictly less than $`1/2`$. The BML model in three dimension, although not relevant for vehicular traffic, has also been studied .
#### 12.1.2 Mean-field theory of the BML model
Recall that the occupation variables $`n(i;t)`$ in the NaSch model describe the state of occupation of the sites $`i`$ ($`i=1,2,\mathrm{}`$) by the vehicles on the one-dimensional highway. A scheme for a truly microscopic analysis of the BML model begins by introducing the corresponding generalized occupation variables $`n_{}(x,y;t)`$ and $`n_{}(x,y;t)`$, which describe the state of occupation of the sites $`(x,y)`$ by the north-bound and east-bound vehicles, respectively, on the two-dimensional street-network. The analogs of the equations (64) and (65) (see Appendix I for details) have analogous physical interpretations. As usual, in the naive SOMF approximation one neglects the correlations between the occupations of different sites . Since none of the sites is allowed to be occupied by more than one particle at a time, Pauli operators may be used to develop an analytical theory of vehicular traffic but one should keep in mind that the particles representing the vehicles are purely classical and the system does not have any quantum mechanical characteristics.
#### 12.1.3 Generalizations and extensions of the BML model
The BML model has been generalized and extended to take into account several realistic features of traffic in cities.
$``$ Asymmetric distribution of the vehicles:
Suppose the vehicles are distributed asymmetrically among the east-bound and north-bound streets , i.e., $`c_{}c_{}`$. For convenience, let us write $`c_{}=cf_a`$ and $`c_{}=c(1f_a)`$ where $`f_a`$ is the fraction of the vehicles moving towards north. Clearly $`f_a=1/2`$ corresponds to the symmetric case $`c_{}=c_{}=c/2`$. On the other hand, $`f_a=0`$ ($`f_a=1`$) correspond to the extreme asymmetric case where all the vehicles are east-bound (north-bound). Obviously, the absence of grid-locking in the extreme limits $`f_a=0`$ and $`f_a=1`$ rules out the possibility of BML-like complete jamming transition, i.e., $`c_{}=1`$ for both $`f_a=0`$ and $`f_a=1`$. Moreover, $`c_{}`$ decreases with decreasing asymmetry in the distribution of the vehicles; $`c_{}`$ is the smallest for $`f_a=1/2`$, i.e., symmetric distribution of the vehicles. These results can be presented graphically by plotting the curve $`c_{}(f_a)`$ in the phase diagram on the $`cf_a`$ plane .
$``$ Unequal maximum velocities:
In the BML model both east-bound and north-bound vehicles can move by a maximum of one lattice spacing at a time and, therefore, the average speeds of both types of vehicles can never exceed unity. On the other hand, recall that in the Fukui-Ishibashi model of highway traffic vehicles can move up to a maximum of $`M`$ lattice sites at a time and, hence, have average velocities larger than unity. Incorporating similar high-speed vehicles Fukui et al. generalized the BML model where the east-bound vehicles are allowed tom ove by $`M`$ sites at a single time step while the north-bound vehicles can move by only one lattice spacing.
$``$ Overpasses or two-level crossings:
The BML model has been extended to take into account the effects of overpasses (or two-level crossings) . A fraction $`f_o`$ of the lattice sites in the BML model are randomly identified as overpasses each of which can accommodate up to a maximum of two vehicles simultaneously. The overpasses weaken the grid-locking in the BML model. Therefore, $`c_{}`$ is expected to increase with increasing $`f_o`$. Besides, $`c_{}`$ is expected to be unity if $`f_o=1`$. Naturally, we address the question: does jamming disappear (i.e., $`c_{}`$ is unity) only at $`f_o=1`$ or for even smaller values of $`f_o`$? In order to answer this question we extend the self-consistent mean-field arguments, which led to the equation (132), incorporating the effects of the overpasses thereby getting the generalized self-consistency equation
$$v_{}=v_{}=v=\frac{1}{2}\left[1+\frac{1f_o}{2}c+\sqrt{\left(1+\frac{1f_o}{2}c\right)^24(1f_o)c}\right]$$
(133)
in the symmetric case $`c_{}=c_{}=c`$. The equation (133) reduces to the equation (132) in the limit $`f_o=0`$. It predicts that, if $`f_o0`$, a moving phase exists in the BML model with overpasses for vehicle densities $`cc_{}=(6\sqrt{32})/(1f_o)`$ and that the jammed phase disappears altogether (i.e., $`c_{}`$ becomes unity) for $`f_o10.343=0.657`$, an underestimate when compared with the corresponding computer simulations. These results can be presented graphically by drawing the curve $`c_{}(f_o)`$ in the mean-field phase diagram on the $`cf_o`$ plane not only for the symmetric distribution (i.e., for $`f=1/2`$) but also for asymmetric distributions of the vehicles among the east-bound and north-bound streets .
$``$ Faulty traffic lights:
The effects of faulty traffic lights have been modeled by generalizing the BML as follows : a fraction $`f_{tl}`$ of the sites are identified randomly as the locations of the faulty traffic lights. Both the north-bound vehicles currently south of the faulty traffic lights and the east-bound vehicles currently west of a faulty traffic light are allowed to hop onto the empty crossings where the faulty traffic lights are located, irrespective of whether the corresponding time step of updating is odd or even. If an east-bound vehicle and a north-bound vehicle simultaneously attempt to enter the same crossing, where the faulty traffic light is located, then only one of then is allowed to enter that crossing by selecting randomly.
Since a north-bound (east-bound) vehicle will be able to move forward although only east-bound (north-bound) vehicles would have moved forward had $`f_{tl}`$ been zero, the average velocity of the vehicles is expected to increase with increasing $`f_{tl}`$. However, with the increase of $`f_{tl}`$ there is an increasing likelihood that an east-bound (north-bound) vehicle would be blocked by a north-bound (east-bound) vehicle at a faulty traffic light. Thus, the increasing density of faulty traffic lights increases the effect of grid-locking thereby decreasing $`c_{}`$.
$``$ Static hindrances:
The BML model has been extended to incorporate the effects of static hindrances or road blocks (e.g., vehicles crashed in traffic accident), i.e., stagnant points . A vehicle, which stays at a normal site for only one time step before attempting to move out of it, stays at a point-like blockage for $`T_p`$ time steps before attempting to move out of it. Obviously, the longer is $`T_p`$, the lower is the corresponding $`c_{}`$. The time-dependent phenomenon of spreading of the jam from the blockage site during the approach of the system to its jammed steady-state configurations has also been investigated .
$``$ Stagnant street:
Let us consider the effects of a stagnant street, where the local density $`c_s`$ of the vehicles is initially higher than that in the other streets , on the traffic flow in the BML model. The stagnant street, effectively, acts like a ”line-defect”, rather than a ”point-defect”. However, in contrast to the static roadblocks, a stagnant street offers a time-dependent hindrance to the vehicles moving in the perpendicular streets. As intuitively expected, the jamming transition has been found to occur at a lower global density when the local density in the stagnant street is higher.
$``$ Independent turning of the vehicles:
Let us now generalize the BML model by assigning a trend or preferred direction of motion, $`W_n(x,y)`$, to each vehicle $`n`$ ($`n=1,2,\mathrm{}`$) located at the site $`x,y`$. According to this definition, the vehicle $`n`$, located at $`x,y`$ jumps to the next site towards east with probability $`W_n(x,y)`$ while $`1W_n(x,y)`$ is the corresponding probability that it hops to the next site towards north . In this generalized model vehicles can take a turn but the processes of turning from east-bound (north-bound) to north-bound (east-bound) streets is stochastic. For simplicity, suppose, $`N/2`$ vehicles are assigned $`W_n(x,y)=\gamma `$ while the remaining $`N/2`$ vehicles are assigned $`W_n(x,y)=1\gamma `$ where $`0\gamma 1/2`$; this implies that $`N/2`$ vehicles move preferentially east-ward whereas the remaining $`N/2`$ vehicles move preferentially towards north. In the limit $`\gamma =0`$ no vehicle can turn and we recover the original BML model with deterministic update rules. The most dramatic effect of the stochastic turning is that the discontinuous jump $`\mathrm{\Delta }v`$ of the average velocity $`v`$ decreases with increasing $`\gamma `$ and, eventually, the first order jamming transition ends at a critical point where $`\mathrm{\Delta }v`$ just vanishes.
$``$ Jam-avoiding turn and drive:
In the model of turning considered in references a vehicle turns stochastically independently of the other vehicles. In real traffic, however, a vehicle is likely to turn if its forward movement is blocked by other vehicles ahead of it in the same street. Therefore, let us now consider a model where an east-bound (north-bound) vehicle turns north (east) with probability $`p_{turn}`$ if blocked by another vehicle in front of it. Computer simulations of this model shows that $`c_{}(p_{turn})`$ increases with increasing $`p_{turn}`$. In a slightly different model , on being blocked by a vehicle in front, an east-bound (north-bound) vehicle hops with probability $`p_{ja}`$ to the next east-bound (north-bound) street towards north (east).
$``$ A single north-bound street cutting across east-bound streets
Let us now consider a special situation where only one north-bound street exists which cuts all of the $`L`$ equispaced mutually parallel east-bound streets of length $`L`$ . This situation can be modeled as a $`L\times L`$ square lattice and each cell on the north-bound street can be in one of the three allowed states, namely, either empty or occupied by a $``$ or a $``$. But, in contrast to the BML model, there are two allowed states for each cell (outside the crossings) on the east-bound street; these can be either empty or occupied by only $``$. Since no grid-locking is possible with only one north-bound street, complete jamming occurs trivially in this case only if each of the cells either on the east-bound streets or on the north-bound street or on all the streets are occupied simultaneously by vehicles. Nevertheless, at any finite nonvanishing density, the crossings of the east-bound and north-bound streets act like hindrances with finite non-vanishing permeability for the flow of the east-bound traffic. Obviously, the higher is the density of the north-bound vehicles, the lower is the permeability and the stronger is the rate-limiting factor of the bottleneck. A comparison of this problem with that of highway traffic in the presence of static hindrances explains not only why the flux along the east-bound streets exhibits a flat plateau over an intermediate density regime but also why the plateau appears at lower values of the flux with increasing density of the north-bound vehicles . A mean-field theory for the macroscopic phase segregation in this model has been developed by appropriately modifying that for the similar phenomenon in the one-dimensional models of highway traffic in the presence of static hindrances. It is worth emphasizing here that in each site on the east-bound streets has been interpreted as a cell, which can accommodate one vehicle at a time, rather than as a crossing of the east-bound street with a north-bound street. Conceptually, this is an extension of the BML model. Finally, we briefly mention that in the phase diagram of a system consisting of one east-bound and one north-bound street with one crossing has been investigated.
$``$ Green-wave synchronization
Often the traffic lights along the main streets in cities are synchronized to allow continuous flow; this is usually referred to as ”green-wave” synchronization. A green-wave (GW) model has been developed by replacing the parallel updating scheme of the BML model by an updating scheme which is partly backward-ordered sequential (see Appendix A for a general explanation of this update scheme). At odd time steps, an east-bound vehicle moves by one lattice spacing if the target site was empty at the end of the previous time step or has become empty in the current time step (this is possible because of the backward-ordered sequential updating at every time-step). Similarly the positions of the north-bound vehicles are updated at every even time step. The main difference between the BML model and the GW model (see Fig. 53), arising from the different updating schemes, is that in the GW model vehicles move as ”convoys” (a cluster of vehicles with no empty cell between them) thereby mimicking the effects of green-wave synchronization of the traffic lights in real traffic. The jamming transition in the GW model has been investigated by a combination of a mean-field argument and numerical input from computer simulations .
$``$ More realistic description of streets and junctions
At first sight the BML model may appear very unrealistic because the vehicles seem to hop from one crossing to the next. However, it may not appear so unrealistic if each unit of discrete time interval in the BML model is interpreted as the time for which the traffic lights remain green (or red) before switching red (or green) simultaneously in a synchronized manner, and over that time scale each vehicle, which faces a green signal, gets an opportunity to move from $`j`$-th crossing to the $`(j+1)`$-th (or, more generally , to the $`(j+r)`$-th where $`r>1`$).
In the original version of the BML model the vehicles are located on the lattice sites which are identified as the crossings. Brunnet and Goncalves considered a modified version where, instead, the vehicles are located on the bonds and, therefore, never block the flow of vehicles in the transverse direction. Consequently, in this version of the CA model of city traffic jams of only finite sizes can form and these jams have finite lifetime after which they disappear while new jams may appear elsewhere in the system; an infinitely long-lived jam spanning the entire system is possible only in the trivial limit $`c=1`$. In contrast, Horiguchi and Sakakibara generalized the BML by replacing each of the bonds connecting the nearest-neighbour lattice sites by a bond decorated with an extra lattice site in between. In a generalization to $`s`$ extra lattice sites between crossings has been presented. However, the vehicles are still allowed to hop forward by only one lattice spacing in the model of Horiguchi and Sakakibara. Moreover, generalizing the rules for turning of the vehicles in Ref. , Horiguchi and Sakakibara also allowed probabilistic turning of the vehicles in their model. The model exhibits a transition from the flowing phase to a completely jammed phase.
The streets in the original BML model were assumed to allow only one-way traffic. This restriction was relaxed in a more realistic model proposed by Freund and Pöschel which allows both-way traffic on all the streets. Thus, each east-west (north-south) street is implicitly assumed to consist of two lanes one of which allows east-bound (north-bound) traffic while the other allows west-bound (south-bound) traffic. Moreover, each site is assumed to represent a crossing of a east-west street and a north-south street where four numbers associated with the site denote the number of vehicles coming from the four nearest-neighbour crossings (i.e., from north, south, east and west) and queued up at the crossing under consideration. So, in this extended version of the BML model, each site can accommodate at most $`4Q`$ particles if each of the four queues of vehicles associated with it is allowed to grow to a maximum length $`Q`$.
In the model proposed by Freund and Pöschel , initially, each of the vehicles is assigned a site selected at random, the queue to which it belongs (i.e., whether it is approaching the crossing from north,south, east or west) and the desired direction (i.e., left, right or straight) for its intended motion at the next time step. At each discrete time step a vehicle is allowed to move forward in its desired direction of motion by one lattice spacing provided (a) it is at the front of the queue in its present location and (b) there are fewer than $`Q`$ vehicles queued up at the next crossing in the same desired direction of motion. Once a vehicle moves to the next crossing it finds a place at the end of the corresponding queue at the new crossing while the vehicles in the queue it left behind are moved ”closer to the crossing” by one position (by mere relabeling as no physical movement of the vehicles in the queue takes place explicitly). Various reasonable choices for the rule which determines the desired direction of each vehicle at every time step have also been considered.
The finite space of the streets in between successive crossings do not appear explicitly in the extension of the BML model suggested by Freund and Pöschel although it is more realistic than the BML model because it implicitly takes into account the possibility of formation of queues by vehicles approaching one crossing from another. Chopard et al. have developed a more realistic CA model of city traffic where the stretches of the streets in between successive crossings appear explicitly. In this model also each of the streets consist of two lanes which allow oppositely directed traffic. The rule for implementing the motion of the vehicles at the crossing is formulated assuming a rotary to be located at each crossing. Depending on the details of the rules to be followed by the vehicles at the rotary, the system can exhibit a variety of phenomena. For example, the flow can be metastable at all densities if each of the vehicles on the rotary is required to stop till the destination cell becomes available for occupation . Moreover, the bottleneck created by the vehicles on the rotaries at the junctions can lead to a plateau in the fundamental diagram which is analogous to that caused by a static hindrances on a highway .
### 12.2 Marriage of the NaSch model and the BML model; a ”unified” CA model of city traffic
If one wants to develop a more detailed ”fine-grained” description of city traffic than that provided by the BML model then one must first decorate each bond with $`D1`$ ($`D>1`$) sites to represent $`D1`$ cells in between each pair of successive crossings thereby modeling each segment of the streets in between successive crossings in the same manner in which the entire highway is modelled in the NaSch model. Then, one can follow the prescriptions of the NaSch model for describing the positions, speeds and accelerations of the vehicles as well as for taking into account the interactions among the vehicles moving along the same street. Moreover, one should flip the color of the signal periodically at regular interval of $`T`$ ($`T1`$) time steps where, during each unit of the discrete time interval every vehicle facing green signal should get an opportunity to move forward from one cell to the next. Such a CA model of traffic in cities has, indeed, been proposed very recently where the rules of updating have been formulated in such a way that, (a) a vehicle approaching a crossing can keep moving, even when the signal is red, until it reaches a site immediately in front of which there is either a halting vehicle or a crossing; and (b) no grid-locking would occur in the absence of random braking.
Let us model the network of the streets as a $`N\times N`$ square lattice. For simplicity, let us assume that the streets parallel to $`\widehat{x}`$ and $`\widehat{y}`$ axes allow only single-lane east-bound and north-bound traffic, respectively, as in the original formulation of the BML model. Next, we install a signal at every site of this $`N\times N`$ square lattice where each of the sites represents a crossing of two mutually perpendicular streets. We assume that the separation between any two successive crossings on every street consists of $`D`$ cells so that the total number of cells on every street is $`L=N\times D`$. Each of these cells can be either empty or occupied by at most one single vehicle at a time. Because of these cells, the network of the streets can be viewed as a decorated lattice. However, unlike the BML model , which corresponds to $`D=1`$, and the model of Horiguchi and Sakakibara , which corresponds to $`D=2`$, $`D(<L)`$ in this model is to be treated as a parameter. Note that $`D`$ introduces a new length scale into the problem.
The signals are synchronized in such a way that all the signals remain green for the east-bound vehicles (and simultaneously, red for the north-bound vehicles) for a time interval $`T`$ and then, simultaneously, all the signals turn red for the east-bound vehicles (and, green for the north-bound vehicles) for the next $`T`$ timesteps. Clearly, the parameter $`T`$ introduces a new time scale into the problem. Thus, in contrast to the BML model, the forward movement of the individual vehicles over smaller distances during shorter time intervals are described explicitly in this ”unified” model.
If no turning of the vehicles is allowed, as in the original BML model, the total number of vehicles on each street is determined by the initial condition, and does not change with time because of the periodic boundary conditions. Following the prescription of the NaSch model, we allow the speed $`v`$ of each vehicle to take one of the $`v_{max}+1`$ integer values $`v=0,1,\mathrm{},v_{max}`$. Suppose, $`v_n`$ is the speed of the $`n`$-th vehicle at time $`t`$ while moving either towards east or towards north. At each discrete time step $`tt+1`$, the arrangement of $`N`$ vehicles is updated in parallel according to the following ”rules”:
Step 1: Acceleration.
$`v_n\mathrm{min}(v_n+1,v_{max})`$.
Step 2: Deceleration (due to other vehicles or signal).
Suppose, $`d_n`$ is the gap in between the $`n`$-th vehicle and the vehicle in front of it, and $`s_n`$ is the distance between the same $`n`$-th vehicle and the closest signal in front of it (see Fig. 54).
Case I: The signal is red for the $`n`$-th vehicle under consideration:
$`v_n\mathrm{min}(v_n,d_n1,s_n1)`$.
Case II: The signal is green for the $`n`$-th vehicle under consideration:
Suppose, $`\tau `$ is the number of the remaining time steps before the signal turns red. Now there are two possibilities in this case:
$`(i)`$ When $`d_ns_n`$, then $`v_n\mathrm{min}(v_n,d_n1)`$. The motivation for this choice comes from the fact that, when $`d_ns_n`$, the hindrance effect comes from the leading vehicle.
$`(ii)`$ When $`d_n>s_n`$, then $`v_n\mathrm{min}(v_n,d_n1)`$ if $`\mathrm{min}(v_n,d_n1)\times \tau >s_n`$; else $`v_n\mathrm{min}(v_n,s_n1)`$. The motivation for this choice comes from the fact that, when $`d_n>s_n`$, the speed of the $`n`$-th vehicle at the next time step depends on whether or not the vehicle can cross the crossing in front before the signal for it turns red.
Step 3: Randomization.
$`v_n\mathrm{max}(v_n1,0)`$ with probability $`p`$ ($`0p1`$); $`p`$, the random deceleration probability, is identical for all the vehicles and does not change during the updating.
Step 4: Vehicle movement.
For the east-bound vehicles, $`x_nx_n+v_n`$ while for the north-bound vehicles, $`y_ny_n+v_n`$.
The rule in case II of step 2 can be simplified without changing the overall behaviour of the model :
Case II’:
If the signal turns to red in the next timestep ($`\tau =1`$):
$`v_n\mathrm{min}(v_n,s_n1,d_n1)`$
else
$`v_n\mathrm{min}(v_n,d_n1)`$.
These rules are not merely a combination of the rules proposed by Biham et al. and those introduced by Nagel and Schreckenberg but also involve some modifications. For example, unlike all the earlier BML-type models, a vehicle approaching a crossing can keep moving, even when the signal is red, until it reaches a site immediately in front of which there is either a halting vehicle or a crossing. Moreover, if $`p=0`$ every east-bound (north-bound) vehicle can adjust speed in the deceleration stage so as not to block the north-bound (east-bound) traffic when the signal is red for the east-bound (north-bound) vehicles.
Initially, we put $`N_{}`$ and $`N_{}`$ vehicles at random positions on the east-bound and north-bound streets, respectively. We update the positions and velocities of the vehicles in parallel following the rules mentioned above. After the initial transients die down, at every time step $`t`$, we compute the average velocities $`v_x(t)`$ and $`v_y(t)`$ of the east-bound and north-bound vehicles, respectively.
A phase transition from the ”free-flowing” dynamical phase to the completely ”jammed” phase takes place in this model at a vehicle density $`c_{}(D)`$. The dependence on the dynamical parameters $`p`$, $`v_{max}`$ and $`T`$ is not clear yet . The data obtained so far from computer simulations do not conclusively rule out the possibility that the transition density only depends on the structure of the underlying lattice, similar to the percolation transition , and is independent of the dynamical parameters. The intrinsic stochasticity of the dynamics, which triggers the onset of jamming, is similar to that in the NaSch model, while the phenomenon of complete jamming through self-organization as well as the final jammed configurations (see Fig. 55) are similar to those in the BML model. The variations of $`v_x`$ and $`v_y`$ with time as well as with $`c`$, $`D`$, $`T`$ and $`p`$ in the flowing phase are certainly more realistic that in the BML model .
The ”unified” model has been formulated intentionally to keep it as simple as possible and at the same time capture some of interesting features of the NaSch model as well as the BML model. We believe that this model can be generalized (i) to allow traffic flow both ways on each street which may consist of more than one lane, (ii) to make more realistic rules for the right-of-the-way at the crossings and turning of the vehicles, (iii) to implement different types of synchronization or staggering of traffic lights (including green-wave), etc.
### 12.3 Practical applications of the models of vehicular traffic; On-line simulation of traffic networks
A large fraction of the available resources are spent by the governments, particularly in the industrialized developed countries, to construct more highways and other infrastructural facilities related to transportation. The car-following models, the coupled-map lattice models as well as the CA models have been used for computer simulation with a hope to utilize the results for on-line traffic control. For planning and design of the transportation network , for example, in a metropolitan area , one needs much more than just micro-simulation of how vehicles move on an idealized linear or square lattice under a specified set of vehicle-vehicle and road-vehicle interactions. For such a simulation, to begin with, one needs to specify the roads (including the number of lanes, ramps, bottlenecks, etc.) and their intersections. Then, times and places of the activities, e.g., working, shopping, etc., of individual drivers are planned. Micro-simulations are carried out for all possible different routes to execute these plans; the results give informations on the efficiency of the different routes and these informations are utilized in the designing of the transportation network . Some socio-economic questions as well as questions on the environmental impacts of the planned transportation infrastructure also need to be addressed during such planning and design. For a thorough discussion of these aspects we refer to the recent review by Nagel et al. .
## 13 Some related systems, models and phenomena
In this section we will briefly present some stochastic models and phenomena which are somehow related to the main topic of this review. For some models the relation to traffic flow is rather obvious, e.g. there are toy models which share some features with traffic models, but which can be solved exactly. We also point out similarities in the description of other phenomena, e.g. granular flow or surface growth. Similarities also exist to systems from solid state physics, namely ionic conductors. This point will not discussed here, instead we refer to refs. . Finally we like to mention that there are some resemblances between the CA models of vehicular traffic and the CA models of driven diffusive Frenkel-Kontorova-type systems .
### 13.1 Generalizations of the TASEP
The TASEP is probably one of the best studied models in nonequilibrium physics. Using powerful methods like MPA or Bethe Ansatz recently it became possible to calculate not only simple expectation values in the stationary states, but also more complicated quantities like diffusion constants or the large-deviation function. For a more detailed discussion and a list of references we refer to the recent reviews on the TASEP (see also ). Several variants of the TASEP have been proposed. Most of those preserve the exact solvability. In the following we will discuss briefly some of the variants and generalizations of the TASEP discussed so far in the literature.
$``$ Partially asymmetric exclusion process
An obvious generalization of the TASEP is to allow hopping processes in both directions . Here only results for the random-sequential update exist, since for parallel dynamics ambiguities in the updating appear when two particles attempt to hop onto the same site. One finds that the phase diagram looks essentially like that of the TASEP with three different phases (see Sec. 8.5). Recently it has been shown , however, that the high- and low-density phases can be divided into three subphases (AI, AII, AIII and BI, BII, BIII) instead of two in the TASEP case. Again the phase boundaries of these subphases are determined by the behaviour of the density profiles and the corresponding correlation lengths.
$``$ Multispecies models I
Karimipour and collaborators have developed a multispecies generalization of the TASEP which retains the solvability by MPA. It is similar to the disordered model discussed in Sec. 10.1 where each particle is characterized by a hopping rate (also called “velocity” in this context) $`v_j`$ ($`j=1,\mathrm{},N_s`$), i.e. there are $`N_s`$ different “species”’ of particles. In contrast to the models with quenched disorder discussed in Sec. 10, however, overtaking of particles is possible, i.e. the ordering of the particles is not fixed.
More specifically the dynamics of the model is defined as follows: A particle is chosen at random (random-sequential update). If the particle is of species $`j`$ and the cell to its right is empty, it hops there with rate $`v_j`$. If the cell is occupied by a particle of species $`l`$ and $`v_l<v_j`$, then they interchange there positions with rate $`v_jv_l`$. This means that a fast particle can overtake a slower one with a rate proportional to their velocity difference.
$``$ Multispecies models II
In several papers multispecies generalizations of the ASEP have been suggested which exhibit phase separation and spontaneous symmetry breaking.
Arndt et al. considered a system of positive and negative charged particles diffusing on a ring in opposite directions. Positive particles move to an empty right neighbour and negative particles move to an empty left neighbour with the same rate $`\lambda `$. Furthermore positive and negative particles on neighbouring sites can exchange their positions. The process $`++`$ occurs with rate 1, and the inverse process $`++`$ with a different rate $`q`$. For equal densities of positive and negative particles the system exhibits three phases. For $`q<1`$ in the thermodynamic limit, the system organizes itself into configurations consisting of blocks of the type $`000\mathrm{}+++\mathrm{}\mathrm{}`$ The dynamics out of theses states is extremely slow. Translational invariance is broken and the current vanishes. This phase is called ’pure phase’. For $`1<q<q_c`$ the system is in the ’mixed phase’ which consists of two coexisting phase, the dense phase and the fluid phase. The dense phase where no vacancies exist covers a macroscopic region which shrinks to 0 for $`qq_c`$. It is remarkable that the current for $`1<q<q_c`$ takes the value $`J=(q1)/4`$ independent of the total density and the rate $`\lambda `$. For $`q>q_c`$ the fluid phase extends through the whole system. There is no charge separation and density profiles are uniform.
Similar results have been found by Evans et al. in a slightly different model. The dynamics of their 3-species model<sup>21</sup><sup>21</sup>21One species may be interpreted as a vacant site. is given by
(134)
i.e. the rates are cyclic in $`A`$, $`B`$ and $`C`$ and the numbers $`N_A`$, $`N_B`$ and $`N_C`$ of particles of each species is conserved.
In the case $`N_A=N_B=N_C`$ the dynamics satisfies the detailed balance condition with respect to a Hamiltonian with long-range asymmetric interactions. Stationary states are of the form $`A\mathrm{}AB\mathrm{}BC\mathrm{}C`$ and exhibit phase separation, i.e. for large separations $`r`$ the two-point function satisfies $`lim_L\mathrm{}\left[A_1A_rA_1A_r\right]>0`$.
### 13.2 Surface growth, KPZ equation and Bethe Ansatz
In Sec. 8.1.2 we have explained how the NaSch model with $`v_{max}=1`$ can be mapped onto a stochastic surface growth model. This connection can be employed to calculate several properties of the noisy Burgers and KPZ equation exactly.
Gwa and Spohn used the Bethe Ansatz (see e.g. ) to determine the spectrum of the stochastic Hamiltionian (see App. F)
$$=\frac{1}{4}\underset{j=1}{\overset{L}{}}\left[𝝈_j𝝈_{j+1}1+iϵ\left(\sigma _j^x\sigma _{j+1}^y\sigma _j^y\sigma _{j+1}^x\right)\right]$$
(135)
corresponding to the ASEP with random-sequential updating and periodic boundary conditions. $`𝝈_j=(\sigma _j^x,\sigma _j^y,\sigma _j^z)`$ are the standard Pauli matrices at site $`j`$ and $`ϵ`$ is the asymmetry of the hopping rates, $`q_{right}=\frac{1}{2}(1+ϵ)`$ and $`q_{left}=\frac{1}{2}(1ϵ)`$.
The “ground state” of $``$ has eigenvalue $`0`$ and is $`L`$-fold degenerate. For a fixed number $`N`$ of up-spins<sup>22</sup><sup>22</sup>22i.e. a fixed number $`N`$ of particles every configuration has equal weight in the ground state. In order to determine the dynamical scaling exponent of the noisy Burgers and KPZ equations, Gwa and Spohn investigated the finite-size behaviour of the energy gap of $``$. Since $``$ is non-hermitian its spectrum is complex. The first excited state is then defined as the eigenvalue with the smallest (positive) real part $`E_{gap}`$. In it was shown that $`E_{gap}L^{3/2}`$ for $`ϵ=1`$ and $`N=L/2`$. This implies that the dynamical exponent $`z`$ for the stationary correlations of the KPZ equation is given by $`z=3/2`$. The dynamical exponent relates temporal and spatial scaling behaviour on large scale. Generalizations and related results can be found in .
By an extension of the Bethe Ansatz method of Gwa and Spohn, Derrida and Lebowitz calculated the large deviation function (LDF) of the time-averaged current of the TASEP. The LDF is related to the total displacement $`Y(t)`$, i.e. the total number of hops to the right minus the total number of hops to the left between time $`0`$ and time $`t`$. In the corresponding growth model $`Y(t)`$ is the total number of particles deposited until time $`t`$. The LDF is then defined as
$$f(y)=\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{ln}\left[Prob\left(\frac{Y(t)}{t}=\overline{v}+y\right)\right]$$
(136)
where $`\overline{v}=lim_t\mathrm{}Y(t)/t=ϵc(1c)L^2/(L1)`$ is the mean current for a ring of finite size $`L`$ and density $`c=N/L`$. The results of have been extended and generalized in .
Apart from the treatment of finite systems, the BA can also be used to solve the master equation for an infinite system with a finite number of particles . This allows e.g. to study the collective diffusion of two single particles.
### 13.3 Protein synthesis
You must have noticed in the earlier sections that some of the models of traffic are non-trivial generalizations or extensions of the TASEP, the simplest of the driven-dissipative systems which are of current interest in nonequilibrium statistical mechanics . Some similarities between these systems and a dynamical model of protein synthesis have been pointed out .
In a simplified picture of the mechanism of biopolymerisation ribosomes read the genetic information encoded in triplets of base pairs. They attach to one end of a messenger-RNA and then move along the chain after adding a monomer to a biopolymer attached to the ribosome. The type of monomer added depends on the genetic information read by the ribosome. When the ribosome reaches the other end of the m-RNA the biopolymer is fully synthesized and the ribosome is released.
MacDonald et al. have described the kinetics of this process using an ASEP-type model. The m-RNA is represented by a chain of $`L`$ sites where each site corresponds to one triplet of base pairs. The ribosome is given by a hard-core particle covering $`r`$ neighbouring sites ($`r2030`$) which moves forward by one site with rate $`q`$. At the beginning of the chain particles are added with rate $`\alpha q`$ and at the end they are released with rate $`\beta q`$. In the idealized case $`r=1`$ this is exactly the TASEP of Sections 8.1.1 and 8.5. The relevant case for the experiments is $`\alpha =\beta <1/2`$. The exact solution of the TASEP allowed for an explanation of many aspects of the experiments .
### 13.4 Granular flow
Another quasi-one-dimensional driven-dissipative system, which is also receiving wide attention of physicists in recent years, is the granular material flowing through a pipe . Since the fascinating phenomena (e.g. size segregation, convection, standing waves, localized excitations) found in granular materials have been subject of several excellent reviews we discuss only briefly the similarities between the clustering of vehicles on a highway and particle-particle (and particle-cluster) aggregation process .
Obviously both highway traffic and granular flow through a pipe may be described as quasi-onedimensional systems consisting of discrete elements (vehicles, grains). The dynamics of these elements is determined by an intricate interplay between a driving force (driver, gravitation) and dissipation (braking, inelastic scattering processes). These similarities already show that both systems can be described by similar approaches. One important difference between traffic flow and granular flow exist, however. In granular flow density waves can move both in and against the direction of the flow whereas in traffic flow they only move backwards.
An important success of a description of granular flow using the optimal velocity model (Sec. 6.2) is the explanation of the experimentally observed $`f^{4/3}`$-behaviour of the power spectrum .
### 13.5 The bus route model
The bus route model (BRM) has been formulated as a one-dimensional lattice with periodic boundary conditions. The sites represent cells, each of which may be thought of as a bus stop and are labeled by an index $`i`$ ($`i=1,2,\mathrm{},L`$) . Two binary variables $`\sigma _i`$ and $`\tau _i`$ are assigned to each cell $`i`$: (i) If the cell $`i`$ is occupied by a bus then $`\sigma _i=1`$; otherwise $`\sigma _i=0`$. (ii) If cell $`i`$ has passengers waiting for a bus then $`\tau _i=1`$; otherwise $`\tau _i=0`$. Since a cell cannot have simultaneously a bus and waiting passengers, let us impose the condition that a cell cannot have both $`\sigma _i=1`$ and $`\tau _i=1`$ simultaneously. Each bus is assumed to hop from one stop to the next.
Next, let us specify the update rules (see Fig. 57): a cell $`i`$ is picked up at random. Then, (i) if $`\sigma _i=0`$ and $`\tau _i=0`$ (i.e, cell $`i`$ contains neither a bus nor waiting passengers), then $`\tau 1`$ with probability $`\lambda `$, where $`\lambda `$ is the probability of arrival of passenger(s) at the bus stop. (ii) If $`\sigma _i=1`$ (i.e., there is a bus at the cell $`i`$) and $`\sigma _{i+1}=0`$, then the hopping rate $`\mu `$ of the bus is defined as follows: (a) if $`\tau _{i+1}=0`$, then $`\mu =\alpha _b`$ but (b) if $`\tau _{i+1}=1`$, then $`\mu =\beta _b`$, where $`\alpha _b`$ is the hopping rate of a bus onto a stop which has no waiting passengers and $`\beta _b`$ is the hopping rate onto a stop with waiting passenger(s). Generally, $`\beta _b<\alpha _b`$, which reflects the fact that a bus has to slow down when it has to pick up passengers. We can set $`\alpha _b=1`$ without loss of generality. When a bus hops onto a stop $`i`$ with waiting passengers $`\tau _i`$ is reset to zero as the bus takes all the passengers. Note that the density of buses $`c=N/L`$ in a conserved quantity whereas that of the passengers is not.
An ideal situation in this bus-route model would be one where the buses are evenly distributed over the route so that each bus picks up roughly the same number of passengers. However, because of some fluctuation, a bus may be delayed and, consequently, the gap between it and its predecessor will be larger than the average gap. Therefore, this bus has to pick up more passengers than what a bus would do on the average, because during the period of delay more passengers would be waiting for it and, as a result, it would get further delayed. On the other hand, the following bus has to pick up fewer passengers than what a bus would do on the average and, therefore, it would catch up with the delayed bus from behind. The slowly moving delayed bus would slow down the buses behind it thereby, eventually, creating a jam. In other words, once a larger-than-average gap opens up between two successive buses, the gap is likely grow further and the steady state in a finite system would consist of a single jam of buses and one large gap. This is very similar to the Bose-Einstein-condensation-like phenomenon we have observed earlier in particle-hopping models with slow impurities. On the basis of heuristic arguments and mean-field approximation it has been argued that this model exhibits a true phase transition from an inhomogeneous low-density phase to a homogeneous (but congested) high-density phase only in the limit $`\lambda 0`$. Finally we mention that the BRM with parallel dynamics has been recently been studied in where also its connection with the NaSch model has been elucidated.
### 13.6 Mobile directional impurities
We have considered the effects of random hopping probabilities, assigned either to the lattice sites or to the particles in the TASEP and in the NaSch-type models of vehicular traffic, on the nature of the corresponding steady states as well as their approach to the steady-states starting from random initial conditions. Toroczkai and Zia solved analytically a model with one ”mobile directional impurity”; this model is also an extension of the ASEP. In this model, $`N`$ particles, labeled by integers $`1`$ to $`N`$ (from left to right), occupy the sites of a one-dimensional lattice of length $`N+1`$ where periodic boundary conditions are applied. Thus, there is a single ”hole” (i.e., empty site) in this model. The shifting of the hole from the site in between the particles $`n`$ and $`n+1`$ to the site in between particles $`n+1`$ and $`n+2`$ is described by the statement ”hole jump from position $`n`$ to $`n+1`$”. In the absence of any impurity, the hole at position $`n`$ can exchange position with either the particle on its left (with probability $`W_{n1,n}`$) or the particle on its right (with probability $`W_{n+1,n}`$). These probabilities are arbitrary and direction-dependent (i.e., in general, $`W_{n1,n}W_{n+1,n}`$) but time-independent. Note that the hopping probabilities of the hole is determined by the particles (rather than the lattice sites) in front and behind it. The general case, where $`W_{n1,n}W_{n+1,n}`$ and the probabilities $`W_{n1,n}`$ as well as $`W_{n+1,n}`$ for different $`n`$ are chosen randomly, is referred to as the random asymmetric case. As the hole wanders, the string of particles also shifts as a whole. However, this system is translationally invariant in the sense that the jump rate of any particular particle-hole pair is independent of its location on the lattice.
The ”directional impurity” is introduced by identifying a specific bond (whose position is fixed with respect to the lattice) as a defect bond such that the time-independent rates of particle-hole exchanges across it are fixed at, say, $`q`$ and $`q^{}`$ irrespective of the particle-hole pair involved. In other words, when the hole is in between the particles $`n`$ and $`n+1`$, the hopping probability to the right (left) is always $`W_{n+1,n}`$ ($`W_{n1,n}`$), except when the defect bond is involved. This model can be represented as in Fig. 58 where the defect bond is shown as a kink; the motivation for such a kink came from an earlier model of gel electrophoresis .
More specifically, suppose, the probabilities for the exchange of the hole with a particle across the defect bond are $`q`$ if the hole moved upward and $`q^{}`$ if the hole moved downward, independent of the particle involved. In the so-called pure limit, we have, for all $`n`$, $`W_{n+1,n}=W_{}`$ and $`W_{n1,n}=W_{}`$. In this limit, the hole can also be regarded as a particle undergoing biased diffusion (if $`W_{}W_{}`$) everywhere except across across a specific defect bond. If we make the further assumption that no backward motion of any particle is allowed, then this model reduces to the model of TASEP with a single blockage .
### 13.7 Computer networks
Inspired by the recent success of the methods of statistical mechanics outside the traditional domain of physics, tools of statistical mechanics have also been applied to analyze fundamental properties of information traffic on the international network of computers (Internet) . Messages in the form of information packets are continuously being emitted from the hundreds of millions of host computers and transported to their destination computers through this network. Each of these packets is relayed through the so-called routers on its way. The routers can deal with the packets one by one. Each router has a finite buffer where the arriving packets get queued up and forwarded one by one from the head of the remaining queue to their respective next destinations. Since packets run with the velocity of light through the cables, information congestion does not take place inside the cables. It is the routers which give rise to the information congestion on the internet. Measuring the fluctuations in the round-trip time taken by a message on the internet (using the ping command of the UNIX operating system), $`1/f`$-like power spectrum has been observed .
In the square lattice model of a computer network developed by Ohira and Sawatari information packets are generated at the sites on the boundary at a rate $`\lambda `$ with the corresponding destination addresses chosen randomly from among the boundary sites. The packets can form queues of unlimited length at the the inner nodes, which act as routers of the network. At every time step, the packets from the heads of the queues at the routers are forwarded to the tail ends of the queues at the next router. Both deterministic and probabilistic strategies have been considered for selecting the next router to which the individual packets are to be forwarded. On reaching their individual destinations the packets die. The average number of time steps between the birth and death of a packet is referred to as the average lifetime of a packet. Computing the average lifetime as a function of the birth rate $`\lambda `$ of the packets, Ohira and Sawatari observed a transition from a low-congestion phase to a high-congestion phase at a non-zero finite value $`\lambda _{}`$.
## 14 Summary and conclusion
In this section we summarize our conclusions regarding the current status of understanding of the statistical physics of vehicular traffic. We also speculate on the future trends of research in this area.
As we stated in the beginnning, one of the main aims of basic research on vehicular traffic, from the point of view of statistical physics, is to understand the nature of the steady-states of the system. We have summarized the empirical evidences available at present in support of the occurrence of three distinct dynamical phases, namely, those corresponding to free-flow, synchronized flow and stop-and-go traffic. Our critical review of the theoretical works has made it clear that, at present, the physical mechanisms at the ”microscopic” level, which give rise to the synchronized traffic, are not as well understood as those responsible for the free-flow and stop-and-go traffic. There are, however, strong indications that for a complete theory, which would account for all these three phases, one must take into account not only the vehicles on a given stretch of the highway but also on the on- and off-ramps.
We have explored the possibility of transitions from one dynamical phase to another in the NaSch model (and TASEP) with periodic boundary conditions. Moreover, we have also presented the generic phase diagram of the TASEP with open boundary conditions and explained the notion of boundary-induced phase transitions in such one-dimensional driven-diffusive lattice gases which are far from equilibrium. Furthermore, we have found that while some models exhibit first-order phase transitions, some others exhibit second order phase transitions and the signatures of ”criticality” while in some rare situations, e.g., in the cruise-control limit of the NaSch model, the system is found to exhibit even ”self-organized” criticality.
In the beginning, we stated that one of the aims of basic research on vehicular traffic is to understand the nature of the dynamical fluctuations around the steady-states. The time-dependent correlations functions and the distributions of the relaxation times have been computed for the NaSch model and some other models but the general questions of the validity of dynamic scaling and dynamic universality classes have not been addressed. Another aim of the statistical mechanical approach to vehicular traffic, as we stated also in the beginning, is to investigate how the system evolves from initial states which are far from the corresponding steady-state. The phenomenon of ”coarsening” of the platoons of vehicles during evolution from random initial states have been studied in some models. But, the questions of ”universality”, if any, of the growth exponents have not been addressed so far. Metastable states have been observed in several CA models. But, to our knowledge, the mechanisms of spontaneous decay of such states (analogue of homogeneous nucleation) has not been investigated so far. Besides, to our knowledge, so far it has not been possible to develope any powerfull analytical technique for calculating the dynamical properties of the traffic models.
While stating the aims of basic research on vehicular traffic, we also mentioned the need to understand the effects of quenched disorder on the steady-states as well as on the dynamical properties of the systems. We have seen that the randomization of the hopping probability of the vehicles can lead to some exotic platooning phenomena which are close analogs of the ”Bose-Einstein-condensation”. An alternative prescription for introducing quenched disorder into the traffic models is to install random bottlenecks on the road and assign a time-independent hopping probability (or, equivalently, a ”permeability”) for hopping across bonds in such locations. It has been found that such localized bottlenecks can lead to global phase-segregation.
The NaSch model is the most extensively studied minimal CA model of vehicular traffic on idealized single-lane highways. We have explained the conceptual framework, and illustrated the use of the mathematical formalism, of the cluster-theoretic analytical calculations for the NaSch model. This formalism, which yields exact results for the NaSch model with $`v_{max}=1`$, gives quite accurrate estimates of various quantities of interest as long as $`v_{max}`$ is not too large. It would be desirable to develope a new formalism to carry out exact analytical calculations for higher velocities too. Here also the limit $`v_{max}=\mathrm{}`$ is interesting since it shows a rather peculiar behaviour.
In addition to the detailed discussions on the NaSch model and its various generalizations, we have also mentioned briefly some other alternative CA models of single-lane highway traffic, e.g., the VDR model, the Fukui-Ishibashi model, the Werth-Froese-Wolf model, etc. so that one can appreciate the ongoing efforts to formulate the most satisfactory minimal model.
From the point of view of practical applications, modelling vehicular traffic on multi-lane highways are more relevant than that on idealized single-lane highways which are, nevertheless, interesting from the point of fundamental understanding of truly non-equilibrium phenomena in driven-diffusive lattice gases. At present, there are several different alternative prescriptions for formulating the CA rules for lane-changing of the vehicles on multi-lane highways. But, in order to pick out the most appropriate one from among these CA models theorists would require input from careful further observations of the phenomenon of lane-changing on real multi-lane highways. Empirical observations may also indicate modifications or extensions of the CA rules necessary for more realistic modelling of the multi-lane traffic.
The generalizations of the CA models of traffic on idealized single-lane highways to those on multi-lane highways may be regarded as extensions of one-dimensional model chains to one-dimensional strips. The BML model of vehicular traffic in cities may be regarded as a further generalizations of these models from one-dimensional chains to two-dimensional lattices, or further, to decorated lattices. A few different CA rules have been considered so far for taking into account the effects of the traffic lights at the crossings of streets in such idealized street networks. We have emphasized the intrinsic differences between percolation clusters and the cluster of jammed vehicles in the BML (and similar) models in spite of some apparent similarities between them.
We have focused attention mainly on the progress made in the recent years using ”particle-hopping” models, formulated in terms of cellular automata, and compared these with several other similar systems. Although this may be a slightly biased overview (as all reviews usually are) of the theory of vehicular traffic, we have also discussed the main ideas behind all the major approaches including the fluid-dynamical, gas-kinetic and car-following theories of vehicular traffic. At present the relationships between different approaches of modelling have not explored in great detail. It would be very useful if the phenomenological parameters of the macroscopic theories can be estimated by utilizing the mathematical formulae relating these with those of the ”microscopic” models.
It is now quite clear that, in order to make significant further progress, we not only need more realistic models and better techniques of calculation but we also need more detailed and accurate empirical data from real traffic on highways as well as more careful re-analysis of the existing data in the light of recently developed concepts. So far as the observation of the real traffic is concerned, a lot can be learnt from a systematic analysis of aereal pictures or video photographs. Alternatively, as the second best choice, a series of counting loops along the highway can give more insight, by providing detailed information on, for example, time-headways (flux), velocity and local density.
If you are a critical thinker (or a pragmatist) you may ask: ”armed with the theoretical tools at our disposal now, can we predict the occurrence of a traffic jam at a specific place on a given highway (or street) at a particular instant of time”? This question sounds similar to questions often asked in the context of some other interdisciplinary topics of current research in the area of complex systems, e.g., ”can we predict an earthquake”, or, ”can we predict a stock market crash”? Of course, we know that, at present, the best we can hope for is to predict (if at all possible) probabilities of occurrences of all these phenomena. But, we must admit that, we have a long way to go before we come even close to this goal. Nevertheless, we hope you have enjoyed the fascinating twists and turns of the way we have covered so far. Our endeveour will be more successful if your interest has been stimulated by the intellectual challenges posed by the open problems and if you are willing to uncover the current mysteries as well as anticipating new surprises that may lie ahead. We are just at the beginning of a long road!
”The volume of vehicular traffic in the past several years has rapidly outstripped the capacities of the nation’s highways. It has become increasingly necessary to understand the dynamics of traffic flow and obtain a mathematical description of the process” \- H. Greenberg (1959)
Acknowledgements: It is our pleasure to thank P. Arndt, R. Barlovic, M. Barma, J.G. Brankov, E. Brockfeld, G. Diedrich, B. Eisenblätter, J. Esser, K. Ghosh, N. Ito, J. Kertész, K. Klauck, W. Knospe, S. Krauss, J. Krug, D. Ktitarev, A. Majumdar, K. Nagel, L. Neubert, A. Pasupathy, V. Popkov, V.B. Priezzhev, N. Rajewsky, M. Schreckenberg, G. Schütz, S. Sinha, D. Stauffer, R.B. Stinchcombe, Y. Sugiyama, P. Wagner, D.E. Wolf and J. Zittartz for enjoyable collaborations the results of some of which have been reviewed here, for useful discussions and for critical comments as well as suggestions on a preliminary draft of this review. One of us (DC) acknowledges warm hospitality of ICTP, Trieste during the preparation of this manuscript. This work is supported by SFB341 Köln-Aachen-Jülich.
## Appendix A Definition of update orders
A dynamical model is not fully defined just by its local transition rules. In addition one has to specify the order in which the rules are applied to the different particles, i.e. the update ordering (sometimes also called ’dynamics’). This is an essential part of the definition of the model since the transient and even the stationary state may differ dramatically .
For the NaSch model one uses a parallel update scheme where the rules are applied to all particles (i.e., vehicles) at the same time. This kind of ordering is sometimes also called synchronous updating.
Among the various types of asynchronous update schemes most frequently the the so-called random-sequential update is used. Here one picks one particle at random and applies the transition rules to it. Then one makes another random choice (which can also be the same particle again) and so on. This update is sometimes called continuous since it can be described by a master equation in continuous time.
Apart from the parallel update there are other updates which are discrete in time. We just mention the ordered-sequential updates. Here one starts by applying the transition rules to one particle. After that the rules are applied to the other particles in a fixed order, e.g. one might continue with the next particle ahead of the first one (forward-ordered) or the next particle behind it (backward-ordered). We would like to point out that, in principle, one has to distinguish two different types of ordered-sequential updates which one could name site-ordered-sequential and particle-ordered-sequential, respectively. In contrast to the particle-ordered-sequential update described above, in the site-ordered-sequential update the rules are applied to all sites consecutively. This might have a strong effect, since a particle might move to next cell ahead which then is updated next (for the forward-particle-ordered-sequential update). Then this particle might move again and so on. This is different from the particle-ordered-sequential case where a particle at most moves once during a sweep through the lattice. As an example consider the extreme case of the NaSch mode with $`v_{max}=1`$ where only one particle is present which moves with probability $`q=1`$ to an empty cell in front. This particle will move through the whole lattice during one sweep! By looking at a lattice with two particles, one can already see that the two different updates might introduce rather different correlations. Starting with particles separated by $`d`$ empty sites, in the site-ordered-sequential update the left particle will move to the right until it reaches the right particle, which then starts to move. On the other hand, in the case of particle-ordered-sequential update the particles will stay always $`d`$ or $`d1`$ sites apart. For general values of $`q`$ the situation is similar.
There are several other updates which can be defined. We refer to the literature (see e.g. ) for a comparison of different update procedures. The parallel update usually produces the strongest correlations and is used for traffic simulations . Note that the forward-particle-ordered-sequential update is almost identical to the parallel update. In the case of periodic boundary conditions a difference only occurs during the update of the last particle. In the forward-ordered case the particle in front of it (i.e. the first particle) might already have moved since it has been updated earlier. Although this difference appears to be minor it can have a large effect. The difference between parallel and forward-particle-ordered update can be viewed as a dynamical defect.
## Appendix B TASEP
This simple model of driven systems of interacting particles is one of the most exhaustively studied prototype models in nonequilibrium statistical mechanics . This model can be divided into four classes on the basis of the boundary conditions and the update scheme for the implementation of the dynamics. In this appendix we consider the TASEP with only random-sequential dynamics.
Let us consider the TASEP with periodic boundary conditions and random-sequential dynamics. Since only two states, namely empty and occupied, are allowed for each site we can use a two-state variable $`n_i`$ to denote the state of the $`i`$-th site where $`n_i=0`$ if the $`i`$-th site is empty and $`n_i=1`$ if the $`i`$-th site is occupied. For any given initial configuration $`\{n_i(0)\}`$, we can write the equations governing the time evolution of $`n_i(t)`$ (and all the correlation functions) by taking into account all the processes during the elementary time interval $`dt`$. It is not difficult to establish that
$$n_i(t+dt)=\{\begin{array}{cc}n_i(t)\hfill & \text{with probability }12dt\hfill \\ n_{i1}(t)+n_i(t)n_{i1}(t)n_i(t)\hfill & \text{with probability }dt\hfill \\ n_i(t)n_{i+1}(t)\hfill & \text{with probability }dt\hfill \end{array}$$
(137)
$$\frac{dn_i}{dt}=n_{i1}(1n_i)n_i(1n_{i+1})$$
(138)
which, upon averaging over the history between times $`0`$ and $`t`$, leads to the equation
$$\frac{dn_i}{dt}=n_{i1}n_in_{i1}n_i+n_in_{i+1}$$
(139)
for $`n_i`$, the average occupation of the $`i`$-th site. Note that the equation for $`n_i`$ involves two-site correlations. Similarly, it is straightforward to see that the equations for the two-site correlations involve three-site correlations and so on. Thus, the problem is an intrinsically $`N`$-body problem! The probability distribution for this system in the steady-state is given by
$$P_{steadystate}(\{n_i\})=\frac{N!(LN)!}{L!}$$
(140)
where $`L`$ and $`N`$ refer to the total number of sites and the total number of particles, respectively. From this distribution it follows that $`n_i=N/L`$ and $`v=\frac{LN}{L1}`$ which lead to $`n_i=c`$ and $`v=(1c)`$ in the thermodynamic limit.
The stationary state of the TASEP with open boundary conditions and random-sequential dynamics has been determined exactly using the so-called Matrix Product Ansatz (MPA) (see Appendix F for a more technical introduction) in and in using recursion relations. This solution has been generalized to different types of discrete dynamics in . The solution for parallel dynamics was obtained recently in and using generalizations of the MPA technique.
## Appendix C Naive site-oriented mean-field treatment of the NaSch model
Suppose, $`c_v(i,t)`$ Probability that there is a vehicle with speed $`v`$ ($`v=0,1,2,\mathrm{},v_{max}`$) at the site $`i`$ at the time step $`t`$. Then, obviously, $`c(i,t)=_{j=0}^{v_{max}}c_j(i,t)`$ Probability that the site $`i`$ is occupied by a vehicle at the time step $`t`$ and $`d(i,t)=1c(i,t)`$ is the corresponding probability that the site $`i`$ is empty at the time step $`t`$. Using the definition
$$J(c,p)=\underset{v=1}{\overset{v_{max}}{}}vc_v.$$
for the flux $`J(c,p)`$ one can get the mean-field fundamental diagram for the given $`p`$ provided one can get $`c_v`$ in the mean-field approximation.
Step I: Acceleration stage ($`tt_1`$)
$`c_0(i,t_1)`$ $`=`$ $`0,`$ (141)
$`c_v(i,t_1)`$ $`=`$ $`c_{v1}(i,t),(0<v<v_{max})`$ (142)
$`c_{v_{max}}(i,t_1)`$ $`=`$ $`c_{v_{max}}(i,t)+c_{v_{max}1}(i,t)`$ (143)
Step II: Deceleration stage ($`t_1t_2`$)
$`c_0(i,t_2)`$ $`=`$ $`c_0(i,t_1)+c(i+1,t_1){\displaystyle \underset{v=1}{\overset{v_{max}}{}}}c_v(i,t_1)`$ (144)
$`c_v(i,t_2)`$ $`=`$ $`c(i+v+1,t_1){\displaystyle \underset{j=1}{\overset{v}{}}}d(i+j,t_1){\displaystyle \underset{v^{}=v+1}{\overset{v_{max}}{}}}c_v^{}(i,t_1)`$ (145)
$`+c_v(i,t_1){\displaystyle \underset{j=1}{\overset{v}{}}}d(i+j,t_1),(0<v<v_{max})`$
$`c_{v_{max}}(i,t_2)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{v_{max}}{}}}d(i+j,t_1)c_{v_{max}}(i,t_1)`$ (146)
Step III: Randomization stage ($`t_2t_3`$)
$`c_0(i,t_3)`$ $`=`$ $`c_0(i,t_2)+pc_1(i,t_2)`$ (147)
$`c_v(i,t_3)`$ $`=`$ $`qc_v(i,t_2)+pc_{v+1}(i,t_2),(0<v<v_{max})`$ (148)
$`c_{v_{max}}(i,t_3)`$ $`=`$ $`qc_{v_{max}}(i,t_2)`$ (149)
Step IV: Movement stage ($`t_3t+1`$)
$$c_v(i,t+1)=c_v(iv,t_3),(0vv_{max})$$
(150)
In the special case $`v_{max}=1`$ the equations (141-150) get simplified and, hence, we get
$`c_0(i,t+1)`$ $`=`$ $`c(i,t)c(i+1,t)+pc(i,t)d(i+1,t),`$ (151)
$`c_1(i,t+1)`$ $`=`$ $`qc(i1,t)d(i,t).`$ (152)
Similarly, in the case of $`v_{max}=2`$, one gets
$`c_0(i,t+1)`$ $`=`$ $`[c(i,t)+pd(i,t)]c_0(i,t)+[1+pd(i,t)]c(i,t)[c_1(i,t)+c_2(i,t)],`$
$`c_1(i,t+1)`$ $`=`$ $`d(i,t)qc_0(i,t)+d(i,t)[qc(i,t)+pd(i,t)][c_1(i,t)+c_2(i,t)],`$
$`c_2(i,t+1)`$ $`=`$ $`qd^2(i,t)[c_1(i,t)+c_2(i,t)].`$ (155)
## Appendix D Paradisical mean-field theory
For $`v_{max}=1`$ the question, whether a state is a GoE state or not, can be decided locally by investigating just nearest-neighbour configurations. By analysing the update rules one finds that all states containing the local configurations $`(0,1)`$ or $`(1,1)`$, i.e. configurations where a moving vehicle is directly followed by another car, are GoE states. This is not possible as can be seen by looking at the previous configurations. The momentary velocity gives the number of cells that the car moved in the previous timestep. In both configurations the first car moved one cell. Therefore it is immediately clear that $`(0,1)`$ is a GoE state since otherwise there would have been a doubly-occupied cell before the last timestep. The configuration $`(1,1)`$ is also not possible since both cars must have occupied neighbouring cells before the last timestep too. Therefore, according to step 2, the second car could not move.
The MFT equations (64) and (65) have to be modified to take into account the existence of GoE states. In general, one has to follow the procedure outlined in Appendix C. A quicker way to derive the paradisical mean-field (pMF) equations is to analyse the MF equations (64) and (65). In (64) the contribution $`c(i;t)c(i+1;t)`$ appears. Since we know that site $`i+1`$ can never be occupied by a car with velocity $`1`$ if site $`i`$ is not empty, this contribution has to be modified to $`c(i;t)c_0(i+1;t)`$ in pMFT. All other contributions are left unchanged compared to MFT.
Due to this modification and the corresponding reduction of the configuration space the normalization $`c_0+c_1=c`$ is no longer satisfied automatically. Therefore a normalization constant $`𝒩`$ has to be introduced. The final equations for a homogeneous stationary state are than given by
$`c_0`$ $`=`$ $`𝒩(c_0+pd)c,`$ (156)
$`c_1`$ $`=`$ $`𝒩qcd,`$ (157)
with the normalization
$$𝒩=\frac{1}{c_0+d}.$$
(158)
Since $`c_0+c_1=c`$ we have only one independent variable for fixed density $`c`$, e.g. $`c_1`$. Solving (156), (157) for $`c_1`$ we obtain
$$c_1=\frac{1}{2}\left(1\sqrt{14q(1c)c}\right).$$
(159)
The flow is given by $`f(c)=c_1`$ and we recover the exact solution for the case $`v_{max}=1`$.
In the case $`v_{max}=2`$ more GoE states exist. In order to identify these it is helpful to note that the rules step 1 – step 4 imply $`d_j(t)=d_j(t1)+v_{j+1}(t)v_j(t)`$ and therefore
$`d_j(t)`$ $``$ $`v_{j+1}(t)v_j(t),`$ (160)
$`v_j(t)`$ $``$ $`d_j(t1)`$ (161)
The second inequality (161) is a consequence of step 2.
In the following we list the elementary GoE states, i.e. the local configurations which are dynamically forbidden (cars move from left to right):
$`(0,1),(0,2),(1,2),(0,,2),`$ (162)
$`(1,1),(2,1),(2,2),(1,,2),(2,,2),`$ (163)
$`(0,,,2).`$ (164)
The numbers give the velocity of a vehicle in an occupied cell and $``$ denotes an empty cell.
The elementary GoE states in (162) violate the inequality (160), and the configurations in (163) violate (161). The state in (164) is a second order GoE state. Going one step back in time leads to a first order GoE state since $`(0,,,2)`$ must have evolved from $`(0,v)`$ (with $`v=1`$ or $`v=2`$).
Again we can derive the pMF equations by modifying the method for the derivation of the MFT. Taking into account only the first order GoE states (162) and (163) one obtains the following pMF equations:
$`c_0`$ $`=`$ $`𝒩\left[c_0c+pd(c_0+c_1c)\right],`$ (165)
$`c_1`$ $`=`$ $`𝒩\left[pd^2(c_1+c_2)+qd(c_0+c_1c)\right],`$ (166)
$`c_2`$ $`=`$ $`𝒩qd^2(c_1+c_2).`$ (167)
The normalization $`𝒩`$ ensures $`c_0+c_1+c_2=c`$ and is given explicitly by
$$𝒩=\frac{1}{c_0+dc_1+d^2c_2}=\frac{1}{c_0+d(1c_2)}.$$
(168)
These equations have been analysed in . After expressing $`c_2`$ through $`c_0`$ and $`c`$ by
$$c_2=\frac{1}{2d}\left(c_0+d\sqrt{(c_0+d)^24qd^3(cc_0)}\right).$$
(169)
Inserting this result into (165) we obtain a cubic equation which determines $`c_0`$ in terms of the parameters $`c`$ and $`p`$ . Results for different values of $`p`$ are shown in Fig. 19. These results are only slightly modified when also the second order GoE state is taken into account .
## Appendix E Equations of car-oriented theory of NaSch Model and COMF approximation
In terms of $`p`$, $`q`$, $`g`$ and $`\overline{g}(t)=1g(t)`$, the equations describing the time evolution of the probabilities $`P_n(t)`$ for the NaSch model with $`v_{max}=1`$ are given by
$`P_0(t+1)`$ $`=`$ $`\overline{g}(t)[P_0(t)+qP_1(t)],`$ (170)
$`P_1(t+1)`$ $`=`$ $`g(t)P_0(t)+[qg(t)+p\overline{g}(t)]P_1(t)+q\overline{g}(t)P_2(t),`$ (171)
$`P_n(t+1)`$ $`=`$ $`pg(t)P_{n1}(t)+[qg(t)+p\overline{g}(t)]P_n(t)+q\overline{g}(t)P_{n+1}(t)(n2).`$
It is worth mentioning here that, for the NaSch model with $`v_{max}=1`$, the 2-cluster probabilities $`P_2(\sigma _i,\sigma _j)`$ of the the SOMF theory are related to the probabilities $`P_n`$ of the COMF theory through
$`P_2(1,1)`$ $`=`$ $`cP_0,`$ (173)
$`P_2(1,0)`$ $`=`$ $`c(1c)P_1,`$ (174)
$`P_2(0,0)`$ $`=`$ $`(1c){\displaystyle \frac{P_{n+1}}{P_n}}(n1)`$ (175)
## Appendix F The matrix-product Ansatz for stochastic systems
For the stochastic systems considered here the time-evolution of the probability $`P(𝝉,t)`$ to find the system in the configuration $`𝝉=(\tau _1,\mathrm{},\tau _L)`$ is determined by the master equation. For random-sequential dynamics it has the form
$$\frac{P(𝝉,t)}{t}=\underset{\stackrel{\mathbf{~}}{𝝉}}{}w(\stackrel{\mathbf{~}}{𝝉}𝝉)P(\stackrel{\mathbf{~}}{𝝉},t)\underset{\stackrel{\mathbf{~}}{𝝉}}{}w(𝝉\stackrel{\mathbf{~}}{𝝉})P(𝝉,t)$$
(176)
with transition rates $`w(\stackrel{\mathbf{~}}{𝝉}𝝉)`$ from state $`\stackrel{\mathbf{~}}{𝝉}`$ to state $`𝝉`$. Eq. (176) can be rewritten in the form of a Schrödinger equation in imaginary time ,
$$\frac{}{t}|P(t)=|P(t),$$
(177)
with the state vector $`|P(t)=_𝝉P(𝝉,t)|𝝉`$. The vectors $`|𝝉=|\tau _1,\mathrm{},\tau _L`$ corresponding to the configurations $`𝝉`$ form an orthonormal basis of the configuration space. The stochastic Hamiltonian $``$ is defined through its matrix elements
$$𝝉||\stackrel{\mathbf{~}}{𝝉}=w(\stackrel{\mathbf{~}}{𝝉}𝝉),𝝉||𝝉=\underset{𝝉\stackrel{\mathbf{~}}{𝝉}}{}w(𝝉\stackrel{\mathbf{~}}{𝝉})(𝝉\stackrel{\mathbf{~}}{𝝉}).$$
(178)
The stationary state of the stochastic process corresponds to the eigenvector $`|P_0`$ of $``$ with eigenvalue $`0`$.
For discrete-time dynamics the master equation takes the form
$$P(𝝉,t+1)=\underset{\stackrel{\mathbf{~}}{𝝉}}{}W(\stackrel{\mathbf{~}}{𝝉}𝝉)P(\stackrel{\mathbf{~}}{𝝉},t)$$
(179)
where $`W(\stackrel{\mathbf{~}}{𝝉}𝝉)=w(\stackrel{\mathbf{~}}{𝝉}𝝉)\mathrm{\Delta }t`$ are transition probabilities. This can be rewritten as
$$|P(t+1)=𝒯|P(t).$$
(180)
Here the stationary state corresponds to the eigenvector $`|P_0`$ of the transfer matrix $`𝒯`$ with eigenvalue $`1`$.
A very powerful method for the determination of stationary solutions of the master equation is the so-called matrix-product Ansatz (MPA). For a system with open boundaries the weights $`P(𝝉)`$ in the stationary state can be written in the form
$$P(\tau _1,\mathrm{},\tau _L)=\frac{1}{Z_L}W|\underset{j=1}{\overset{L}{}}\left[\tau _jD+(1\tau _j)E\right]|V.$$
(181)
For periodic boundary condition the MPA takes the form
$$P(\tau _1,\mathrm{},\tau _L)=\mathrm{Tr}\left(\underset{j=1}{\overset{L}{}}\left[\tau _jD+(1\tau _j)E\right]\right).$$
(182)
For simplicity we have assumed a two-state system where e.g. $`\tau _j=0`$ corresponds to an empty cell $`j`$ and $`\tau _j=1`$ to an occupied cell. $`Z_L`$ is a normalization constant that can be calculated as $`Z_L=W|C^L|V`$. In (181), (182) $`E`$ and $`D`$ are matrices and $`W|`$ and $`|V`$ are vectors characterizing the boundary conditions. The explicit form of these quantities has to be determined from the condition that (181) or (182) solves the master equation. This leads in general to a algebraic relations between the matrices $`E`$ and $`D`$ and the boundary vectors $`W|`$ and $`|V`$. Once one these have been determined one has a simple recipe for determining $`P(\tau _1,\mathrm{},\tau _L)`$: First, translate the configuration $`\tau _1,\mathrm{},\tau _L`$ into a product of matrices by identifying each empty cell ($`\tau _j=0`$) with a factor $`E`$ and each occupied cell ($`\tau _j=1`$) with $`D`$. In that way the configuration $`011001\mathrm{}`$ corresponds to the product $`EDDEED\mathrm{}=ED^2E^2D\mathrm{}`$. The weight of the configuration is then just the matrix element with the vectors $`w|`$ and $`|v`$.
A simple example is the ASEP discussed in Sect. 8.1.1. Here these quantities have to satisfy
$`pDE`$ $`=`$ $`D+E,`$ (183)
$`\alpha W|E`$ $`=`$ $`W|,`$ (184)
$`\beta D|V`$ $`=`$ $`|V.`$ (185)
If one is able to find explicit representations for this algebra one can determine in principle all expectation values in the stationary state exactly. For (183)–(185) one can show that all representations are infinite-dimensional . Only on the line $`\alpha +\beta =p`$ one-dimensional representations (with $`E,D`$ and $`W|`$, $`|V`$ being real numbers) are possible<sup>23</sup><sup>23</sup>23This can be seen easily from (183)–(185)..
At this point it might appear that the MPA works only in very special cases. However, it can be shown that the stationary state of one-dimensional stochastic processes is generically of matrix-product form <sup>24</sup><sup>24</sup>24We have to mention here that these results up to now do not include the case of parallel dynamics.. Even if it is not straightforward to find general representations of the resulting algebras, one can at least search systematically for finite-dimensional representations on special lines in the parameter space of the model. Furthermore, since the mathematical structure of the stationary state is known it is sometimes possible to derive rather general results. As an example in interesting relations between expectation values for ordered-sequential and sublattice-parallel dynamics have been derived.
For a more detailed description of the MPA for different types of dynamics and its relation with the MPA technique for quantum-mechanical spin systems we refer to . A review of the treatment of the ASEP using the MPA is given in . The MPA has also been extended to treat the full dynamics, not only the stationary state . Using time-dependent matrices $`D(t)`$ and $`E(t)`$ one obtains the Bethe Ansatz equations for the corresponding stochastic Hamiltonian .
## Appendix G Two schemes for solving the mean-field approximation of the DTASEP
A mean-field approximation scheme for this model has also been developed . The time-averaged steady-state current $`J_{j,j+1}`$ in the bond $`(j,j+1)`$ is given by $`J_{j,j+1}=q_{j,j+1}n_j(1n_{j+1})`$. In the mean-field approximation, $`n_j(1n_{j+1})=n_j(1n_{j+1})`$ and, hence,
$$J=J_{j,j+1}=q_{j,j+1}c_j(1c_{j+1})$$
(186)
where $`c_j=n_j`$. In order to calculate the steady-state flux $`J`$ as a function of the mean density $`c`$ of the particles, Tripathy and Barma used two different iteration schemes based on the equation (186).
(i) Constant-current iteration scheme: In this scheme, for a given system length $`L`$ and a fixed flux $`J=J_0`$, one starts with some value of $`c_1`$ and, computes all the other $`c_j`$ $`(j>1)`$ using the equation (186), i.e.,
$$c_{j+1}=1\frac{J_0}{q_{j,j+1}c_j},j=1,2,\mathrm{},L$$
(187)
together with the periodic boundary condition $`c_{j+L}=c_j`$. If the iteration converges, i.e., one gets all the site densities in the physically acceptable range $`[0,1]`$, one accepts the average of these final site densities to be the global mean density of the particles corresponding to the flux $`J_0`$. (ii) Constant-density iteration scheme: In this scheme, for a given system length $`L`$ and fixed global average density $`c`$, one begins by assigning the site densities $`0c_j(0)1`$ to the lattice sites subject to the global constraint $`\frac{1}{L}_jc_j(0)=c`$. Then, the site densities are updated in parallel according to
$$c_j(t+1)=c_j(t)+J_{j1,j}(t)J_{j,j+1}(t),j=1,2,\mathrm{},L$$
(188)
which follows from the equation (186). It is straightforward to verify that the this iteration scheme keeps the average global density $`c`$ unchanged at every step of updating and hence the name. After sufficient number of iterations the set of densities converge to a set $`\{c_j\}`$ and the flux on each bond converge to the steady-state flux $`J_0`$.
## Appendix H Self-consistent equations for $`v_x`$ and $`v_y`$ in the mean-field approximation of the BML model
Suppose, $`v_x`$ and $`v_y`$ denote the average speeds of east-bound and north-bound vehicles, respectively. Then, on the average, an east-bound vehicle spends a time $`1/v_x`$ at a site whereas a north-bound vehicle spends a time $`1/v_y`$ at a site. The north-bound vehicles lead to a reduction of the speed of the east-bound vehicles by $`n_y/v_y`$. Moreover, because of the hindrance of the east-bound vehicles by other east-bound vehicles ahead of it there will be further reduction of the speed of the east-bound vehicles by $`n_x[\frac{1}{v_x}1]`$. Furthermore, if the density of the overpasses is $`f_o`$, then
$$v_x=1(1f_o)\left[\frac{n_y}{v_y}+n_x\left(\frac{1}{v_x}1\right)\right]$$
(189)
Similarly, the corresponding equation for $`v_y`$ is given by
$$v_y=1(1f_o)\left[\frac{n_x}{v_x}+n_y\left(\frac{1}{v_y}1\right)\right]$$
(190)
In the special case $`n_x=n_y=n/2`$ both the equations (189) and (190) reduce to the form
$$v=1(1f_o)\left[\frac{1}{v}\frac{1}{2}\right]n$$
(191)
where $`v_x=v_y=v`$. The solution of the quadratic equation (191) for $`v`$ is
$$v=\frac{1}{2}\left[1+\frac{1f_o}{2}n+\sqrt{\left(1+\frac{1f_o}{2}n\right)^24(1f_o)n}\right]$$
(192)
## Appendix I Derivation of the equations in the microscopic theory of the BML model
By definition, $`n_{}(x,y;t)`$ ($`n_{}(x,y:t)`$) is unity if the site $`(x,y)`$ is occupied at time $`t`$ by a north-bound (east-bound) vehicle and zero if the site $`(x,y)`$ is not occupied by a north-bound (east-bound) vehicle. Normalization requires $`n_{}(x,y;t)+n_{}(x,y;t)=1n_{empty}(x,y;t)`$ where the two-state variable $`n_{empty}(x,y;t)`$ is unity if, at time $`t`$, the site $`(x,y)`$ is empty and zero if the site $`(x,y)`$ is not empty. In order to describe the state of the signals at time $`t`$, one also defines a two-state variable $`S(t)`$: $`S(t)=1(0)`$ if the signal is green (red) for the vehicles under consideration. The space-average of $`n_{}(x,y;t)`$ and $`n_{}(x,y;t)`$ are $`c_{}(t)`$ and $`c_{}(t)`$, respectively. Besides, the time-average of $`S(t)`$ is $`1/2`$.
The updating rules of the BML model lead to the equations
$`n_{}(x,y;t+1)=n_{}(x,y;t)\left[n_{}(x+1,y;t)+n_{}(x+1,y;t)\right]S(t)`$
$`+n_{}(x1,y;t)\left[1n_{}(x,y;t)n_{}(x,y;t)\right]S(t)+n_{}(x,y;t)[1S(t)]`$
(193)
The first term on the right hand side of equation (193) describes the situation when the east-bound vehicle, which was at the site $`(x,y)`$ at time $`t`$, finds a green signal but cannot move because the next site towards east is occupied by another vehicle. The second term on the right hand side of (193) corresponds to the situation where the east-bound vehicle, which was at the site $`(x1,y)`$ at time $`t`$, finds a green signal and moves to the next site towards east, which was empty. The last term on the right hand side of equation (193) arises from the possibility that the east-bound vehicle, which was at $`(x,y)`$ at time $`t`$, could not move because of a red signal, irrespective of the state of occupation of the next site towards east. Following the similar arguments, one can also write down the corresponding equation for the north-bound vehicles.
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# Effective description of quark mixing
## Acknowledgements
It is a pleasure to thank J.A. Aguilar-Saavedra, F. Cornet, J.L. Cortés and J. Prades for discussions. JS thanks the Dipartimento di Fisica “Galileo Galilei” and INFN sezione di Padova for their hospitality. This work has been supported by CICYT and Junta de Andalucía. MPV and JS also thank MECD for financial support.
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# The Liar-paradox in a Quantum Mechanical Perspective11footnote 1Appeared in Foundations of Science, 4(2), 1999, 115-132.
## 1 Introduction.
Our approach, to analyze the Liar paradox as a cognitive entity, emerges naturally from research on integrating world views (Worldviews Group, 1994, 1995; Broekaert, 1999). The justification is based on their necessary inclusion in an encompassing model. An integrating world view can be based on a model of interacting layers epistemically corresponding to the various contemporary sciences (CLEA,1997). Basically, in the same way as we view social entities in a social layer, or quantum entities in the quantum — which in this context can be viewed as the pre-material — layer, we place cognitive entities in a cognitive layer of reality. Which is considered the expanse where the personal and interpersonal cognitive interactions are taking place, and which will be elaborated in the next section.
In this paper, we will work with versions of the Liar Sentence with index or sentence pointers followed by the sentence this index points at. Typically we will either concentrate on the single sentence version;
‘Single Liar’
(1) sentence (1) is false
or the two sentence version;
‘Double Liar’
(1) sentence (2) is false
(2) sentence (1) is true
First, we stress that it is not our intention to construct a solution for the paradox, but we will concentrate in this paper on a model of the self-referential circularity — more precisely, the truth-value dynamics — involved. In this way, we are able to understand the nature of the liars’ paradoxicality by modelling the typical truth-value oscillation in the encompassing framework. By taking this position there is no need to go into details of ‘solving’-strategies that have been proposed in the literature. Actually, the most important of these strategies are specifically constructed to avoid the paradox by, for example, declaring it ‘ungrammatical’, or ‘meaningless’ or ‘ungrounded’. All these strategies, of which we only mentioned some, and those that introduced truth value ‘gaps’ or ‘gluts’, have proven valuable, and serve here to distinguish the present position. For more details on the mentioned solving strategies, see for example (Grim 1991, chp 1).
Finally, we mention some issues concerning the recourse taken to the quantummechanical formalism in describing cognitive entities. Contrary to the present paper, many attempts have been made to approach ‘emergent’ phenomena — like cognition — through complex dynamics. Focusing on the collective dynamics of quite different entities, some very similar patterns are encountered regardless of the nature of the substrate. This has been pointed out in research pertaining to the theories of nonlinearity and complexity. Because of the efficacy of these theories, tentative conclusions about the nature of — quite disparate — processes in reality can be inferred from structural-dynamical similarities. Although, we have to stress that these similarities generally involve only so called ‘noncontextual’ entities, as the theories of nonlinearity and complexity are fundamentally classical theories. By noncontextuality we mean that the entity will ‘not’ be influenced by an act of observation (measurement of the entity). A direct consequence of the fact that the theories of nonlinearity and complexity are classical deterministic. It is now our claim that cognitive entities — especially of the type of the liar paradox — share with quantum entities a variable, and sometimes high, contextual nature: i.e. also cognitive entities, as it is the case for quantum entities, are influenced by the act of measurement, which in the case of a cognitive entity is generated by the cognitive interaction. Therefore the theory of quantum mechanics can be used in analyzing the nature of cognitive entities and more specifically the nature of the liar paradox. In (Aerts, Broekaert, Smets, 1999) we built a quantum mechanical model that fits this purpose. With respect to the high contextual nature of intricate cognitive and also social entities, we can point out an approach where the probability model that results in an opinion pole — with important influence of the interviewer on the interviewee — is of a quantum mechanical nature (Aerts,1998; Aerts and Aerts,1995,1996; Aerts, Coecke and Smets,1999). In a similar way, we encounter an important contextual influence of an observer reading (observing) the liar paradox sentences. In this case the high contextuality will occur through the observer — the cognitive person — reasoning through the self-referent entity.
## 2 Cognitive Entities
The existence of a cognitive entity is recognized by its aptitude of being generally and cognitively influenced on as a practically stable configuration, e.g. in the reasoning on it and communicating about it, and by the limited number of different states that it can be in. The cognitive entity is endowed with properties and relations with the other elements of its layer. Their internal coherence relates to language as well as conceptions of experience. The extent to which it is related to pendants in the physical and other layers enhances its identity and its granted coincidence with the physically real entity. This variable correspondence relation is put forward between entities of the ontological cognitive layer and their pendants in physical or other layers.
In order to give a more precise and complete characterization of the cognitive entity, we need to reconstruct its interaction profile. This is done by specifying the different kinds of measurements or experiments appropriate to the intended entity. This necessarily happens with the intervention of the cognitive person — the conscious human being conceiving the cognitive entity.
We distinguish between those entities that refer to elements of layers other than the cognitive one and entities that do not. Arguably, such a distinction may only be possible in specific phases of cognition and theory forming, respecting personal psychological evolution. In both cases, the influencing or intervening of the cognitive person will be different; engendering physical interventions or only cognition, reasoning.
We define such an intervention as a kind of measurement on the entity with the aim of gaining information or increasing knowledge on its precise state. In order to keep our analogy with the physical layer, we restrict our exposition here to measurements that characterize the state of a cognitive entity.
The state formalizes specific — or all — possible actualisations according the interaction profile. The exhaustive description being practically impossible, an idealized state description can only be obtained for the overall state. We mention however that this problem appears in an analogous way in physics, where in each model also only an idealized state description is obtained, depending on the possible experiments that are available. Hence, a specific property state corresponds, as in physics, to appropriate measurements that are available. In the case of the liar paradox entity, the exact meaning of the truth-state will be explained in the next section. The measurements that characterize the state of an entity referring to the classical-material layer are mostly straightforward. For example, analyzing if a piece of chalk is breakable has a fixed measurement-procedure namely, breaking the piece of chalk. This gives the corresponding entity a classical behavior in the sense that each time we perform the measurement we obtain, with certainty, the same result. In a quantum mechanical context we say that such entities are in an ‘eigenstate’ for the corresponding measurement. Also cognitive entities that do not refer to other layers can have straightforward measurement-procedures. Especially in the case of e.g. a mathematical theorem, there is often a straightforward procedure to establish a proof for it. And indeed each time we apply the same procedure, we obtain the same result, so we say again that those entities are in an eigenstate corresponding to this measurement procedure. Naturally, we do not expect everybody to know these mathematics, so the state of that same theorem will depend on what we call the ‘cognitive background’ of each cognitive person. Even mathematicians will have different opinions on the status of some theorems, therefore it is clear that the general cognitive background of people, over history, is important and plays a role in the nature of the states that will be attributed to cognitive entities. As this example points out, the nonfixed character of knowledge implies we accord states of nonfixed character to cognitive entities. Again however, we remark that this is analogous to the physical situation. The state of a physical entity, defined correspondingly to measurement procedures, depends also on the general body of experimental possibilities at a certain epoch of time in the history of human culture, and hence is not fixed once and forever. Depending on whether this state has been defined by more ‘universal’ methods of experimentation, the state will approach in a deeper way the ontology of the entity. The same holds true for the state of a cognitive entity.
In the case of the liar paradox as a cognitive entity, for any measurement procedure, the truth-value will give us different results each time we intervene with it. Here, the state describing the truth-value will be called a superposition state as related to the measurement procedure in question, in analogy with the quantum mechanical concept. This means that we can not obtain a ‘certain’ prediction — in a classical sense — of properties, in this case ‘thruth’ of ‘false’, of that cognitive entity. The input of the quantummechanical formalism is therefore appropriate.
How does one measure the liar paradox? The measurement here consists of two part-processes, ’reading the sentence’ and ‘making a sentence true or false’. This means that in our description the liar paradox within the cognitive layer of reality is ‘in general’ — before the measurement — not in a state such that a reading would give true or false. The ‘true state ’ and the ‘false state ’ of the sentence are specific states; ‘eigenstates’ of the measurement. In general, the state of the liar paradox is not one of these two eigenstates. Due to the act of measurement, and in analogy with what happens during a quantum measurement, the state of the sentence changes (‘collapses’) into one of the two possible eigenstates, the ‘true state’ or the ‘false state’. This act of making a sentence true or false can be specifically described as ‘read it and make an hypothesis about its truth or falsehood’. In the next section we will apply this approach to the Double Liar, and see that an initial measurement followed by the sequence of logical inferences puts into work an oscillation dynamics that we can describe by a Schrödinger evolution over reasoning-time. We will also see in the next section, that the change of state due to measurement can be described by a projection operator in the quantum mechanical Hilbert space where the Schrödinger evolution is defined.
Once we make an hypothesis about one sentence the whole entity starts changing from one truth-state into another by continued reading with logical inference. When we stop this process, by means of not ‘looking at’ or ‘reasoning on’ the sentences any more, the entity — we hypothesise — reestablishes its original superposition state of indefiniteness. This superposition state does not correspond to what has been called a truth value gap and does not assign a third truth value to the liar. In this sense, in our model, the set of semantical truth values, of which only one can be assigned to the cognitive entity, does not contain a third truth value or a value gap.
Finally, we remark some interpretative issues concerning the origin of the entity’s dynamics and the nature of the cognitive layer. From the quantummechanical analogy, the temporal evolution of the entity is expected to originate intrinsically, still the construction of the evolution — as will be clear from the next section — supposes the cognitive person’s motivation by reasoning. We interpret the latter to be reflected in the autonomous dynamics, as such, the cognitive entity obtains its cognitive essence. The precise origin of temporal evolution has thereby become less transparent; the entity as well as the cognitive person will engender identical evolution.
The nature of the cognitive layer, is essentially different from the sphere spanned by common material objects. In our approach — akin to ‘effective’ realism — the cognitive layer is ‘Hilbert-space like’, a personal and mental construct with social and cultural conditioning, a collective dynamic emergent layer carried by its cognitive participators. The extent and subtlety of this issue, allows in the present context merely explanatory simplifications. More detailed elaborations of the cognitive layer as an emergent, Hilbert-space like sphere in social groups are due (Aerts, Broekaert and Gabora, 1999).
## 3 The Liar-paradox: A Quantum Description of its Truth Behavior.
We will first discuss the Single Liar entity. By ‘measuring’ the single sentence we attribute a chosen truth-value to the sentence, immediately and logically inferring from its lecture the opposite truth-value. The cognitive person is therefore inclined to attribute in an alternating manner opposite truth values to the Single Liar sentence, until the ‘measuring’ process is chosen to be stopped. Subsequently no decisive and unambiguous truth-value can be attributed to the Single Liar sentence as it is. The application of the quantummechanical formalism suggests to describe this situation by a superposition of opposite truth-value states. We remark the striking correspondence between truth values and the two-fold eigenvalues of a spin-1/2 state of some quantum particles (e.g. an electron), and the oscillatory dynamics present in the reasoning dynamics respectively the evolution dynamics of a spin-1/2 particle in a constant magnetic field. This formal correspondence will be used to construct a dynamical representation of the cognitive entity. Recall that this correspondence is possible, due to the fact that the spin of a particle is a quantized property, it exposes itself by means of distinct spin values. For a spin-1/2 particle there are only two distinct spin values namely ‘up’ and ‘down’ — appropriately corresponding to ‘true’ and ‘false’ values for the cognitive entity. The quantum mechanical formalism enables to express a superposition of these ‘up’ and ‘down’ states of a spin-1/2 particle by simple addition. When we apply the same idea to the single liar sentence, we obtain a state $`\mathrm{\Psi }`$ described by a pondered superposition of the two states of opposite truth-value:
$$\mathrm{\Psi }=c_{true}\left(\begin{array}{c}1\\ 0\end{array}\right)+c_{false}\left(\begin{array}{c}0\\ 1\end{array}\right)$$
The measurement itself, namely the interaction of the cognitive person on the entity when the sentence is being made true or being made false is described respectively by the true-projector $`P_{true}`$ or false-projector $`P_{false}`$.
$$P_{true}=\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)P_{false}=\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)$$
In this quantum mechanical description, the true-measurement (false measurement) on the superposed state $`\mathrm{\Psi }`$ results in the true state (resp. false state). The true state is represented as follows :
$$P_{true}\mathrm{\Psi }=c_{true}\left(\begin{array}{c}1\\ 0\end{array}\right)$$
where the square modulus of the corresponding pondering factor $`c_{true}`$ gives the statistical probability of finding the entity in the true-state. An unequivocal result is therefore not obtained when the superposition does not leave out one of the states completely, i.e. either $`c_{true}`$ or $`c_{false}`$ is zero. Only in those instances do we unambiguously attribute to a sentence its truth or falsehood.
We now look at the Double Liar sentences in more formal and mathematical detail. We consider three situations:
$$\mathrm{A}\{\begin{array}{cc}(1)\hfill & \mathrm{sentence}(2)\mathrm{is}\mathrm{false}\hfill \\ (2)\hfill & \mathrm{sentence}(1)\mathrm{is}\mathrm{true}\hfill \end{array}$$
$$\mathrm{B}\{\begin{array}{cc}(1)\hfill & \mathrm{sentence}(2)\mathrm{is}\mathrm{true}\hfill \\ (2)\hfill & \mathrm{sentence}(1)\mathrm{is}\mathrm{true}\hfill \end{array}$$
$$\mathrm{C}\{\begin{array}{cc}(1)\hfill & \mathrm{sentence}(2)\mathrm{is}\mathrm{false}\hfill \\ (2)\hfill & \mathrm{sentence}(1)\mathrm{is}\mathrm{false}\hfill \end{array}$$
From the case of the Single Liar we expect here a representation of the truth-behaviour by coupled <sup>2</sup> vectors, one for each sentence. Closer inspection of the coupled sentences of the two-sentence liar paradox of type (B) and (C), shows this is possible. The measurement of the Double Liar (B) allways will couple true-states of (B1) and (B2), and false-states of (B1) and (B2). In the case (C) on the other hand, measurement will couple the true-state of (B1) to the false-state of (B2), and the false-state of (B1) to the true-state of (B2). From a formal point of view, the equivalent spin-states in quantum mechanics would be described by the so called ‘singlet state’ and a ‘triplet state’ respectively. The singlet state indicates that two spin-1/2 particles are anti-alined and in an anti-symmetrical state , while a triplet state indicates that two spin-1/2 particles are alined and in a symmetrical state.
Whereas the Single Liar has been mathematically represented in a <sup>2</sup> finite dimensional complex Hilbert Space, we now need a $`{}_{}{}^{2}`$<sup>2</sup> space for the description of the (B) or (C) Double Liar. The tensorproduct $``$ connects the two sentences into one composed entity.
In the specific case of (C), and taking into account the anti-symmetric spin analog, $`\mathrm{\Psi }`$ is written as:
$$\frac{1}{\sqrt{2}}\left\{\left(\begin{array}{c}1\\ 0\end{array}\right)\left(\begin{array}{c}0\\ 1\end{array}\right)\left(\begin{array}{c}0\\ 1\end{array}\right)\left(\begin{array}{c}1\\ 0\end{array}\right)\right\}$$
Still, in the application to the cognitive entity other choices of the pondering coefficients are possible. The only constraints on the coefficients are: equal amplitude and addition of the squared amplitudes to unity.
The state-vector for the liar paradox in case (B) can be constructed in a similar manner:
$$\frac{1}{\sqrt{2}}\left\{\left(\begin{array}{c}1\\ 0\end{array}\right)\left(\begin{array}{c}1\\ 0\end{array}\right)+\left(\begin{array}{c}0\\ 1\end{array}\right)\left(\begin{array}{c}0\\ 1\end{array}\right)\right\}$$
The projection operators which make sentence one and respectively sentence two true are now:
$$P_{1,true}=\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)\mathrm{𝟏}_2P_{2,true}=\mathrm{𝟏}_1\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)$$
The projection operators that make the sentences false are obtained by switching the elements $`1`$ and $`0`$ on the diagonal of the matrix. These four projection operators represent the possible ‘logical’ interactions between the cognitive person and the cognitive entity.
During the continued measurement on the entities (B) and (C), the sequence of logical inferences results in a repetitive pattern of consecutive true-false states. These patterns are not very complicated, in case of entity (B) it will be a repetition of true-states (resp. false-states) depending on whether we presupposed an initial true (resp. false) state. While in the case of entity (C) it will always be an alternation between true-states and false states, no matter which state we presupposed.
Finally, we describe in detail the original double liar paradox, case (A). In this case we will show how the true-false cycle originates from the Schrödinger time-evolution of the appropriate initial state.
Instead of working within the coupled Hilbert space $`{}_{}{}^{2}`$<sup>2</sup>, as in cases (B) and (C), we have to use a space of higher dimension for (A). This complexification is due to the fact that no initial state can be found in the restricted space $`{}_{}{}^{2}`$<sup>2</sup>, such that application of the four true-false projection operators results in four orthogonal states respectively representing the four truth-falsehood states. The existence of such a superposition state — with equal amplitudes of its components — is required to describe the entity prior to, and after, any measurement procedure. If we perform the continued measurement on (A), by consecutive logical inference, the dynamical pattern is not anymore a two-step process like in the previous cases (B) and (C),instead we have a four-step process. Starting from the initial superposition state this four-step process can not be described by the coupled spin-1/2 models any more. In order to resolve this problem, recourse has to be taken to a 4 dimensional Hilbert-space for each sentence. The Hilbert-space needed to describe the Double Liar (A) is therefore $`{}_{}{}^{4}`$<sup>4</sup>.
The initial un-measured superposition state — $`\mathrm{\Psi }_0`$ — of the Double Liar (A) is given by any equally pondered superposition of the four true-false states:
$$\frac{1}{2}\left\{\left(\begin{array}{c}0\\ 0\\ 1\\ 0\end{array}\right)\left(\begin{array}{c}0\\ 1\\ 0\\ 0\end{array}\right)+\left(\begin{array}{c}0\\ 1\\ 0\\ 0\end{array}\right)\left(\begin{array}{c}0\\ 0\\ 0\\ 1\end{array}\right)+\left(\begin{array}{c}0\\ 0\\ 0\\ 1\end{array}\right)\left(\begin{array}{c}1\\ 0\\ 0\\ 0\end{array}\right)+\left(\begin{array}{c}1\\ 0\\ 0\\ 0\end{array}\right)\left(\begin{array}{c}0\\ 0\\ 1\\ 0\end{array}\right)\right\}$$
Each next term in this superposition state is the consecutive state which is reached in the course of time, when the paradox is reasoned through. The truth-falsehood values attributed to these states, refer to the chosen measurement projectors.
Making a sentence true or false in the act of measurement, will be described by the appropriate projection operators in $`{}_{}{}^{4}`$<sup>4</sup>. In the case we make sentence 1 (resp. sentence 2) true we get:
$$P_{1,true}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 0\end{array}\right)\mathrm{𝟏}_2P_{2,true}=\mathrm{𝟏}_1\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 0\end{array}\right)$$
The projectors for the false-states are constructed by placing the $`1`$ on the final diagonal place:
$$P_{1,false}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 1\end{array}\right)\mathrm{𝟏}_2P_{2,false}=\mathrm{𝟏}_1\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 1\end{array}\right)$$
As a consequence of making a freely chosen sentence of (A) either true or false, by logical inference the four consecutive states are repeatedly run trough. In order to give a time-ordered description of this cyclic change of state, a continous time $`t`$ is introduced as an ordering parameter. The time-odering parameter extrapolates the discrete moments of consecutive outcomes of the logical inferences, and as such relates to the physical time of reasoning. Under these interpretative restrictions a Schrödinger evolution over ‘time’ can be constructed.
Essentially, a Hamiltonian $`H`$ can be constructed, such that the unitary evolution operator $`U(t)`$ — with $`U(t)=e^{iHt}`$ — describes the cyclic change of logical inferences.
The construction of the evolution operator $`U(t)`$ is more easily accomplished by switching temporarily to an equivalent representation in a larger Hilbert space. Switching to a <sup>16</sup> Hilbert-space is done without any modification or alteration of the problem, as it is isomorphic to the original $`{}_{}{}^{4}`$<sup>4</sup> coupled Hilbert space.
A new basis in <sup>16</sup> is constructed from the basis of $`{}_{}{}^{4}`$<sup>4</sup> ( $`i`$ and $`j`$ from 1 to 4 ) :
$$e_ie_j=e_{\kappa (i,j)}\mathrm{and}\kappa (i,j)=4(i1)+j$$
Where $`\kappa `$ is the natural basis transformation function from the $`C^4C^4`$ to the <sup>16</sup> Hilbert space, and the index from the new basis states in <sup>16</sup>.
The initial superposition state $`\mathrm{\Psi }_0`$ can now be represented in <sup>16</sup> by:
$$\mathrm{\Psi }_0=\frac{1}{2}\{e_{10}+e_8+e_{13}+e_3\}$$
For notational ease we continue to work further in a 4-dimensional subspace of <sup>16</sup>, namely this subspace generated by the basis $`(e_{10},e_8,e_{13},e_3)`$. Obviously we do not loose any information by this restriction. The 4 by 4 submatrix — $`U_D`$ — of the discrete unitary evolution operator, which describes the time-evolution at discrete instants of time when consecutive outcomes of logical inferences have been reached, is:
$$U_D=\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 1& 0\\ 1& 0& 0& 0\\ 0& 1& 0& 0\end{array}\right)$$
In order to obtain a description at every instance of the time-ordering parameter, a procedure of diagonalisation on the submatrix $`U_D`$ is performed, i.e. $`U_D|_{\mathrm{diag}}`$. The diagonalisation procedure allows to solve the matrix equation by breaking it into four uncoupled scalar equations. From the Schrödinger evolution and Stone’s Theorem we obtain:
$$H_{sub}|_{\mathrm{diag}}=i\mathrm{ln}U_D|_{\mathrm{diag}}$$
Inverting the procedure of diagonalisation, the infinitesimal generator of the time-evolution — the submatrix hamiltonian — is obtained :
$$H_{sub}=\left(\begin{array}{cccc}1/2& 1/2& (1i)/2& (1+i)/2\\ 1/2& 1/2& (1+i)/2& (1i)/2\\ (1+i)/2& (1i)/2& 1/2& 1/2\\ (1i)/2& (1+i)/2& 1/2& 1/2\end{array}\right)$$
Or in terms of gamma matrices (these are merely introduced for shorthand notation for there is no implication of relativistic nature):
$$H_{sub}=\frac{1}{2}\left(\gamma _0\gamma _5+\gamma _0\gamma _1\right)+\frac{i}{2}\left(\gamma _1+\gamma _2\gamma _3+\gamma _0\gamma _5\right)$$
The submatrix of the evolution operator $`U(t)`$, valid at all intermediary times $`t`$ is then given by the expression:
$$U_{sub}(t)=e^{iH_{sub}t}$$
The time evolution operator $`U_{sub}(t)`$ in the 4-dimensional subspace of <sup>16</sup> is:
$`U_{sub}(t)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\{(1+e^{it}+e^{it}+e^{2it})1+i(1e^{it}e^{it}+e^{2it})\gamma _2\gamma _3`$
$`(1e^{2it})\gamma _5+i(e^{it}e^{it})\gamma _0\gamma _5i(e^{it}+e^{it})\gamma _1\}`$
In order to finalise our initial claim, we should bring back the hamiltonian $`H`$ as well as the time-evolution operator $`U(t)`$ in the original $`{}_{}{}^{4}`$<sup>4</sup> Hilbert space. Although the outcome is straightforward to obtain by using the basis transformation function $`\kappa (i,j)`$, the complexity of the expression necessitates shorthand notation:
$$H=\underset{\kappa ,\lambda =1}{\overset{16}{}}H_{sub}^{}{}_{\kappa (i,j)\lambda (u,v)}{}^{}O_{iu}O_{jv}$$
and
$$U(t)=\underset{\kappa ,\lambda =1}{\overset{16}{}}U_{sub}^{}{}_{\kappa (i,j)\lambda (u,v)}{}^{}(t)O_{iu}O_{jv}$$
with;
$$O_{iu}O_{jv}=\{e_i.e_u^t\}\{e_j.e_v^t\}$$
For example, the term $`\kappa =3`$ , $`\lambda =10`$ of the time evolution operator $`U(t)`$ is given by;
$$\frac{1}{4}(1ie^{it}+ie^{it}ie^{2it})\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right)\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 0\end{array}\right)$$
Finally, we complete the dynamical picture of the Double Liar cognitive entity (A); when submitted to any measurement at choice, the entity starts its truth-falsehood cycle, when left un-measured the entity remains statically in its undifferentiated superposition state. The latter statement follows immediately from the fact that the initial state $`\mathrm{\Psi }_0`$ is left unchanged by the dynamical evolution $`U(t)`$;
$$\mathrm{\Psi }_0(t)=\mathrm{\Psi }_0$$
$`\mathrm{\Psi }_0`$ is a time invariant, as it is an eigenstate of the Hamiltonian $`H`$. Exactly this time invariance points to the fact that the state $`\mathrm{\Psi }_0`$ describes the existence of the cognitive entity (A) in its cognitive space, independent of any observer. The highly contextual nature of the Double Liar (A) — its unavoidable dynamics engendered by measuring it — implies, intrinsically it can not expose its complete nature, analogous to the quantum entities of the micro-physical world.
The evolution over the time-ordering parameter $`t`$ of the truth behaviour of the cognitive entity during the measurement process can be illustrated graphically.
Temporal evolution of the cognitive evaluation of the Double Liar (A). At time-ordering parameter $`t=0`$, a ‘false’-measurement on (A1) has been executed on the initial state $`\mathrm{\Psi }_0`$ (short-dashed line). The second consecutive step in the logical inference, the ‘true’-state of (A2) at time-ordering parameter $`t=\frac{Pi}{2}`$, and the intermediary times, is shown by the long-dashed line. For clarity, consecutive steps have been deleted. Discrete moments of outcomes of logical inferences are at time-ordering parameter $`t=n\frac{\pi }{2}`$
## 4 Conclusion.
We found that the quantum mechanical formalism can be applied to self-referent cognitive entities as the single liar and double liar paradox. Essentially the two necessary features of the dynamics were constructed; when measured the entity starts its truth-falsehood cycle, when left un-measured the entity remains invariantly in its initial state. A measurement of the entity — engendered by appropriating one sentence its truth or falsehood — sets into action the dynamical evolution which attributes, alternatively over time, truth and falsehood to the coupled sentences. The unmeasured entity $`\mathrm{\Psi }_0`$ does not change over time, its invariance reflects its emergent nature.
The truth-falsehood alternation, due to a measurement process, and the time-invariance of the initial state are both derived from constructed evolution operator $`U(t)`$, driven by the Hamiltonian $`H`$. The quantum formalism therefore has proven an appropriate tool to describe the liar paradox entity.
We set out with the idea of realist cognitive entity, if it can be generally and cognitively influenced on as a practically stable configuration, and by the limited number of different states that it can be in. We put forward that exactly this time invariance points to the fact that the state $`\mathrm{\Psi }_0`$ describes the existence of the cognitive entity in its cognitive space, independent of any observer. For this reason we take the position at liberty to put forward a realistic picture of the cognitive reality, in the sense that $`\mathrm{\Psi }_0(t)`$ represents the state of the real entity.
The nature of the cognitive layer, as here proposed, is essentially different from the space of common macroscopic material objects. The cognitive layer is ‘Hilbert-space like’, and originates from the cognitive person. Quite different from the complexity-theory approach, which emphasizes dynamics, the present formalism clearly describes the emergent entity as an ontological state. An important question in our research therefore remains; the relation between the quantum mechanical state description and the non-linear approach of complex dynamics. Rather speculatively some relation between privileged states in both models could be expected. A model exposing common grounds to both formalisms could in first approach allow inquiry into the relation between the quantum mechanical eigenstates of the Hamiltonian and the phase-space attractors of the corresponding non-linear evolution equations. Eventually, more detailed elaborations of the cognitive layer as an emergent, Hilbert-space like sphere in social groups, can relate both approaches.
The generalisation to other cognitive entities with reference to non-cognitive entities or processes can most probably, to extent, be covered by the $`ϵ`$-model formalism (Aerts, 1986, Aerts and Durt 1994, Aerts, Durt and Van Bogaert 1993). This model allows to describe effects of intermediary contextuality, which is expected in common cognitive entities. Further development needs to bring into the formalism, a general dynamic, and contextual influences other than those of the cognitive person on the truth behavior of the cognitive entities.
Philosophical questions, quite speculative at this stage of our research, can be put forward: e.g. Can we, from the example of the Liar paradox, learn something in general about the nature and origin of dynamical change? What real life experiments can underpin the realistic approach to the cognitive entities?
## 5 References.
Aerts, D., (1986), “A possible explanation for the probabilities of quantum mechanics”, J. Math. Phys., 27, 202.
Aerts, D. (1992), “ Construction of Reality and its Influence on the Understanding of Quantum Structures”, Int. J. Theor. Phys., Vol 31, 10, p. 1813.
Aerts, D. (1994), “The biomousa: a new view of discovery and creation”, in Perspectives on the World, Aerts, D., Apostel, L. et all., VUBPress.
Aerts, D. (1998), “The entity and modern physics: the creation-discovery- view of reality”, in Collection On Objects, ed. Castelani, E., Princeton University Press, Princeton.
Aerts, D. (1999), “The creation discovery view and the layered structure of reality”, in Worldviews and the Problem of Synthesis, the Yellow Book of Einstein meets Magritte, ed. Aerts, D., Van Belle, H. and Van der Veken, J., Kluwer Academic, Dordrecht.
Aerts, D. and Aerts, S. (1995-1996), “Applications of Quantum Statistics in Psychological Studies of Decision Processes”, Foundations of Science 1, 85.
Aerts, D. and Aerts, S. (1987), “Application of Quantum Statistics in Psychological Studies of Decision Processes”, in Foundations of Statistics, eds. Van Fraassen B., Kluwer Academic, Dordrecht.
Aerts, D. Broekaert, J. Ganora L. (1999), “The Quantum Nature of Common Processes”, (in preparation)
Aerts, D. Coecke, B. and Smets, S. (1999), ‘On the Origin of Probabilities in Quantum Mechanics: Creative and Contextual Aspects”, in Metadebates on Science, the Blue Book of Einstein meets Magritte, ed. Cornelis, G., Smets, S. and Van Bendegem, J.P., Kluwer Academic, Dordrecht New York.
D. Aerts, D. and Durt, T., (1994), “Quantum. Classical and Intermediate, an illustrative example”, Found. Phys. 24, 1353.
Aerts, D., Durt, T. and Van Bogaert, B. (1993), “Quantum Probability, the Classical Limit and Non-Locality”, in the proceedings of the International Symposium on the Foundations of Modern Physics 1992, Helsinki, Finland, ed. T. Hyvonen, World Scientific, Singapore, 35.
Broekaert, J. (1999), “World Views. Elements of the Apostelian and General Approach”, Foundations of Science, Vol. 3 ,1 , pp.
CLEA (1996), Research Project: “Integrating Worldviews: Research on the Interdisciplinary Construction of a Model of Reality with Ethical and Practical Relevance” Ministry of the Flemish Community, dept. Science, Innovation and Media.
Grim, P. (1991), The Incomplete Universe, Totality, Knowledge, and Truth, MIT Press, Massachusetts.
Worldviews Group: Aerts, D., Apostel, L., De Moor, B., Maex, E., Hellemans, S., Van Belle, H., Van der Veken, J. (1994);Worldviews: From Fragmentation to Integration, VUB Press, Brussels
Worldviews Group: Aerts, D., Apostel, L., De Moor, B., Maex, E., Hellemans, S., Van Belle, H., Van der Veken, J. (1995);Perspectives on the World: An Interdisciplinary Reflection, VUB Press, Brussels
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# From Crêpes to Pancakes in the MV Model1footnote 11footnote 1Talk presented at the 7th Conference on the Intersections of Particle and Nuclear Physics, Québec City, Québec, Canada, May 22–28, 2000.
## Abstract
The McLerran-Venugopalan model provides a framework which allows one to compute the gluon distribution function of a very large nucleus from the equations of QCD, provided that the longitudinal momentum fraction, $`x__F`$, is sufficiently small. The source of color charge for this computation may be thought of as a crêpe moving along the $`z`$ axis at the speed of light. We refine the MV model by allowing for the presence of non-trivial longitudinal correlations between the color charges that comprise the nucleons. We find that a consistent treatment forces us to consider a pancake-like source which moves at slightly less than the speed of light. Our calculation allows us to consider larger values of $`x__F`$ than were allowed in the original MV model.
Several years ago, McLerran and Venugopalan realized that for large enough nuclei at small enough values of the longitudinal momentum fraction $`x__F`$, it ought to be possible to compute the gluon distribution function using QCDMVmodel1 . Based on this observation, the framework known as the McLerran-Venugopalan model (MV model) was subsequently developedMVmodel2 ; MVmodel3 ; MVmodel4 ; MVmodel9 . Recently, it was shownColorNeut that the infrared divergences present in the MV model may be cured by capturing the physics of confinement via a color-neutrality condition to be imposed on the charge density correlation function used as input to the MV model. In this talk, I will describe workPancakes on extending the MV model to larger values of $`x__F`$.
The MV model as originally formulated in Refs. MVmodel1 ; MVmodel2 ; MVmodel3 ; MVmodel4 ; MVmodel9 , is restricted to very small longitudinal momentum fractions $`x__FA^{1/3}/(ma)`$, where $`A`$ is the number of nucleons, $`m`$ is the nucleon mass, and $`a\mathrm{\Lambda }_{\mathrm{QCD}}^1`$ is the nucleon radius. In this regime, the longitudinal resolution of the gluons is so poor that they probe distances which are much longer than the (Lorentz-contracted) thickness of the nucleus. Thus, all of the quarks inside the nucleus effectively have the “same” value of the longitudinal coördinate $`x^{}`$. At each value of the transverse coördinate x describing this crêpe-shaped (i.e. very thin) nucleus, the color charges from a large number of valence quarks must be summed, resulting in a color charge density which is in a high-dimensional representation of the gauge group. This is a necessary condition for a classical treatment to be valid. In going to larger values of $`x__F`$, we find that the longitudinal resolution of the gluons improves, and we begin to see the longitudinal structure of the nucleus. In order to include this longitudinal structure, we are naturally led to a pancake-shaped geometry (i.e. one with a finite non-zero thickness). The details of this fully 3-dimensional calculation are contained in Ref. Pancakes .
For a classical treatment to be valid, not only must the color charge be in a large representation of the gauge group, but the gauge coupling $`\alpha _s`$ should also be weak. McLerran and VenugopalanMVmodel1 argue that the running coupling ought to be evaluated at the scale $`\mu ^2`$, which is set by the charge-squared per unit transverse area. For large enough nuclei, $`\mu ^2\mathrm{\Lambda }_{\mathrm{QCD}}^2`$, implying that $`\alpha _s(\mu ^2)1`$.
Thus, we begin our computation of the gluon distribution function for a large nucleus by solving the classical Yang-Mills equations describing a pancake-shaped distribution of color charge moving along the $`z`$ axis at nearly the speed of light. The result is a non-linear expression for the vector potential $`A(x^{},\text{x})`$ in terms of the charge density $`\rho (x^{},\text{x})`$. In the spirit of the Weizsäcker-Williams approximation WW , we extract the gluon number density from the two-point correlation function, $`A(x^{},\text{x})A(x^{},\text{x}^{})`$. We replace the quantum mechanical average implied by the angled brackets with a classical average over an ensemble of nuclei. This ensemble is specified by inputting the two-point charge-density correlator $`\rho \rho 𝒟(x^{},\text{x})`$. Furthermore, we assume that the correlations are Gaussian. Confinement is incorporated into the calculation at this stage in the form of a color neutrality condition on $`𝒟`$ColorNeut . When $`𝒟`$ satisfies the color neutrality condition, the two-point correlation function is infrared finite, and may be Fourier-transformed to momentum space, producing a gluon number density $`dN/dx__Fd^2\text{q}`$ which is differential not only in $`x__F`$, but in the transverse momentum q as well.
In the limit $`A^{1/3}1`$, the result of this rather lengthy calculation reads<sup>2</sup><sup>2</sup>2Eq. (1) has been written assuming cylindrical geometry for the nucleus. For a full discussion of the (rather weak) geometric dependence of the result, see Ref. Pancakes .
$$\frac{dN}{dx__Fd^2\text{q}}=3AC_F\frac{2\alpha _s}{\pi ^2}\frac{1}{x__F}d^2𝚫e^{i\text{q}𝚫}(x__F,𝚫)\frac{\mathrm{exp}[N_c𝒳_{\mathrm{}}L(𝚫)]1}{N_c𝒳_{\mathrm{}}L(𝚫)}.$$
(1)
Although complicated in appearance, Eq. (1) is made up of several easily-understood parts. The prefactor shows that at lowest order the number of gluons is simply proportional to the number of quarks. The lowest order result is governed by the function
$$(x__F,𝚫)\frac{1}{2}\frac{d^2\text{q}}{4\pi ^2}e^{i\text{q}𝚫}\frac{\text{q}^2\stackrel{~}{𝒟}(x__F,\text{q})}{[\text{q}^2+(x__Fm)^2]^2},$$
(2)
and is what would be obtained by considering an Abelian theory. The non-Abelian corrections to this result are contained in the exponential factor, which depends on two quantities. First, the the spatial dependence is determined by
$$L(𝚫)\frac{d^2\text{q}}{4\pi ^2}\frac{\stackrel{~}{𝒟}(0,\text{q})}{\text{q}^4}\left[e^{i\text{q}𝚫}1\right].$$
(3)
The strength of the non-Abelian corrections is set by the factor
$$𝒳_{\mathrm{}}=\frac{1}{\pi R^2}\frac{3Ag^4C_F}{N_c^21}8\pi \alpha _s^2A^{1/3}\mathrm{\Lambda }_{\mathrm{QCD}}^2.$$
(4)
These effects are most prominent in very large nuclei. For uranium we have $`𝒳_{\mathrm{}}5`$ or $`6\mathrm{\Lambda }_{\mathrm{QCD}}^2`$. To obtain $`𝒳_{\mathrm{}}=20\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ (as is employed in the plots below) requires of order $`10^4`$ nucleons. Finally, we note that if we set $`x__F=0`$ in Eqs. (1)–(3), we explicitly reproduce the original MV result MVmodel9 .
The functions appearing in Eqs. (2) and (3) are not finite at $`\text{q}=0`$ unless the charge density correlator $`𝒟`$ satisfies the requirement of color neutrality. As explained in Ref. ColorNeut , a key consequence of confinement is the appearance of color neutral nucleons. Mathematically, this consequence may be implemented as a constraint on $`𝒟`$:
$$𝑑\mathrm{\Delta }^{}d^2𝚫𝒟(\mathrm{\Delta }^{},𝚫)=0,\text{or}\stackrel{~}{𝒟}(0,\mathrm{𝟎})=0.$$
(5)
Eq. (5) ensures that widely-separated nucleons are essentially uncorrelated. In order for (5) to be true, $`𝒟`$ must contain a length scale. Not surprisingly, this scale turns out to be the nucleon radius $`a\mathrm{\Lambda }_{\mathrm{QCD}}^1`$. When combined with the assumption of rotational symmetry in the transverse plane, Eq. (5) is sufficient to render the functions contributing to the gluon number density completely infrared finite: the integrals get cut off at $`qa^1\mathrm{\Lambda }_{\mathrm{QCD}}`$.
Aside from the neutrality condition (5), the correlation function $`𝒟`$ is unspecified. In order to illustrate the general features of the gluon number density (1), it is convenient to choose the following form for $`\stackrel{~}{𝒟}`$:
$$\stackrel{~}{𝒟}(x__F,\text{q})=1\frac{1}{1+(a\text{q})^2+(x__Fma)^2}.$$
(6)
In the context of Kovchegov’s nuclear modelModel , the correlation function given in Eq. (6) corresponds to quarks which are distributed within the nucleons according to the weight
$$|\psi (\stackrel{}{r})|^2=\frac{1}{2\pi ^2a^3}\frac{a}{r}K_1\left(\frac{r}{a}\right).$$
(7)
At large distances, the modified Bessel function produces a Yukawa-like behavior in this probability distribution.
Fig. 2 exhibits the behavior of the fully differential gluon number density as a function of the transverse momentum $`\text{q}^2`$. We see that there is saturation: the number of gluons increases as $`\text{q}^20`$ to some maximum value and then stops growing. Qualitatively, these distributions are very much like those obtained by Muller from the point of view of onium scatteringMuller1 ; Muller2 .
The gluon distribution function resolved at the scale $`Q^2`$ is related to the fully differential gluon number density by
$$g_A(x__F,Q^2)_{|\text{q}|Q}d^2\text{q}\frac{dN}{dx__Fd^2\text{q}}.$$
(8)
In Fig. 2 we multiply the fully differential distribution by $`\text{q}^2`$: on the semi-log scale used this produces a “true” representation of where the important contributions to Eq. (8) are. We see that the very low $`\text{q}^2`$ region does not play a significant role provided $`Q^2`$ is not too small. From the plot we see that for $`Q^2\mathrm{}`$, the gluon distribution function becomes insensitive to the presence or absence of the non-Abelian corrections. In fact, we can prove under rather general circumstances that
$$d^2\text{q}\{\frac{dN}{dx__Fd^2\text{q}}|_{\genfrac{}{}{0pt}{}{\mathrm{all}}{\mathrm{orders}}}\frac{dN}{dx__Fd^2\text{q}}|_{\genfrac{}{}{0pt}{}{\mathrm{lowest}}{\mathrm{order}}}\}=0,$$
(9)
independent of $`𝒟`$ ColorNeut .
In Fig. 4, we exhibit $`x__Fg_A(x__F,Q^2)`$ as a function of $`x__F`$ for several different values of $`Q^2`$. At low values of $`Q^2`$, the addition of the non-Abelian corrections reduces the number of gluons from the Abelian result, while by the time $`Q^2=2500\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ is reached, the effect of the non-Abelian terms is negligible. Although we have plotted out to $`x__F=1`$, the calculation is not reliable beyond $`x__F0.25`$Pancakes . For small $`x__F`$ the gluon distribution function exhibits the pure $`1/x__F`$ dependence characteristic of the original MV model.
Finally, we present Fig. 4, which illustrates the deviation of the gluon distribution function from the naïve expectation that for $`A`$ nucleons we should obtain $`A`$ times the result for a single nucleon. We see that at low values of $`Q^2`$ the distribution function grows less rapidly than the number of nucleons as $`A`$ is increased, whereas at large $`Q^2`$, the simple scaling expectation holds.
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# Collective excitation frequencies and damping rates of a two-dimensional deformed trapped Bose gas above the critical temperature
## I introduction
There has been renewed interest in the Bose-Einstein condensation(BEC) after its experimental demonstration by Anderson et.al in a magnetically trapped rubidium gas. In other experiments the Bose-Einstein condensate of trapped lithium and sodium vapor is also observed. There has been much interest in the theoretical understanding of this system . The low-lying collective excitations above the critical temperature in the hydrodynamic region has been discussed by Griffin et.al using the kinetic theory. Damping of the hydrodynamic modes in a trapped Bose gas above the Bose-Einstein critical temperature ($`T_c`$) is also discussed by Odelin et.al and Kavoulakis et.al . Frequencies of the low-lying excitation modes at $`T=0`$ have been discussed by S. Stringari .
After the discovery of BEC in a trapped alkali atom, the influences of the dimension of a Bose systems has been a subject of extensive studies . In our present technology one can frezze the motion of the trapped particles in one direction to create a quasi-2D Bose gas. In the frozen direction the particles executes the zero point motion. To achive this quasi-2D system, the frequency in the frozen direction should be much larger than the frequency in the X-Y plane and the mean field interactions among the particle. It has been shown by V. Baganato et.al that for an ideal 2D Bose gas under harmonic trap a macroscopic ocupation of the ground state can exist at temperature $`T<T_c=\sqrt{\frac{N\lambda }{\zeta (2)}}\frac{\mathrm{}\omega }{k_B}`$. $`\lambda =\frac{\omega _y}{\omega _x}`$ is the deformation parameter. Some experiment shows the possibility of creating a quasi-2D trapped Bose gas .
It is well known that the excitation frequencies for monopole and quadrupole modes are $`2\omega _0`$ and $`\sqrt{2}\omega _0`$ respectively in a 2D isotropic trapped Bose gas. Using the approximation, $`\omega _z>>\omega `$, the dispersion relation of the excitation frequencies , does not produce the correct frequencies for monopole and quadrupole modes in a 2D trapped Bose system. There has been no systematic study on the collective excitations of a 2D deformed trapped Bose gas above the critical temperature.
The purpose of this paper is to give analytic results for the dispersion law of low-lying collective modes in 2D deformed trapped Bose gas and their damping rates in both regime, hydrodynamic and collisionless.
Above the critical temperature, one can distinguish two regimes, the hydrodynamic(collisional) one where collisions ensure the local thermal equlibrium and collisionless where the motion is described by the single particle hamiltonian. In the hydrodynamic region, the characteristic mode frequency is small compared to the collision frequecy ($`\omega \tau <<1)`$. In the collisionless region ($`\omega \tau >>1`$ ), the collision are not important.
The paper is organised as follows. We derive in sec\[II\] a closed equation of motion for the velocity fluctuations of a 2D deformed trapped Bose gas just above the critical temperature ($`T>T_c`$) using the kinetic theory. We make use of this equation in sec\[III\] to calculate the excitation frequencies for a few low-lying collective modes and the corresponding density fluctuations. In sec\[IV\], we derive a dispersion relation in a 2D deformed trap at very high temperature using the method of averages that interpolates between the collisionless and hydrodynamic regimes. From this dispersion relation, we calculate the eigenfrequencies and damping rates for monopole and quadrupole mode. We discuss the evolution of the wave packet width of a Bose gas in a time independent as well as time dependent trap. In sec\[V\], we presents a summary of our work.
## II Hydrodynamic Equation of motion for the velocity fluctuations
We shall discuss the collective modes of a 2D deformed trapped Bose gas in the hydrodynamic regime just above the critical temperature $`T>T_c`$ using the kinetic theory. In the low energy excitations, we can use the semiclassical approximation for the dynamics of a Bose gas, using the following Boltzmann equation for the phase space distribution function $`f(\stackrel{}{r},\stackrel{}{p},t)`$
$$\frac{f}{t}+\stackrel{}{v}._\stackrel{}{r}f+\frac{\stackrel{}{F}}{m}._\stackrel{}{v}f=I_{coll}(f)$$
(1)
where $`I_{coll}`$ is the collisional integral and $`\stackrel{}{F}=U_0(\stackrel{}{r})`$. The trap potential is $`U_0(r)=\frac{1}{2}m(\omega _x^2x^2+\omega _y^2y^2)`$.
In the hydrodynamic regime, collisions ensures the local thermodynamic equilibrium. To lowest order, the perturbed distribution function produced by a slowly varying external potential is the equlibrium Bose distribution function
$$f(\stackrel{}{r},\stackrel{}{p},t)=[\mathrm{exp}(\beta (\stackrel{}{r},t)\eta (\stackrel{}{r},t)1]^1$$
(2)
$$\eta (\stackrel{}{r},t)=\frac{[\stackrel{}{p}m\stackrel{}{v(\stackrel{}{r},t)}]^2}{2m}\mu (\stackrel{}{r},t)$$
(3)
$`\mu (\stackrel{}{r},t)`$ is the chemical potential. The inverse temperature is $`\beta (\stackrel{}{r},t)=\frac{1}{k_BT(\stackrel{}{r},t)}`$. We are only interested in a small perturbations around the equilibrium states.
The conservation laws are :
$$\frac{n(\stackrel{}{r},t)}{t}+.[n_0(\stackrel{}{r})\delta \stackrel{}{v}(\stackrel{}{r},t)]=0$$
(4)
$$mn_0(\stackrel{}{r})\frac{\delta \stackrel{}{v}}{t}=[P(\stackrel{}{r},t)+n(\stackrel{}{r},t)U_0(\stackrel{}{r})]$$
(5)
$$\frac{E(\stackrel{}{r},t)}{t}=.[(P(\stackrel{}{r})+E(\stackrel{}{r}))\delta \stackrel{}{v}]n_0\delta \stackrel{}{v}.U_0(\stackrel{}{r})$$
(6)
These conservation laws are obtained from Eq. (1) multiplying by 1, $`\stackrel{}{p}`$, $`\frac{\stackrel{}{p}^2}{2m}`$ and integrating the resulting equation over $`\stackrel{}{p}`$. During collisions, the total number of particles N, momentum $`\stackrel{}{p}`$, and energy $`\frac{\stackrel{}{p}^2}{2m}`$ are conserved, so the collisional term vanishes. $`\delta \stackrel{}{v}`$ is the velocity fluctuation around the equilibrium states.
Using the quantum statistical mechanics, pressure and density can be written as
$$\frac{P}{k_BT}=\frac{g_2(z)}{\mathrm{\Lambda }^2}$$
(7)
$$n=\frac{g_1(z)}{\mathrm{\Lambda }^2}$$
(8)
where $`g_n(z)=_{i=1}^{\mathrm{}}(\frac{z^i}{i^n})`$ are the Bose-Einstein functions. $`z(r,t)=e^{\beta (r,t)\mu (r,t)}`$ is the local thermodynamic fugacity which is always less than one. $`\mathrm{\Lambda }=\sqrt{\frac{2\pi \mathrm{}^2}{mk_bT}}`$ is the thermal de-Broglie wave length.
One can easily get the relation,
$$P(\stackrel{}{r},t)=E(\stackrel{}{r},t)$$
(9)
in 2D . Using Eq. (9), Eq. (6) can be written as
$$\frac{P(\stackrel{}{r},t)}{t}=2.[(P_0(\stackrel{}{r})\delta \stackrel{}{v}(\stackrel{}{r},t)]n_0\delta \stackrel{}{v}(\stackrel{}{r},t).U_0(\stackrel{}{r})$$
(10)
Taking the time derivative of Eq. (5) and using Eqs. (4) and (10), we get
$$m\frac{^2\delta \stackrel{}{v}}{t^2}=2\frac{P_0(\stackrel{}{r})}{n_0(\stackrel{}{r})}[.\delta \stackrel{}{v}][.\delta \stackrel{}{v}]U_0(\stackrel{}{r})[\delta \stackrel{}{v}.U_0(\stackrel{}{r})]$$
(11)
The closed equation of motion for the velocity fluctuations has been derived by Griffin et.al for 3D trapped Bose system. The term $`\frac{P_0(\stackrel{}{r})}{n_0(\stackrel{}{r})}`$ of (11) is associated with the Bose statistics.
Without any external potential $`U_0=0`$, the Eq. (11) becomes
$$m\frac{^2\delta \stackrel{}{v}}{t^2}=2\frac{P_0(\stackrel{}{r})}{n_0(\stackrel{}{r})}[.\delta \stackrel{}{v}]$$
(12)
It has the plane wave solution with the dispersion relation $`\omega ^2=c^2k^2`$. The sound velocity is $`c^2=\frac{2P_0(\stackrel{}{r})}{mn_0(\stackrel{}{r})}`$ or $`c^2=\frac{2k_BT_0g_2(z_0)}{mg_1(z_0)}`$ where $`z_0=e^{\frac{\mu _0(r)}{k_BT_0}}`$ and $`\mu _0(r)=\mu U_0(r)`$. At high temperature ($`z<<1`$ ), the sound velocity becomes $`c^2=\frac{2k_BT_0}{m}`$. This sound velocity exactly matches with known result.
From the continuity Eq. (4) we have,
$$\frac{\delta n(\stackrel{}{r},t)}{t}=(.\delta \stackrel{}{v})n_0(\stackrel{}{r})\delta \stackrel{}{v}(\stackrel{}{r},t).n_0(\stackrel{}{r},t)$$
(13)
The density fluctuation is given by $`\delta n(\stackrel{}{r},t)=\delta n(\stackrel{}{r})e^{i\omega t}`$. In classical limit, the static density profile is $`n_0(\stackrel{}{r})=n_0(\stackrel{}{r}=0)e^{\frac{m(\omega _x^2x^2+\omega _y^2y^2)}{2\theta }}`$ , where $`\theta =k_BT`$.
## III Eigenfrequencies and The corresponding density fluctuations in the hydrodynamic regime
1) The normal mode solution of (11) is $`\delta \stackrel{}{v}(\stackrel{}{r})=(z^l)`$, here $`z=(x+iy)`$ and $`l>0`$. The excitation frequencies and the associated density fluctuations are $`\omega ^2=l\omega _x^2`$ , $`\delta n_x\omega _x^2xz^{(l1)}n_0(\stackrel{}{r})`$ and $`\omega ^2=\omega _x^2+(l1)\omega _y^2`$ , $`\delta n_y\omega _y^2yz^{(l1)}n_0(\stackrel{}{r})`$.
For isotropic trap, the frequency is $`\omega =\omega _0\sqrt{l}`$. The corresponding density fluctuation is $`\delta n(\stackrel{}{r})n_0(\stackrel{}{r})r^le^{il\theta }`$. At $`r=0`$ there is no density fluctuation. There is a maximum density fluctuation at $`r=\sqrt{\frac{l\theta }{m\omega _0^2}}`$.
2) The other solution of Eq. (11) is $`\delta v(\stackrel{}{r})=[\alpha x^2\pm \beta y^2]`$. The positive sign is for the monopole mode and the negative sign is for quadrupole mode. In a deformed trap, the excitation frequencies are
$$\omega ^2=\frac{1}{2}[3(\omega _x^2+\omega _y^2)\pm \sqrt{9(\omega _x^2+\omega _y^2)^232\omega _x^2\omega _y^2}]$$
(14)
For an isotropic tarp, it becomes $`\omega =2\omega _0`$ or $`\omega =\sqrt{2}\omega _0`$. Hence in the anisotropic trap the monopole mode is coupled to the quadrupole mode. If $`\omega _x<<\omega _y`$, the lowest excitation frequency is $`\omega =\sqrt{\frac{8}{3}}\omega _x`$. If $`\omega _x>>\omega _y`$, the lowest excitation frequency is $`\omega =\sqrt{\frac{8}{3}}\omega _y`$. The density fluctuation for the monopole mode is $`\delta n(\stackrel{}{r})[2\frac{m(\omega _x^2x^2+\omega _y^2y^2)}{\theta }]n_0(\stackrel{}{r})`$ where as the density fluctuation for quadrupole mode is $`\delta n(\stackrel{}{r})(\omega _y^2y^2\omega _x^2x^2)n_0(\stackrel{}{r})`$.
3) There is another quadrupole mode which has velocity field $`\delta \stackrel{}{v}(\stackrel{}{r})=(xy)`$. This is also called scissors mode . The excitation frequency is $`\omega ^2=\omega _x^2+\omega _y^2`$ and the corresponding density fluctuation is $`\delta n(\stackrel{}{r})(\omega _x^2+\omega _y^2)xyn_0(\stackrel{}{r})`$. In isotropic trap $`\omega ^2=2\omega _0^2`$, which agrees with that for the scissors mode in hydrodynamic regime above $`T_c`$ .
## IV Method of Averages
At very high temperature, the dynamical behaviour of a dilute gas is described by the Boltzmann tarnsport equation. Here we include the collisional term in the Boltzmann transport equation and study the eigenfrequencies for monopole and quadrupole mode using the method of averages . These two modes are coupled in a deformed trap.
From Eq. (1), one can get the equations for the average of a dynamical quantity $`\chi (\stackrel{}{r},\stackrel{}{v})`$ is ,
$$\frac{d<\chi >}{dt}<\stackrel{}{v}._\stackrel{}{r}\chi ><\frac{\stackrel{}{F}}{m}._\stackrel{}{v}\chi >=<I_{coll}\chi >$$
(15)
where the average is taken in phase space and $`<\chi >`$ can be written as
$$<\chi >=\frac{1}{N}d^2rd^2vf(\stackrel{}{r},\stackrel{}{v},t)\chi (\stackrel{}{r},\stackrel{}{v})$$
(16)
$`<\chi I_{coll}>`$ can be defined as
$$<\chi I_{coll}>=\frac{1}{4N}d^2rd^2v[\chi _1+\chi _2\chi _1^{}\chi _2^{}]I_{coll}(f)$$
(17)
If $`\chi =a(\stackrel{}{r})+\stackrel{}{b}(\stackrel{}{r}).\stackrel{}{v}+c(\stackrel{}{r})\stackrel{}{v}^2`$ , for elastic collision the collisional term is zero , . a, $`\stackrel{}{b}`$, and c are arbitrary functions of the position.
Now we define the following quantities,
$$\chi _1=x^2+y^2$$
(18)
$$\chi _2=y^2x^2$$
(19)
$$\chi _3=xv_x+yv_y$$
(20)
$$\chi _4=yv_yxv_x$$
(21)
$$\chi _5=v_x^2+v_y^2$$
(22)
$$\chi _6=v_y^2v_x^2$$
(23)
Using the Boltzmann kinetic equation (15), we get the following closed set of equations.
$$<\ddot{\chi }_1>=2<\chi _5>t<\chi _1>+ϵ<\chi _2>$$
(24)
$$<\ddot{\chi }_2>=2<\chi _6>t<\chi _2>+ϵ<\chi _1>$$
(25)
$$<\ddot{\chi }_3>=2ϵ<\chi _4>2t<\chi _3>$$
(26)
$$<\ddot{\chi }_4>=2ϵ<\chi _3>2t<\chi _4>\frac{<\chi _6>}{\tau }$$
(27)
$`<\ddot{\chi }_5>`$ $`=`$ $`ϵ<\chi _6>t<\chi _5>ϵt<\chi _2>`$ (28)
$`+`$ $`{\displaystyle \frac{ϵ^2+t^2}{2}}<\chi _1>`$ (29)
$`<\ddot{\chi }_6>`$ $`=`$ $`ϵt<\chi _1>+{\displaystyle \frac{ϵ^2+t^2}{2}}<\chi _2>`$ (30)
$`+`$ $`ϵ<\chi _5>{\displaystyle \frac{<\dot{\chi }_6>}{\tau }}t<\chi _6>`$ (31)
where double dot indicates the double derivative with respect to time. $`t=\omega _x^2+\omega _y^2`$ and $`ϵ=\omega _x^2\omega _y^2`$. The $`\chi _6`$ is not a conserved quantity, so the collisional contribution comes only through the $`\chi _6`$ term. We have used the fact that $`<\chi _6I_{coll}>=\frac{\chi _6}{\tau }`$, where $`\tau `$ is the relaxation time. This relaxation time $`\tau `$ can be computed by a gaussian ansatz for the distribution function. The relaxation time $`\tau `$ is order of the inverse of the collision rate $`\gamma _{coll}n(0)v_{th}\sigma _0`$, where $`v_{th}=\sqrt{\frac{\pi k_BT}{2m}}`$ is the mean thermal velocity and $`n(0)=\frac{Nm\omega _0^2\lambda }{2\pi k_BTa_z}`$ is the central density for a quasi-2D system. $`a_z`$ is the osscilator length in the z-direction. Hence $`\tau \frac{4a_z}{N\sigma _0\omega _0^2\lambda }\sqrt{\frac{2\pi k_BT}{m}}`$. $`\sigma _0=8\pi a^2`$ is the 3-D scattering cross-section. It can be written in terms of $`T_c`$ as
$$\tau \sqrt{\frac{h}{m\pi ^2\sqrt{\zeta (2)}}}\frac{1}{(N\lambda )^{\frac{3}{4}}}(\frac{a_z}{a^2\omega _0^{3/2}})\sqrt{\frac{T}{T_c}}$$
(32)
The relaxation time $`\tau `$ varies as $`\sqrt{T}`$ in a quasi-2D where as in 3D it varies as T . Now we are looking for a solutions of Eqs. (24)-(30) as $`e^{i\omega t}`$. We have the following dispersion relation
$`(\omega ^24\omega _x^2)(\omega ^24\omega _y^2)+{\displaystyle \frac{i}{\omega \tau }}[\omega ^43\omega ^2(\omega _x^2+\omega _y^2)+`$ (33)
$`8\omega _x^2\omega _y^2]`$ $`=`$ $`0`$ (34)
This dispersion relation interpolates between the collisionless and hydrodynamic regimes. In the hydrodynamical regime ( $`\omega \tau 0`$ ), the first term does not contribute. It gives $`\omega ^2=\frac{1}{2}[3(\omega _x^2+\omega _y^2)\pm \sqrt{9(\omega _x^2+\omega _y^2)^232\omega _x^2\omega _y^2}]`$. This eigen frequencies exactly matches with Eq. (14), a result we found using the equation of motion for the velocity fluctuations even in deformed trap also. We have considered a few low energy excitation modes for which $`.\delta \stackrel{}{v}`$ is constant.The first term of the right-hand side of Eq. (11) does not contribute in the excitation spectrum. Thats why the frequencies of these normal modes are same for a Bose gas just above $`T_c`$ and at very high temperature. In pure collisionless regime ( $`\omega \tau \mathrm{}`$ ), it gives $`\omega _C=2\omega _x`$ and $`\omega _C=2\omega _y`$.
We can write phenomenological interpolation formula for the frequency and the damping rate of the modes in the following form -,
$$\omega ^2=\omega _C^2+\frac{\omega _H^2\omega _C^2}{1i\omega \tau }$$
(35)
The imaginary part of the above equation gives for the damping rate
$$\mathrm{\Gamma }=\frac{\tau }{2}\frac{d}{1+(\omega \tau )^2}$$
(36)
where $`d=(\omega _C^2\omega _H^2)`$. In the hydrodynamic limit ($`\omega \tau 0`$), the damping rate is
$$\mathrm{\Gamma }_{HD}=\frac{\tau }{2}d$$
(37)
where as in the collisionless region ($`\omega \tau \mathrm{}`$),
$$\mathrm{\Gamma }_{CL}=\frac{d}{2\omega _C^2\tau }$$
(38)
The damping rate depends on the difference between the square of the frequencies in the collisional and hydrodynamical regime. The damping rates can be calculated for different values of temperature, number of trapped atoms as well as of the trapping parameters and scattering length through the relaxation time $`\tau `$ (32) . For monopole mode in an isotropic trap the difference $`d`$ is zero. So there is no damping in the monopole mode in a 2D isotropic trapped Bose system when the tempareture is very high. It was first shown by Boltzmann and later Odelin et.al in a 3D trapped Bose system at very high temperature.
For isotropic harmonic trap, Eqs. (24) - (30) decouples into two subsystem, one is for monopole and other one for quadrupole mode. The closed set of equations for monopole mode are
$$<\ddot{\chi _1}>=2<\chi _5>2\omega _0^2<\chi _1>$$
(39)
$$<\ddot{\chi _3}>=4\omega _0^2<\chi _3>$$
(40)
$$<\ddot{\chi _5}>=2\omega _0^4<\chi _1>+2\omega _0^2<\chi _3>$$
(41)
There is no collisional term in the above equations. So there is no damping for the monopole mode of a classical dilute gas confined in a isotropic trap. We are looking for solutions of Eqs. (39) - (41) as $`e^{i\omega t}`$, we get $`\omega =2\omega _0`$.
The Eqs. (39) - (41) can be re-written as
$$\ddot{<\chi _1>}\frac{1}{2<\chi _1>}(\dot{<\chi _1>}^2)+2\omega _0^2\chi _1=\frac{Q}{\chi _1}$$
(42)
where $`Q=2(<\chi _1><\chi _5><\chi _3>^2)`$ is invariant quantity under time evolution. Define $`X(t)=\sqrt{<\chi _1>}`$ which is the wave packet width and substituting it into Eq. (42) gives,
$$\ddot{X}+\omega _0^2X=\frac{Q}{X^3}$$
(43)
This is a nonlinear singular Hill equation. The same equation is obtained at $`T=0`$ in 2D by Garcia Ripoll et.al . At equlibrium, $`X_0^4=\frac{Q}{\omega _0^2}`$. We linearized the Eq. (43) around the equilibrium point $`X_0`$, we get
$$\ddot{\delta X}+4\omega _0^2\delta X=0$$
(44)
One obtains an oscillation frequancy of the gas is $`\omega =2\omega _0`$, corresponding to the frequency of a single particle excitation in the gas.
For time dependent trap, the equation of motion for the width of the wave packet is
$$\ddot{X}+\omega _0^2(t)X=\frac{Q}{X^3}$$
(45)
The general solution is $`X(t)=\sqrt{u^2(t)+\frac{Q}{W^2}v^2(t)}`$ where $`u(t)`$ and $`v(t)`$ are two linearly independent solutions of the equation $`\ddot{p}+\omega _0^2(t)p=0`$ which satisfy $`u(t_0)=X(t_0)`$, $`\dot{u}(t_0)=X^{}(t_0)`$, $`v(t_0)=0`$, $`v^{}(t_0)0`$. W is the Wronskian.This time dependent Hill equation (45) can be solved explicitly only for a paricular choices of $`\omega _0(t)`$.
The closed set of equation for quadrupole mode in isotropic trap are
$$\ddot{<\chi _2>}=2<\chi _6>2\omega _0^2<\chi _2>$$
(46)
$$\ddot{<\chi _4>}=4\omega _0^2<\chi _4>\frac{<\chi _6>}{\tau }$$
(47)
$$\ddot{<\chi _6>}=2\omega _0^4<\chi _2>\frac{<\dot{\chi }_6>}{\tau }2\omega _0^2<\chi _6>$$
(48)
Solving this set of equation, we get damped quadrupole mode,
$$(\omega ^24\omega _0^2)+\frac{i}{\omega \tau }(\omega ^22\omega _0^2)=0$$
(49)
In the hydrodynamic regime, the oscillation frequency is $`\omega ^2=2\omega _0^2`$ where as in the collisionless region, the frequency is just a single particle oscillator frequency. The damping rate can be calculated by using the Eqs. (37) and (38).
## V SUMMARY
In this work, we derived the equations of motion for velocity fluctuations of a Bose gas in a 2D deformed trap potential just above the critical temperature. When $`U_0=0`$, it becomes a wave equation, from which we found the exact sound velocity at high temperature. We have also computed the frequency of the scissors mode in hydrodynamic regime above $`T_c`$ which agrees with the result obtained by Odelin et.al . We have also calculated the frequencies for monopole and quadrupole mode and the corresponding density fluctuations in a deformed trap above $`T_c`$.
Using the method of averages, we obtained a dispersion relation that interpolates between the collisionless and hydrodynamic regimes at very high temperature. In a deformed trap as well as an isotropic trap, we have found the frequencies and the damping rates (in terms of the relaxation time) for monopole and quadrupole mode in both the regimes. In hydrodynamical regime, the excitation frequencies for monopole and quadrupole mode exactly matches with the previous result that we have found from equatin (11). We have also shown that the relaxation time $`\tau `$ varies as $`\sqrt{T}`$ in quasi-2D Bose gas.
We have shown that there is no damping for monopole mode in a 2D isotropic trapped Bose gas when the temperature is very high. It was shown by Boltzmann , later Odelin et.al for 3D isotropic trap.
We have discussed about the time evolution of the wave packet width of a Bose gas in a time independent as well as time dependent isotropic trap. It can be described by the solution of the Hill equation.
## VI ACKNOWLEDGEMENT
I would like to thank G. Baskaran and Subhasis Sinha for helpful discussions.
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# Optical phonons in a quarter-filled 1D Hubbard model
## I Introduction
Electron-phonons interactions have always attracted a lot of attention. Indeed, they are responsible for a large number of phase transitions, such as BCS supra-conductivity or Peierls transitions. Among the materials for which electron-phonons coupling are important, molecular materials constitute a special class. A trivial statement is that molecular crystals differ from simple crystals by the fact that their basic units are molecules. These basic units are therefore structured systems with a large number of internal degrees of freedom that may interact with the valence or conduction electrons. In particular, molecular crystals have two kinds of electron-phonons interactions. The first one is inter-molecular and the phonons essentially couple to the electronic structure through a modulation of the hopping parameter between two molecular sites. This electron-phonons coupling is responsible for the Peierls transitions and has been extensively studied. The second type of interaction is intra-molecular. The totally-symmetric molecular vibrational modes couple to the electronic structure, essentially through a modification of the on-site parameters such as molecular orbital energies and on-site repulsions. In the $`60`$’s, Little suggested that intra-molecular vibrations could be responsible for supra-conductivity in organic conductors, more recently they were proposed as mediators for supra-conductivity properties in materials such as fullerides .
Indeed in systems such as the organic conductors, there is a broad spectrum of intra-molecular vibrations ranging from $`200cm^1`$ to $`2000cm^1`$ , that is between $`t/5`$ and $`t`$, where $`t`$ is the hopping inter-molecular integral responsible for the conduction. Moreover it was shown that these vibrations couple quite well with the electronic degrees of freedom both in the hight and low vibrational range of the Raman spectrum . This fact can easily be understood from simple chemical considerations. Indeed, the molecules acting as basic units in these systems share number of characteristics such are being large, strongly conjugated, and built from pentagonal cycles. These characteristics allow them to adjust their geometry to their electronic charge at a low energetic cost, essentially by a modification of the angles in the pentagonal cycles. This mechanism leads to an important electron intra-molecular vibration coupling.
Previous studies on correlated electronic systems coupled to intra-molecular vibrations have essentially explored the two asymptotic regimes, (i) the weak coupling regime that can be treated by perturbative expansion from the electronic model and (ii) the strong coupling or polaronic regime . It is clear that the study of the intermediate coupling regime cannot be done by analytic treatments and requires up to date numerical techniques. The purpose of this work is to fill this gap and systematically explore the phase diagram of a one dimensional, correlated, electronic system coupled to intra-molecular vibrations, both as a function of the coupling constant and the correlation strength. The present paper systematically studies the Hubbard-Holstein model for a one dimensional quarter-filled chain for two values of the vibration frequency ,$`\omega =0.2t`$ and $`\omega =t`$, respectively corresponding to the top and the bottom part of the organic conductors Raman spectra. For each phase, properties such as spin and charge gaps, distance dependence of the spin, charge and singlet correlation functions, etc, are reported.
The next section analyses the Hubbard-Holstein model and develops the computational choices. The third section reports and discusses the results and the last section is devoted to the conclusion.
## II Model analysis and Computational details
### A The Hubbard-Holstein model
The Hubbard-Holstein Hamiltonian associates a Hubbard Hamiltonian, which includes short range electron correlations, with dynamical phonons. The latter are linearly coupled with the electronic degrees of freedom as in the Holstein model .
$`H`$ $`=`$ $`H_e+H_{ph}+H_{eph}`$
$`\mathrm{with}`$
$`H_e`$ $`=`$ $`ϵ{\displaystyle \underset{i,\sigma }{}}n_{i,\sigma }+t{\displaystyle \underset{i,\sigma }{}}(c_{i+1,\sigma }^{}c_{i,\sigma }+c_{i,\sigma }^{}c_{i+1,\sigma })+U{\displaystyle \underset{i}{}}n_{i,}n_{i,}`$
$`H_{ph}`$ $`=`$ $`\omega {\displaystyle \underset{i}{}}(b_i^{}b_i+1/2)`$
$`H_{eph}`$ $`=`$ $`g{\displaystyle \underset{i}{}}n_i(b_i^{}+b_i)`$
where $`c_{i,\sigma }^{}`$, $`c_{i,\sigma }`$ and $`n_{i,\sigma }`$ are the usual creation, annihilation and number operators of electrons of spin $`\sigma `$ on site $`i`$ ($`n_i=n_{i,}+n_{i,}`$). $`b_i^{}`$ and $`b_i`$ are the intra-molecular phonons creation and annihilation operators.
From the point of view of the isolated molecule the Hubbard-Holstein (HH) model tries to mimic the relaxation of the molecular geometry as a function of the ionicity. Indeed the on-site part of the HH model can be rewritten as
$`H_i`$ $`=`$ $`\epsilon n_{i,\sigma }+Un_in_i+\omega \left(b_i^{}b_i+1/2\right)+gn_i\left(b_i^{}+b_i\right)`$ (1)
$`=`$ $`\omega [((b_i^{}+n_i{\displaystyle \frac{g}{\omega }})(b_i+n_i{\displaystyle \frac{g}{\omega }})+{\displaystyle \frac{1}{2}}]+(U2{\displaystyle \frac{g^2}{\omega }})n_in_i+n_i(ϵ{\displaystyle \frac{g^2}{\omega }})`$ (2)
The above formulation points out the three effects treated in the HH model.
1. The modification of the molecular orbital energy: $`\epsilon \epsilon g^2/\omega `$. This effect is very important (i) in multi-band systems since it strongly affects the relative filling in the different bands and (ii) in opened systems where electrons can jump in and out from an external bath. In our case it just changes the energy reference.
2. The decrease of the effective on-site bi-electronic repulsion: $`UU2g^2/\omega `$. In the strong coupling regime the effective electron-electron interaction can become attractive due to the electron-phonons interaction, and two electrons held together via molecular vibrations.
3. The displacement of the harmonic oscillator describing the intra-molecular vibrations as a function of the molecular charge $`n_i`$: $`\omega \left[b_i^{}b_i+1/2\right]\omega \left[\left(b_i^{}+n_ig/\omega \right)\left(b_i+n_ig/\omega \right)+1/2\right]`$. This term mimics the relaxation of the molecular geometry as a function of the molecule ionicity. The $`\lambda =n_ig^2/\omega `$ term acts as an effective molecular coordinate for which the equilibrium geometry is linearly shifted from $`\lambda =0`$, when the molecular site does not carry any electron, up to $`\lambda =2g^2/\omega `$, when it carries $`2`$ electrons.
One sees from equation 2 that the vibronic molecular states are coherent phonons states, eigenstates of the shifted harmonic oscillators. They can be referred as $`|n_i,Sz_i,\nu _j^{n_i}`$, where $`\nu _j^n`$ is the vibrational quantum number of the molecule $`i`$ when it supports $`n_i`$ electrons,
$`H_i|n_i,Sz_i,\nu _j^{n_i}`$ $`=`$ $`\left[n_i\left(ϵ{\displaystyle \frac{g}{\omega }}\right)+\delta (n_i2)\left(U2{\displaystyle \frac{g^2}{\omega }}\right)+\omega \left(\nu _j^{n_i}+{\displaystyle \frac{1}{2}}\right)\right]|n_i,Sz_i,\nu _j^{n_i}`$ (3)
where $`\delta `$ is the Dirac function.
### B Computational details
The calculations on the infinite chain are performed using the infinite system Density Matrix Renormalization Group method . The main problem risen by the HH model is the infinite number of vibronic states on each sites. In order to render the calculations feasible, the basis set have been truncated to the lowest vibronic states of each molecular site, that is $`|n_i,Sz_i,\nu _j`$ such that $`\nu _j=0,1`$. This choice is physically reasonable since (i) we work at $`T=0`$ and therefore only the lowest vibronic states are expected to be involved, (ii) the molecules form well defined entities that are only perturbatively modified by the presence of their neighbors. The drawback of this choice is that we directly work in the vibronic basis set and the truncation of the basis set destroys any further possibility to separate the electronic degrees of freedom from the vibrational ones. We will see later, from a wave function analysis of the different phases, that the truncation does not affect the results of our calculations as long as we are not too close to a phase transition.
In the phase diagram explorations we kept 100 states per renormalized block, while in the properties calculations 256 states were kept. The charge and spin gaps where computed using a double extrapolation (i) on the system size and (ii) on the number of states kept $`m`$. Typically the extrapolations over $`m`$ were done from three DMRG calculations with respectively $`100`$, $`150`$ and $`256`$ states kept. The maximum number of sites is $`84`$ and the correlation functions presented below are done for this chain length. In some cases, exact diagonalizations on small clusters have also been performed in order to better analyze either the wave function or the energy spectrum.
## III Results
The present work explores in a systematic and unbiased way the whole range of electron-phonons coupling regime (from $`g/t=0`$ to $`g/t=1.25`$) as well as the whole range of correlation strength (from $`U/t=0`$ to $`U/t=16`$), for two values of the phonons frequency: $`\omega =0.2t`$ and $`\omega =t`$. These frequencies have been chosen in order to approximatively correspond to the low and high part of the Raman spectrum of the $`TTF`$ , $`TMTSF`$ or $`M(dmit)_2`$ molecules. The adiabatic regime corresponds to the parameter domains $`\omega /t1`$ and $`g/\omega 1`$, and the strong coupling regime corresponds to $`g/\omega 1`$ and $`\omega /t1`$.
### A $`\omega =0.2t`$
Figure 1 reports the phase diagram for $`\omega =0.2t`$ as a function of $`g/\omega `$ and $`U/t`$ which are, with $`\omega /t`$, the three relevant parameters in the HH model (see eq. 2).
Five different phases have been characterized: two insulating phases where the electrons are strongly localized, a polaronic phase (diamonds) and a bi-polaronic phase (circles) and three metallic phases (crosses, stars and plus). These phases were characterized using the following set of tools
(i) the variation of the energy with the number of sites,
(ii) the spin and charge gaps,
(iii) the spin, charge and singlet-singlet correlation functions,
(iV) density matrices at the central sites.
#### 1 The bi-polaronic phase
In the strong coupling regime, $`U_{eff}=U2\omega \left(g/\omega \right)^2`$ becomes negative and a bi-polaronic phase was found. As expected, the attractive character of $`U_{eff}`$ strongly couples the electrons in pairs. For instance, for $`U/t=4`$, $`g/t=1`$ and $`84`$ sites ($`U_{eff}=6t`$, $`g/\omega =5`$), the probability of having a lonely electron on the central site is smaller than $`10^{12}`$. The electron pairing induced by the intra-molecular vibrations is very strong, however, due to the Franck-Condon factors, this phase does not correspond to the singlet super-conducting phase but rather to a localized bi-polaronic phase. Let us suppose that an electron on a site $`i`$ would like to hop on a neighboring $`j`$, omitting the spin degree of freedom the hopping term can be written as
$`a_j^{}a_i|n_i,\nu _\alpha ^{n_i};n_j,\nu _\beta ^{n_j}`$ $`=`$ $`t|n_i1,\nu _\alpha ^{n_i};n_j+1,\nu _\beta ^{n_j}`$ (4)
$`=`$ $`t{\displaystyle \underset{\gamma ,\delta }{}}\nu _\alpha ^{n_i}|\nu _\gamma ^{n_i1}\nu _\beta ^{n_j}|\nu _\delta ^{n_j+1}|n_i1,\nu _\gamma ^{n_i1};n_j+1,\nu _\delta ^{n_{i+1}}`$ (5)
that is the hopping integral between the molecular vibronic ground states is rescaled by the product of the Franck-Condom factors on the two sites $`i`$ and $`j`$, $`\nu _0^{n_i}|\nu _0^{n_i1}\nu _0^{n_j}|\nu _0^{n_j+1}`$. The relaxation energy or self-trapping energy (due to the vibrations) of the electron pair on a site can be evaluated as the difference between the vertical ionization potential (or electron affinity) and the adiabatic one,
$`E_{relax}(i)`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}\nu _\alpha ^{n_i\pm 1}\omega \nu _0^{n_i}|\nu _\alpha ^{n_i\pm 1}^2`$
For $`U/t=4`$ and $`g/t=1`$ this relaxation energy is as large as $`5.00t`$. Figure 2 shows the overlap between the vibrational ground state corresponding to $`n_i`$ electrons on site $`i`$ ($`\nu _0^{n_i}`$) and the different vibrational states corresponding to $`n_i\pm 1`$ electrons ($`\nu _\alpha ^{n_i\pm 1}`$). One sees immediately that for large or even intermediate values of $`g/\omega `$ the overlap between the low energy vibrational states is very small. The consequence is that the electron hopping between two neighboring sites is strongly hindered ; either the transfer takes place towards small quantum-number vibrational states and the transfer integral is strongly renormalized by the small Franck-condom factors, or the transfer takes place towards large quantum-number vibrational states and it is hindered by the vibrational energetic cost. It is clear that the same phenomenon hinders even more the pair hopping since the displacement is twice as large and the Franck-Condon factor is squared ($`\nu _0^{n_i}|\nu _0^{n_i\pm 2}=\mathrm{exp}(2g^2/\omega ^2)`$).
If one looks at the energy per site as a function of the system size (or DMRG iteration number) , one sees that it is nearly constant. The amplitude of variation between the $`4`$ sites system and the $`84`$ sites system is smaller than $`5\times 10^{12}`$ for $`U/t=4`$ and $`g/\omega =5`$, the typical example we have chosen for this phase. This fact could have been surprising but it is easily understood in the context of a strongly hindered hopping. Taking into account the electron pairing and the strong localization one can reasonably approximate the system wave-function by $`\mathrm{\Psi }_{BP}`$
$$\begin{array}{ccccccccccc}& & & & & & & & & & \\ & & & & & & & & & & \end{array}$$
of energy
$$E(\mathrm{\Psi }_{BP})=N_{sites}\left(\frac{1}{4}U_{eff}\frac{1}{2}\frac{g^2}{\omega }+\frac{1}{2}\omega \right)$$
The quality of this approximation can be checked on the energy. The computed DMRG energy per site differs from $`E(\mathrm{\Psi }_{BP})=3.9`$ by at the most $`7\times 10^{12}`$ in $`t`$ units. Exact diagonalization on a $`4`$-sites system, where the wave function can be explicitly analyzed, confirms this result. The projection of the exact ground state on $`\mathrm{\Psi }_{BP}`$ being larger than $`110^7`$. In fact the ground state of the $`4`$-sites system is $`4`$ times quasi-degenerated and the infinite system ground state infinitely degenerated. Indeed, in the absence of inter-site Coulomb repulsion, any choice for the localization of the electron pairs is equivalent. Even with a $`1/r`$ Coulomb repulsion one sees that the infinite system should remain $`4`$ times quasi-degenerated, the degeneracy lifting being of the order of magnitude of the rescaled hopping integral between the low energy vibrational states, that is of the order of $`t\mathrm{exp}(g^2/\omega ^2)`$ ($`1.4\times 10^{11}`$ for $`U/t=4`$ and $`g/\omega =5`$). The charge gap is therefore exponentially small, the exact diagonalization of the 4-sites system shows that the first exited state with a ”real” gap is a one-boson vibrational excited state at $`\omega `$. The spin gap is of the order of magnitude of $`U_{eff}`$ ($`6.0`$ in our example) since it necessitates the breaking of an electron pair.
#### 2 The Lüttinger liquid phase
In the weak coupling regime — for small values of $`g/\omega `$ — up to the intermediate coupling regime for intermediate values of the correlation strength, one finds a phase which is essentially a Tomanaga-Lüttinger liquid , with parameters slightly rescaled by the presence of the vibrations compare to the purely electronic system. This result was expected, from continuity from the pure Hubbard model, and from previous works from Voit and Schulz (within a renormalization group (Rg) scheme in an incommensurate system). We will further refer to this phase as the Tomanaga-Lüttinger (TL) phase. While restricted to a very small range of $`g/\omega `$, for large values of the correlation strength, it is worth to notice that this phase extends up to values of $`g/\omega `$ larger than $`1`$ for intermediate values of $`U`$ (see figure 1). For very small values of $`U`$, the competition between the TL phase and a bi-polaronic phase limits the extension of the former to a parameter range for which the effective repulsion remains positive, that is $`g/\omega <\sqrt{U/2\omega }`$. The study of the charge-charge and spin-spin correlation functions yield that the main effect of the vibrational degrees of freedom is to rescale the TL liquid parameters — in particular the value of $`K_\rho `$ — compared to the purely electronic system. This result is in total agreement with the conclusions derived by Voit et al. .
A simple approach would bring us to think that the value of $`K_\rho `$ is increased by the vibrational degrees of freedom since the effective electron-electron repulsion is strongly reduced by the electron-phonons interactions. Let us see what really comes out of the calculations. The correlation functions have been computed on a $`84`$ sites system with $`256`$ states kept in the DMRG procedure. They exhibit the expected behavior as a function of the inter-site distance with a power law decay, in agreement with the ungaped nature of both the spin and charge channel. Figure 3 shows the charge structure factor for the pure Hubbard model as well as the Hubbard-Holstein one for the $`U/t=1`$ and $`g/\omega =1.25`$ set of parameters.
We can see on figure 3 that the numerical calculations contradict the simple prediction. The structure factor derivative at $`q=0`$ is directly proportional to $`K_\rho `$
$`K_\rho `$ $`=`$ $`{\displaystyle \frac{d}{dq}}\left(\pi {\displaystyle _0^+\mathrm{}}\left(ne(0)ne(0)\right)\left(ne(r)ne(r)\right)e^{iqr}𝑑r\right)|_{q=0}`$ (6)
and is clearly smaller for the $`U/t=1`$ and $`g/\omega =1.25`$ set of parameters than for the $`U/t=1`$ Hubbard model. The $`K_\rho `$ values computed from the structure factors are $`0.79`$ for $`g/\omega =1.25`$ and $`0.89`$ for the Hubbard Hamiltonian — in total agreement with the values found in the literature from numerical resolution of the Bethe ansatz solution . This unexpected reduction of $`K_\rho `$ was already noticed by Voit et al. from Rg considerations. One should notice that it cannot be attributed to the rescaling of the hopping integral since
$$\frac{U_{eff}}{t_{eff}}=\frac{U2g^2/\omega }{t\mathrm{exp}\left(g^2/\omega ^2\right)}=1.31$$
(7)
and the corresponding value of $`K_\rho `$ is much larger than $`0.85`$ .
#### 3 The $`4k_F`$ CDW phase
In a regime of intermediate coupling and intermediate to large correlation strength, we found a phase for which the charge correlation functions are dominated by $`4k_F`$ charge density wave (CDW) fluctuations. The multiplication of the correlation functions by $`r^2`$ ($`r`$ being the inter-sites distance) allows the elimination, when the Fourier transform is performed, of the $`1/r^2`$ term that contribute at all frequencies. The frequency analysis of $`r^2\times `$the correlation function is therefore much clearer, and the relative importance of the $`2k_F`$ and $`4k_F`$ terms enhanced. Figure 5 reports this Fourier analysis for $`U/t=4`$ and $`g/\omega =1.25`$. One can see that while the spin-spin correlation function is change very little compared to the Hubbard model, the $`2k_F`$ contribution to the charge-charge correlation function has disappeared and a strong $`4k_F`$ contribution has set place. Figure 5 reports the same correlation functions as a function of the inter-sites distance. Both charge and spin correlation functions decrease as a power law as a function of the inter-sites distance speaking in favor of a gap-less system in both the spin and charge channels. The charge and spin gaps have been computed independently from double extrapolations (i) on the chain length and on the number of states kept in the DMRG calculations, i.e. $`100`$, $`150`$ and $`256`$. While the charge channel clearly extrapolates toward a null gap, in the spin channel the question is not as clear. Indeed, the gap is found to be $`\mathrm{\Delta }_\sigma (100)=8.1\times 10^3`$ for a DMRG calculation where $`100`$ states are kept, $`\mathrm{\Delta }_\sigma (150)=8.3\times 10^3`$ for $`150`$ states kept and $`\mathrm{\Delta }_\sigma (256)=8.6\times 10^3`$ for $`256`$. These gap values are very small however they do not extrapolate toward a null value when the quality of the calculation increases. This behavior pleads for a non null but very small gap in the spin channel. One should note that this result is not incoherent with the spin-spin correlation function behavior since such a small gap means a very large coherence length, of the order of magnitude of $`\mathrm{\Delta }_\sigma ^1`$, i.e. the exponential behavior of the correlation function should take place at inter-sites distances larger than the chain length.
In the Lüttinger liquid theory the $`4k_F`$ phase is supposed to occur for values of the $`K_\rho `$ parameter smaller than $`1/3`$, that is (i) in very strongly correlated systems and (ii) for values of $`K_\rho `$ unreachable in the Hubbard model — for which $`0.5K_\rho 1`$. The presence of intra-molecular vibrations not only allows the existence of a $`4k_F`$ CDW phase, but this phase can be reached for values of the bi-electronic repulsion as small as $`U/t=2.5`$ (and most probably even lower). We computed the $`K_\rho `$ values from the structure factors and found that similarly to what happened in the TL phase, the $`K_\rho `$ value is strongly diminished in comparison to the purely electronic case. For $`U/t=4`$ and $`g/t=0.25`$ the value of $`K_\rho `$ is $`0.56`$ while it is $`0.71`$ for the $`U/t=4`$ Hubbard model. One should however note that the $`0.56`$ value is in good agreement with the $`K_\rho `$ expected for the effective parameters, $`U_{eff}/t_{eff}=16.1`$.
#### 4 The small polarons phase
For positive $`U_{eff}`$ but relatively large values of $`g/\omega `$, the system is in a small polarons phase (diamonds on fig. 1). Indeed the electrons are no longer paired but are still strongly localized. The probability of having two electrons on the central site is smaller than $`10^9`$ for $`U=8`$, $`g/t=0.75`$ and $`82`$ sites ($`g/\omega =3.75`$ and $`U_{eff}=2.375`$). These lonely electrons remain strongly coupled to the intra-molecular vibrations. The on-site vibrational relaxation energy is as large as $`2.8t`$ and the rescaled hopping integral between the ground vibrational states as low as $`7.8\times 10^7`$. The ground-state wave-function can be approximated by the totally localized wave-function $`\mathrm{\Psi }_P`$
$$\begin{array}{ccccccccccc}& & & & & & & & & & \\ & & & & & & & & & & \end{array}$$
and the residual hopping term treated perturbatively. The overlap of the exact ground-state wave-function of an $`4`$-sites system with $`\mathrm{\Psi }_P`$ is as large as $`0.99`$ for $`U/t=8`$ and $`g/t=0.75`$ and even as large as $`0.98`$ for $`U=4`$ and $`g/t=0.5`$ which is very closed to the phase transition. Coherently, the DMRG energies per site remains nearly constant as the system size increases, with a maximal variation of $`3\times 10^7`$ from $`E(\mathrm{\Psi }_P)/N_{sites}=1/2g^2/\omega +1/2\omega `$, for $`U=8`$ and $`g/t=0.75`$. As for the bi-polaronic phase the absence of Coulombic inter-site repulsion in our model induces a strong quasi-degeneracy of the ground-state: $`12`$ times for the $`4`$-sites system but infinitely for the infinite system. Even in the presence of a $`1/r`$ repulsion, the system would remain at least twice quasi-degenerated due to the even sites versus odd sites equivalence with an additional factor of 4 due to the spin quasi-degeneracy. Capone et al. have studied the conditions of existence of this phase by exact diagonalization on small clusters (4 and 6 sites). They found that the localized polarons phase is stable as long as the polarization energy defined as $`\epsilon _{pol}=N_{sites}\rho (1\rho )g^2/\omega `$ is larger than the electronic energy $`\epsilon _{elec}`$. $`\rho `$ is the filling of the system. The curve $`\epsilon _{pol}/\epsilon _{elec}=1`$ is reported on fig. 1 with a dashed-dotted line and fits relatively well with the phase boundary in the infinite system. In the same spirit curves of constant $`U_{eff}/t_{eff}`$ have been plotted. This small polarons phase corresponds to very large values of the $`U_{eff}/t_{eff}`$ ratio.
#### 5 The Luther-Emery phase
Finally for negative values of $`U_{eff}`$, but large values of $`t_{eff}`$ a delocalized phase takes place for which all spin, charge and on-site singlet correlation functions seem to decrease with the inter-site distance as a power law. Figure 6 reports the absolute values of the correlation functions for $`U/t=0.05`$ and $`g/t=0.1`$, that is $`g/\omega =0.5`$, $`U_{eff}=0.05`$ and $`t_{eff}=0.78`$. In all computed cases ($`U/t=0.2`$, $`g/\omega =1`$ and $`1.3`$ ; $`U/t=0.05`$, $`g/\omega =0.5`$) the charge density fluctuations dominate, however as $`U_{eff}/t_{eff}`$ decreases the on-site singlet fluctuations increases compared to the charge ones. A direct calculation of the charge and spin gaps yield a clearly ungaped charge channel and a slightly gaped spin channel. The computed spin gaps are $`\mathrm{\Delta }_\sigma =0.05`$ for $`U/t=0.05`$ and $`g/t=0.1`$, and $`\mathrm{\Delta }_\sigma =0.06`$ for $`U/t=0.2`$ and $`g/t=0.2`$. These very small values are compatible with the the power law decrease of the spin correlation functions since with such small gaps, the correlation length is very large — of the order of magnitude of $`\mathrm{\Delta }_\sigma ^1`$ — and therefore the exponential decay of the correlation function cannot take place at distances smaller than the order of magnitude of the chain length. One can therefore identify this phase with a Luther-Emery model . The values of $`K_\rho `$ extracted from the structure factors are again smaller than the values of the purely electronic model, with $`K_\rho =0.99`$ for $`U/t=0.05`$ and $`g/t=0.1`$, and $`K_\rho =0.89`$ for $`U/t=0.2`$ and $`g/t=0.2`$. Of course, they are as well smaller than the $`K_\rho `$ values corresponding to the negative $`U_{eff}`$.
### B $`\omega =t`$
Figure 7 reports the phase diagram for $`\omega =t`$ as a function of $`g/\omega `$ and $`U/t`$.
For this intermediate value of the phonons frequency, one finds only three phases, a slightly perturbed Lüttinger liquid phase for $`U_{eff}>0`$, a Luther-Emery phase for $`U_{eff}<0`$ and small $`g/\omega `$, a localized bi-polarons phase for $`U_{eff}<0`$ and larger $`g/\omega `$.
For $`U_{eff}>0`$, the phase is delocalized, dominated by $`2k_F`$ SDW fluctuations. Both charge and spin correlation functions decrease as power laws as a function of the inter-sites distance and both spin and charge channels are ungaped to numerical accuracy. The structure factors for the $`U/t=4`$ Hubbard model and for the Hubbard-Holstein model with $`U/t=4`$ $`g/\omega =0.5`$ and $`U/t=4`$ $`g/\omega =1`$ are reported in figure 8. As can be seen the HH model is in this case indistinguishable from the pure Hubbard model and, as a consequence, the $`K_\rho `$ values as not sensitive (to numerical accuracy) to the electron-phonons coupling. In fact everything goes as if we were in the adiabatic regime while the model parameters can be as far from it as $`\omega /t=1`$ and $`g/\omega =1.25`$. This is typically a case where the Born-Oppenheimer approximation is fully valid and the electronic and vibrational degrees of freedom are nearly independent.
For $`U_{eff}<0`$ one has a Luther-Emery phase for small values of $`g/\omega `$. This phase is very similar to the one found for $`\omega /t=0.2`$, with an effective attractive on-site interaction and large effective hopping integrals. The CDW and the on-site singlet fluctuations dominate, and become of the same order of magnitude for $`g/\omega =0.5`$ and $`U/t=0.2`$. For larger values of the electron-phonons coupling, the CDW are larger than the singlet fluctuations ($`g/\omega =1`$ and $`U/t=1`$). As for the $`\omega =0.2`$ case, all three correlation functions — spin-spin, charge-charge and singlet-singlet — present a power law behavior as a function of the inter-sites distance for the $`84`$ sites system. The direct calculation of the charge and spin gaps yield and ungaped charge channel and a slightly gaped spin channel with $`\mathrm{\Delta }_\sigma =0.04t`$ for $`U/t=1`$ and $`g/t=1`$, as well as for $`U/t=0.2`$ and $`g/\omega =0.5`$. Again, the smallness of the spin gap explains the power law behavior of the spin-spin correlation functions at the computed distances. This time however the values of $`K_\rho `$ are larger than in the Hubbard model, in agreement with the attractive effective on-site interaction. $`K_\rho 1.1`$ for both $`U/t=1`$ $`g/t=1`$ and $`U/t=0.2`$ $`g/t=0.5`$.
When $`g/\omega `$ increases — but still for negative $`U_{eff}`$ — one goes from the Luther-Emery phase toward a localized bi-polaronic phase through a cross-over. Indeed, for $`U/t=1`$ and $`g/\omega =1.5`$ the system is clearly localized in a bi-polaronic phase with the usual characteristics : nearly constant energy per site, large projection of the wave function on $`\mathrm{\Psi }_{BP}`$ ($`0.94`$ for the 4 sites system), etc… For $`U/t=1`$ and $`g/\omega =1.25`$ however, the system is in an intermediate regime. The delocalization is still reasonably large with an effective hopping of $`0.21t`$ and a projection of the 4 sites system wave function on $`\mathrm{\Psi }_{BP}`$ of only $`0.72`$. In this cross over regime, the system is gaped both in the spin and charge channels, with a spin gap of $`\mathrm{\Delta }_\sigma =0.43`$ and a charge gap of $`\mathrm{\Delta }_\rho =0.84`$. Figure 9 reports the charge, spin, and singlet fluctuations. As expected all of them decrease exponentially with the inter-site distance. When the electron-phonons coupling increases, the system localization increases and the wave function tends toward the totally localized wave function$`\mathrm{\Psi }_{BP}`$. Coherently, the spin gap increases and is expected to follow the same variations as $`|U_{eff}|=2g^2/\omega U`$, since one needs $`|U_{eff}|`$ in order to break the singlet pairing on the sites. In totally localized systems ($`g=+\mathrm{}`$), the charge channel is ungaped since in the absence of any delocalization integral, the different possible choices for the localization of the pairs have the same energy. The charge gap can therefore be expected to be scaling as the rescaled hopping integral $`t_{eff}=t\mathrm{exp}(g^2/\omega ^2)`$. That is exponentially decreasing when the electron-phonons coupling increases.
## IV Conclusion
We have investigated the phase diagram of the Hubbard-Holstein Hamiltonian in a quarter-filled chain for two values of the phonons frequency. These frequencies have be chosen in order to correspond to the low ($`\omega =0.2t`$) and hight ($`\omega =t`$) frequencies of the intra-molecular totally-symmetric vibrations of $`TTF`$, $`TMTSF`$ and related molecules from which the organic conductors are built.
For the $`\omega =t`$ frequency, the phase diagram is very simple. For a positive effective on-site interactions ($`U_{eff}=U2g^2/\omega >0`$), the system is in a Lüttinger-Liquid phase with slightly rescaled parameters compared to the purely electronic system. For negative effective on-site interactions ($`U_{eff}<0`$) and small electron-phonons coupling, the system is in a Luther-Emery phase, however the spin gap remains very small and the coherence length large. The $`K_\rho `$ parameter is again rescaled compared to the purely electronic system, in agreement with the effective model $`U_{eff}=U2g^2/\omega `$, $`t_{eff}=t\mathrm{exp}(g^2/\omega ^2)`$. When the electron-phonons coupling increases, the system goes toward a localized bi-polaronic phase through a soft cross-over.
At low frequency ($`\omega =0.2t`$) however, a rich phase diagram have been found with 5 different phases, some delocalized and some strongly localized. For positive effective interactions one has three different phases. For small electron-phonons coupling and intermediate coupling for small correlation strength, the system is in a Lüttinger Liquid phase. For intermediate electron-phonons coupling, the system is in a phase for which the $`2k_F`$ charge fluctuations disappear and the charge correlation functions have only $`4k_F`$ CDW fluctuations. For large electron-phonons coupling, the system is strongly localized in a polaronic phase. When the effective interaction become negative, one finds a Luther-Emery phase for small electron-phonons coupling, and a strongly localized bi-polaronic phase for intermediate to large couplings. The Luther-Emery phase has a very small spin gap (of the order of $`0.05t`$ and the spin-spin correlation functions exhibit a power law decrease up to at least 60 to 80 sites). Unlike for $`\omega /t=1`$, for $`\omega /t=0.2`$ the delocalized phases do not behave as would the $`U_{eff}=U2g^2/\omega `$, $`t_{eff}=t\mathrm{exp}(g^2/\omega ^2)`$ effective Hubbard model. Indeed the $`K_\rho `$ parameter is strongly decreased due to the electron-phonons interactions while the decrease of the effective repulsion (compared to the pure Hubbard model) would have led us to expect an opposite behavior.
This surprising decrease of the $`K_\rho `$ parameters should be put in context with values of the $`K_\rho `$ parameters found in photo-emission experiments on $`TMTSF`$ and $`TMTTF`$ quasi-one-dimensional organic compounds . Indeed Zwick et al. found a density of states exponent $`\alpha =1/4\left(K_\rho +1/K_\rho 2\right)`$ slightly larger than 1 ($`K_\rho 0.25)`$ in the metallic phase. According to the Lüttinger Liquid theory the system should be insulating for such large values of $`\alpha `$. Tentative explanations have been made by Schwartz et al , involving an effective doping of a Mott insulator due to the inter-chain hopping. On another hand, the electron intra-molecular vibrations interactions are never considered in these systems. For the low frequency modes of these systems, the typical energy scale of the electron-phonons interaction : $`2g^2/\omega `$ is of the order of magnitude of one-half to one-fourth the intra-chain hopping amplitude , that is in the intermediate regime of our phase diagram. For such values of the electron-phonons coupling constants and values of the electronic repulsion between one to four $`t`$, the system is in a Lüttinger liquid phase for which the intra-molecular vibrations strongly decrease the $`K_\rho `$ parameter compared to its electronic value. It is clear that the Hubbard model lacks at least nearest-neighbor bi-electronic repulsion in order to well represent these systems. It can however be expected that the trend toward a strong reduction of the $`K_\rho `$ parameters, observed in the Hubbard-Holstein model, would be similar with longer range interactions, and that the coupling of the electronic degrees of freedom with the intra-molecular vibrations would lead to values of the $`\alpha `$ parameter unreachable in a purely electronic model. Considering these trends, it is reasonable to think the the electron intra-molecular vibrations coupling will have to be considered in order to obtain a realistic description of the observed photo-emission experimental data.
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# Evidence of New States Decaying into Ξ_𝑐'𝜋
## Abstract
Using $`13.7fb^1`$ of data recorded by the CLEO detector at CESR, we report preliminary evidence for two new charmed baryons; one decaying into $`\mathrm{\Xi }_c^0\pi ^+`$ with the subsequent decay $`\mathrm{\Xi }_c^0\mathrm{\Xi }_c^0\gamma `$, and its isospin partner decaying into $`\mathrm{\Xi }_c^+\pi ^{}`$ followed by $`\mathrm{\Xi }_c^+\mathrm{\Xi }_c^+\gamma `$. We measure the following mass differences for the two states: $`M(\mathrm{\Xi }_c^0\gamma \pi ^+)M(\mathrm{\Xi }_c^0)=318.4\pm 1.5\pm 2.9`$ MeV, and $`M(\mathrm{\Xi }_c^+\gamma \pi ^{})M(\mathrm{\Xi }_c^+)=323.9\pm 1.4\pm 3.0`$ MeV. We interpret these new states as the $`J^P=\frac{1}{2}^{}`$ $`\mathrm{\Xi }_{c1}`$ particles, the charmed-strange analogues of the $`\mathrm{\Lambda }_{c1}^+(2593)`$.
preprint: CLEO CONF 00-5 ICHEP00 863
..
Submitted to XXXth International Conference on High Energy Physics, July 2000, Osaka, Japan
P. Avery,<sup>1</sup> C. Prescott,<sup>1</sup> A. I. Rubiera,<sup>1</sup> J. Yelton,<sup>1</sup> J. Zheng,<sup>1</sup> G. Brandenburg,<sup>2</sup> A. Ershov,<sup>2</sup> Y. S. Gao,<sup>2</sup> D. Y.-J. Kim,<sup>2</sup> R. Wilson,<sup>2</sup> T. E. Browder,<sup>3</sup> Y. Li,<sup>3</sup> J. L. Rodriguez,<sup>3</sup> H. Yamamoto,<sup>3</sup> T. Bergfeld,<sup>4</sup> B. I. Eisenstein,<sup>4</sup> J. Ernst,<sup>4</sup> G. E. Gladding,<sup>4</sup> G. D. Gollin,<sup>4</sup> R. M. Hans,<sup>4</sup> E. Johnson,<sup>4</sup> I. Karliner,<sup>4</sup> M. A. Marsh,<sup>4</sup> M. Palmer,<sup>4</sup> C. Plager,<sup>4</sup> C. Sedlack,<sup>4</sup> M. Selen,<sup>4</sup> J. J. Thaler,<sup>4</sup> J. Williams,<sup>4</sup> K. W. Edwards,<sup>5</sup> R. Janicek,<sup>6</sup> P. M. Patel,<sup>6</sup> A. J. Sadoff,<sup>7</sup> R. Ammar,<sup>8</sup> A. Bean,<sup>8</sup> D. Besson,<sup>8</sup> R. Davis,<sup>8</sup> N. Kwak,<sup>8</sup> X. Zhao,<sup>8</sup> S. Anderson,<sup>9</sup> V. V. Frolov,<sup>9</sup> Y. Kubota,<sup>9</sup> S. J. Lee,<sup>9</sup> R. Mahapatra,<sup>9</sup> J. J. O’Neill,<sup>9</sup> R. Poling,<sup>9</sup> T. Riehle,<sup>9</sup> A. Smith,<sup>9</sup> C. J. Stepaniak,<sup>9</sup> J. Urheim,<sup>9</sup> S. Ahmed,<sup>10</sup> M. S. Alam,<sup>10</sup> S. B. Athar,<sup>10</sup> L. Jian,<sup>10</sup> L. Ling,<sup>10</sup> M. Saleem,<sup>10</sup> S. Timm,<sup>10</sup> F. Wappler,<sup>10</sup> A. Anastassov,<sup>11</sup> J. E. Duboscq,<sup>11</sup> E. Eckhart,<sup>11</sup> K. K. Gan,<sup>11</sup> C. Gwon,<sup>11</sup> T. Hart,<sup>11</sup> K. Honscheid,<sup>11</sup> D. Hufnagel,<sup>11</sup> H. Kagan,<sup>11</sup> R. Kass,<sup>11</sup> T. K. Pedlar,<sup>11</sup> H. Schwarthoff,<sup>11</sup> J. B. Thayer,<sup>11</sup> E. von Toerne,<sup>11</sup> M. M. Zoeller,<sup>11</sup> S. J. Richichi,<sup>12</sup> H. Severini,<sup>12</sup> P. Skubic,<sup>12</sup> A. Undrus,<sup>12</sup> S. Chen,<sup>13</sup> J. Fast,<sup>13</sup> J. W. Hinson,<sup>13</sup> J. Lee,<sup>13</sup> D. H. Miller,<sup>13</sup> E. I. Shibata,<sup>13</sup> I. P. J. Shipsey,<sup>13</sup> V. Pavlunin,<sup>13</sup> D. Cronin-Hennessy,<sup>14</sup> A.L. Lyon,<sup>14</sup> E. H. Thorndike,<sup>14</sup> C. P. Jessop,<sup>15</sup> M. L. Perl,<sup>15</sup> V. Savinov,<sup>15</sup> X. Zhou,<sup>15</sup> T. E. Coan,<sup>16</sup> V. Fadeyev,<sup>16</sup> Y. Maravin,<sup>16</sup> I. Narsky,<sup>16</sup> R. Stroynowski,<sup>16</sup> J. Ye,<sup>16</sup> T. Wlodek,<sup>16</sup> M. Artuso,<sup>17</sup> R. Ayad,<sup>17</sup> C. Boulahouache,<sup>17</sup> K. Bukin,<sup>17</sup> E. Dambasuren,<sup>17</sup> S. Karamov,<sup>17</sup> G. Majumder,<sup>17</sup> G. C. Moneti,<sup>17</sup> R. Mountain,<sup>17</sup> S. Schuh,<sup>17</sup> T. Skwarnicki,<sup>17</sup> S. Stone,<sup>17</sup> G. Viehhauser,<sup>17</sup> J.C. Wang,<sup>17</sup> A. Wolf,<sup>17</sup> J. Wu,<sup>17</sup> S. Kopp,<sup>18</sup> A. H. Mahmood,<sup>19</sup> S. E. Csorna,<sup>20</sup> I. Danko,<sup>20</sup> K. W. McLean,<sup>20</sup> Sz. Márka,<sup>20</sup> Z. Xu,<sup>20</sup> R. Godang,<sup>21</sup> K. Kinoshita,<sup>21,</sup><sup>*</sup><sup>*</sup>*Permanent address: University of Cincinnati, Cincinnati, OH 45221 I. C. Lai,<sup>21</sup> S. Schrenk,<sup>21</sup> G. Bonvicini,<sup>22</sup> D. Cinabro,<sup>22</sup> S. McGee,<sup>22</sup> L. P. Perera,<sup>22</sup> G. J. Zhou,<sup>22</sup> E. Lipeles,<sup>23</sup> S. P. Pappas,<sup>23</sup> M. Schmidtler,<sup>23</sup> A. Shapiro,<sup>23</sup> W. M. Sun,<sup>23</sup> A. J. Weinstein,<sup>23</sup> F. Würthwein,<sup>23,</sup>Permanent address: Massachusetts Institute of Technology, Cambridge, MA 02139. D. E. Jaffe,<sup>24</sup> G. Masek,<sup>24</sup> H. P. Paar,<sup>24</sup> E. M. Potter,<sup>24</sup> S. Prell,<sup>24</sup> D. M. Asner,<sup>25</sup> A. Eppich,<sup>25</sup> T. S. Hill,<sup>25</sup> R. J. Morrison,<sup>25</sup> R. A. Briere,<sup>26</sup> G. P. Chen,<sup>26</sup> B. H. Behrens,<sup>27</sup> W. T. Ford,<sup>27</sup> A. Gritsan,<sup>27</sup> J. Roy,<sup>27</sup> J. G. Smith,<sup>27</sup> J. P. Alexander,<sup>28</sup> R. Baker,<sup>28</sup> C. Bebek,<sup>28</sup> B. E. Berger,<sup>28</sup> K. Berkelman,<sup>28</sup> F. Blanc,<sup>28</sup> V. Boisvert,<sup>28</sup> D. G. Cassel,<sup>28</sup> M. Dickson,<sup>28</sup> P. S. Drell,<sup>28</sup> K. M. Ecklund,<sup>28</sup> R. Ehrlich,<sup>28</sup> A. D. Foland,<sup>28</sup> P. Gaidarev,<sup>28</sup> L. Gibbons,<sup>28</sup> B. Gittelman,<sup>28</sup> S. W. Gray,<sup>28</sup> D. L. Hartill,<sup>28</sup> B. K. Heltsley,<sup>28</sup> P. I. Hopman,<sup>28</sup> C. D. Jones,<sup>28</sup> D. L. Kreinick,<sup>28</sup> M. Lohner,<sup>28</sup> A. Magerkurth,<sup>28</sup> T. O. Meyer,<sup>28</sup> N. B. Mistry,<sup>28</sup> E. Nordberg,<sup>28</sup> J. R. Patterson,<sup>28</sup> D. Peterson,<sup>28</sup> D. Riley,<sup>28</sup> J. G. Thayer,<sup>28</sup> D. Urner,<sup>28</sup> B. Valant-Spaight,<sup>28</sup> and A. Warburton<sup>28</sup>
<sup>1</sup>University of Florida, Gainesville, Florida 32611
<sup>2</sup>Harvard University, Cambridge, Massachusetts 02138
<sup>3</sup>University of Hawaii at Manoa, Honolulu, Hawaii 96822
<sup>4</sup>University of Illinois, Urbana-Champaign, Illinois 61801
<sup>5</sup>Carleton University, Ottawa, Ontario, Canada K1S 5B6
and the Institute of Particle Physics, Canada
<sup>6</sup>McGill University, Montréal, Québec, Canada H3A 2T8
and the Institute of Particle Physics, Canada
<sup>7</sup>Ithaca College, Ithaca, New York 14850
<sup>8</sup>University of Kansas, Lawrence, Kansas 66045
<sup>9</sup>University of Minnesota, Minneapolis, Minnesota 55455
<sup>10</sup>State University of New York at Albany, Albany, New York 12222
<sup>11</sup>Ohio State University, Columbus, Ohio 43210
<sup>12</sup>University of Oklahoma, Norman, Oklahoma 73019
<sup>13</sup>Purdue University, West Lafayette, Indiana 47907
<sup>14</sup>University of Rochester, Rochester, New York 14627
<sup>15</sup>Stanford Linear Accelerator Center, Stanford University, Stanford, California 94309
<sup>16</sup>Southern Methodist University, Dallas, Texas 75275
<sup>17</sup>Syracuse University, Syracuse, New York 13244
<sup>18</sup>University of Texas, Austin, TX 78712
<sup>19</sup>University of Texas - Pan American, Edinburg, TX 78539
<sup>20</sup>Vanderbilt University, Nashville, Tennessee 37235
<sup>21</sup>Virginia Polytechnic Institute and State University, Blacksburg, Virginia 24061
<sup>22</sup>Wayne State University, Detroit, Michigan 48202
<sup>23</sup>California Institute of Technology, Pasadena, California 91125
<sup>24</sup>University of California, San Diego, La Jolla, California 92093
<sup>25</sup>University of California, Santa Barbara, California 93106
<sup>26</sup>Carnegie Mellon University, Pittsburgh, Pennsylvania 15213
<sup>27</sup>University of Colorado, Boulder, Colorado 80309-0390
<sup>28</sup>Cornell University, Ithaca, New York 14853
The $`\mathrm{\Xi }_c`$ states consist of a combination of a charm quark, a strange quark and an up or down quark. Each angular momentum combination of these quarks exists as an isospin pair. The ground states, the $`\mathrm{\Xi }_c^0`$ and $`\mathrm{\Xi }_c^+`$, are the only members of the group that decay weakly, and over the past decade their masses, lifetimes, and many of their decay modes have been measured. In 1995 and 1996 CLEO found evidence for a pair of excited states that were interpreted as the $`J^P=\frac{3}{2}^+`$ $`\mathrm{\Xi }_c^{}`$ states, and one of these observations has since been confirmed . In 1999, CLEO discovered the $`\mathrm{\Xi }_c^{}`$ states, which like the ground states have $`J^P=\frac{1}{2}^+`$, but have a wave-function that is symmetric under interchange of the two lighter quarks, and are the charmed-strange analogs of the $`\mathrm{\Sigma }_c`$. Also in 1999, CLEO reported the discovery of a pair of states decaying into $`\mathrm{\Xi }_c^{}\pi `$ with a mass about 348 MeV above the ground states. These are interpreted as the $`J^P=\frac{3}{2}^{}`$ $`\mathrm{\Xi }_{c1}`$ states, the analogs of the $`\mathrm{\Lambda }_{c1}^+(2630)`$, where the numerical subscript refers to the light quark angular momentum. Here, using data from the CLEO II and CLEO II.V detector configurations, we present the first evidence of two peaks corresponding to particles decaying to $`\mathrm{\Xi }_c^{}\pi `$. This is the expected decay mode of the $`J^P=\frac{1}{2}^{}`$ $`\mathrm{\Xi }_{c1}`$ states. That fact, in addition to our measured masses, lead us to identify our peaks with these particles.
The data presented here were taken by the CLEO detector operating at the Cornell Electron Storage Ring. The sample used in this analysis corresponds to an integrated luminosity of 13.7 $`fb^1`$ from data taken on the $`\mathrm{{\rm Y}}(4S)`$ resonance and in the continuum at energies just below the $`\mathrm{{\rm Y}}(4S)`$. Of this data, $`4.7fb^1`$ was taken with the CLEO II configuration and the remainder with the CLEO II.V configuration which includes a silicon vertex detector in its charged particle measurement system. We detected charged tracks with a cylindrical drift chamber system inside a solenoidal magnet. Photons were detected using an electromagnetic calorimeter consisting of 7800 cesium iodide crystals.
We first obtain large samples of reconstructed $`\mathrm{\Xi }_c^+`$ and $`\mathrm{\Xi }_c^0`$ particles, using their decays into $`\mathrm{\Lambda }`$, $`\mathrm{\Xi }^{}`$, $`\mathrm{\Omega }^{}`$ and $`\mathrm{\Xi }^0`$ hyperons as well as kaons, pions and protons<sup>*</sup><sup>*</sup>*Charge conjugate states are implied throughout.. The analysis chain for reconstructing these particles follows closely that presented in our previous publications . We fitted the invariant mass distributions for each decay mode to a sum of a Gaussian signal function and a second order polynomial background. $`\mathrm{\Xi }_c`$ candidates were defined as those combinations within $`2\sigma `$ of the known mass of the $`\mathrm{\Xi }_c^+`$ or $`\mathrm{\Xi }_c^0`$, where $`\sigma `$ is the detector resolution for the detector configuration, calculated mode-by-mode by a GEANT-based Monte Carlo simulation program. To illustrate the good statistics and signal to noise ratio of the $`\mathrm{\Xi }_c`$ signals, and to reduce the combinatorial background, we have placed a cut $`x_p>0.6`$, where $`x_p=p/p_{max}`$, $`p`$ is the momentum of the charmed baryon, $`p_{max}=\sqrt{E_{beam}^2M^2},`$ where $`M`$ is the calculated $`\mathrm{\Xi }_c`$ mass. Table 1 details the number of signal and number of background events obtained from each decay mode. The $`x_p`$ cut used to obtain the results in Table 1 was not used in the final analysis, as we prefer to apply an $`x_p`$ cut only on the $`\mathrm{\Xi }_c^{}\pi `$ combinations.
The $`\mathrm{\Xi }_c`$ candidates defined above were then combined with a photon and the mass differences $`M(\mathrm{\Xi }_c^0\gamma )M(\mathrm{\Xi }_c^0)`$ and $`M(\mathrm{\Xi }_c^+\gamma )M(\mathrm{\Xi }_c^+)`$ were calculated. The transition photons were required to each have energy in excess of 100 MeV, to come from the part of the detector that had the best resolution (cos$`\theta <0.7`$, where $`\theta `$ is the polar angle), and to have an energy profile consistent with being due to an isolated photon. Any photon which, when combined with another photon, made a combination consistent with being a $`\pi ^0`$, was rejected. Those combinations with calculated mass differences within 8 MeV ($`2\sigma `$) of the measured mass differences for the $`\mathrm{\Xi }_c^{}`$ particles, were retained for further analysis. As the $`\mathrm{\Xi }_c^{}`$ decays electromagnetically, its instrinsic width is negligible, and so the candidates were kinematically fit to the $`\mathrm{\Xi }_c^{}`$ masses using these measured mass differences.
We then combine these $`\mathrm{\Xi }_c^{}`$ candidates with an appropriately charged track in the event and plot $`M(\mathrm{\Xi }_c^{}\pi )M(\mathrm{\Xi }_c)`$ for each isospin state. Figure 1(a) shows $`M(\mathrm{\Xi }_c^0\pi ^+)M(\mathrm{\Xi }_c^0)`$, and Figure 1(b) shows $`M(\mathrm{\Xi }_c^+\pi ^{})M(\mathrm{\Xi }_c^+)`$, each with a requirement of $`x_p>0.7`$ on the combination. In both figures there is a peak at about 320 MeV, indicative of the decay of a $`\mathrm{\Xi }_{c1}^+`$ (Figure 1a) and a $`\mathrm{\Xi }_{c1}^0`$ (Figure 1b). We fit each of the two peaks to a sum of Gaussian signal function of floating width, and a polynomial background function. For the $`\mathrm{\Xi }_c^0\pi ^+`$ case, we find a signal of $`18.4_{4.9}^{+5.6}`$ events with a width, $`\sigma `$, of $`5.6\pm 1.7`$ MeV. For the $`\mathrm{\Xi }_c^+\pi ^{}`$ case, we find an excess of $`14.2_{3.9}^{+4.6}`$ events and a width, $`\sigma `$, of $`3.9\pm 1.5`$ MeV. The mass resolution of the detector for these decays are from our Monte Carlo simulation program to be about 1.2 MeV in the CLEO II.V data, and around 1.4 MeV in the CLEO II data. This indicates the likelihood that the states have non-negligible intrinsic widths. We have also fit the plots to Breit-Wigner functions convolved with a Gaussian resolution function. The results of this fit are, $`M(\mathrm{\Xi }_c^0\pi ^+)M(\mathrm{\Xi }_c^0)=318.4\pm 1.5`$ MeV, and $`\mathrm{\Gamma }=7.9_{3.9}^{+5.7}`$ MeV, for Figure 1a, and $`M(\mathrm{\Xi }_c^+\pi ^{})M(\mathrm{\Xi }_c^+)=323.9\pm 1.4`$ MeV, and $`\mathrm{\Gamma }=6.5_{2.8}^{+4.4}`$ MeV for Figure 1b. It is these fits that are superimposed on the data in Figure 1. The results for $`\mathrm{\Gamma }`$ are limited in their accuracy by the low statistics, but indicate that it is very likely that these states have an instrinsic width of the order of MeV. However, we prefer to place 90% confidence level upper limits on the width of these states, of $`\mathrm{\Gamma }(\mathrm{\Xi }_{c1}^+)<16`$ MeV and $`\mathrm{\Gamma }(\mathrm{\Xi }_{c1}^0)<13`$ MeV, dominated by the statistical error.
In order to check that all the $`\mathrm{\Xi }_{c1}`$ decays proceed via an intermediate $`\mathrm{\Xi }_c^{}`$, we release the cuts on $`M(\mathrm{\Xi }_c\gamma )M(\mathrm{\Xi }_c)`$, select combinations within 8 MeV of our final signal peaks and plot $`M(\mathrm{\Xi }_c\gamma )M(\mathrm{\Xi }_c)`$. The plots (Figures 2a and b) show signals which are consistent with our published results for the $`\mathrm{\Xi }_c^{}`$ pair. It is clear that our data are consistent with all the $`J^P=\frac{1}{2}^{}\mathrm{\Xi }_{c1}`$ decays proceeding via an intermediate $`\mathrm{\Xi }_c^{}`$.
When quoting our results as a mass difference with respect to a ground state, our systematic uncertainty is dominated by the uncertainty in the $`M(\mathrm{\Xi }_c^{})M(\mathrm{\Xi }_c)`$ mass differences. Alternatively, we can quote the mass difference with respect to the $`\mathrm{\Xi }_c^{}`$ states, and we find $`M(\mathrm{\Xi }_c^+\pi ^{})M(\mathrm{\Xi }_c^+)=216.1\pm 1.4\pm 1.0`$ MeV and $`M(\mathrm{\Xi }_c^0\pi ^+)M(\mathrm{\Xi }_c^0)=211.4\pm 1.5\pm 1.0`$ MeV. The quoted systematic uncertainties include the spread in our results obtained with different fitting procedures, and an estimate of the systematic uncertainty of our mass difference scale.
The decay patterns of the $`\mathrm{\Xi }_{c1}`$ states should be closely analogous to those of the $`\mathrm{\Lambda }_{c1}^+`$. The preferred decay of the $`J^P=\frac{1}{2}^{}`$ $`\mathrm{\Xi }_{c1}`$ should be to $`\mathrm{\Xi }_c^{}\pi `$ because the spin-parity of the baryons allows this decay to proceed via an $`S`$-wave decay, whereas decays to $`\mathrm{\Xi }_c^{}`$ would have to proceed via a $`D`$-wave. In Heavy Quark Effective Theory (HQET), where the angular momentum and parity of the light di-quark degrees of freedom must be considered separately from those of the heavy quark, decays of the $`\mathrm{\Xi }_{c1}`$ directly to ground state $`\mathrm{\Xi }_c`$ baryons are not allowed. Thus we identify our two peaks as the $`J^P=\frac{1}{2}^{}`$ $`\mathrm{\Xi }_{c1}^0`$ and $`\mathrm{\Xi }_{c1}^+`$. Also in the HQET picture, the splitting between the $`J^P=\frac{1}{2}^{}`$ and $`J^P=\frac{3}{2}^{}`$ states is expected to be similar to that between the two $`\mathrm{\Lambda }_{c1}^+`$ states and also the analagous splitting in charmed meson states. Combining our results found here with our previous results on the $`J^P=\frac{3}{2}^{}\mathrm{\Xi }_{c1}`$ states, and using world average values for the isospin splitting of the ground state $`\mathrm{\Xi }_c`$ baryons, we find the splitting between the $`J^P=\frac{3}{2}^{}`$ and $`J^P=\frac{1}{2}^{}`$ of $`24.8\pm 1.8\pm 3.2`$MeV in the charged case, and $`28.8\pm 1.6\pm 4.0`$MeV for the neutral case. These are similar to the analogous splittings in the $`\mathrm{\Lambda }_c^+`$ and charmed meson systems. We can also measure the isospin splitting between the new states we have find, $`M(\mathrm{\Xi }_{c1}^0)M(\mathrm{\Xi }_{c1}^+)=0.0\pm 2.0\pm 4.5`$MeV, where the quoted systematic uncertainties include the systematic uncertainties of our mass difference measurement as well as the uncertainty in the mass difference of the ground states.
Although the uncertainties are large, this confirms the picture that the excited states of the $`\mathrm{\Xi }_c`$ have smaller isospin splittings than that of the ground states. The $`J^P=\frac{1}{2}\mathrm{\Xi }_{c1}`$ particles are the analogs of the $`\mathrm{\Lambda }_{c1}^+(2593)`$. Although the latter decays with very little phase space, it has been measured to have an instrinsic width of a few MeV. It is not surprising, therefore, that the $`J^P=\frac{1}{2}^{}\mathrm{\Xi }_{c1}`$ pair, which has more phase space available for two-body decays, should also appear as wide peaks in our data.
In conclusion, we present evidence for the production of two new states. The first of these states decay into $`\mathrm{\Xi }_c^0\pi ^+`$ with measured mass given by $`M(\mathrm{\Xi }_c^0\gamma \pi ^+)M(\mathrm{\Xi }_c^0)=318.4\pm 1.5\pm 2.9`$ MeV. The second state decays into $`\mathrm{\Xi }_c^+\pi ^{}`$ with a mass given by $`M(\mathrm{\Xi }_c^+\gamma \pi ^{})M(\mathrm{\Xi }_c^+)`$ = $`323.9\pm 1.4\pm 3.0`$ MeV. Although we do not measure the spin or parity of these states, the observed decay modes, masses, and widths are all consistent with the new states being the $`J^P=\frac{1}{2}^{}`$ $`\mathrm{\Xi }_{c1}^+`$ and $`\mathrm{\Xi }_{c1}^0`$ states, the charmed-strange analogues of the $`\mathrm{\Lambda }_{c1}^+(2593)`$.
We gratefully acknowledge the effort of the CESR staff in providing us with excellent luminosity and running conditions. This work was supported by the National Science Foundation, the U.S. Department of Energy, the Research Corporation, the Natural Sciences and Engineering Research Council of Canada, the A.P. Sloan Foundation, the Swiss National Science Foundation, the Texas Advanced Research Program, and the Alexander von Humboldt Stiftung.
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# QCD SUM RULES FOR HEAVY FLAVORS aafootnote aTalk given by O.Yakovlev at the 4th Workshop on Continuous Advances in QCD, Minneapolis, May 12-14, 2000
## 1 Introduction
The accurate study of $`B`$ meson decays is a main source of information for understanding $`CP`$ violation and the physics of heavy quarks. In particular, experiments at $`B`$ factories will allow measurements of $`B`$ decay properties with good precision $`^\mathrm{?}`$. On the theoretical side, the method of QCD sum rules$`^\mathrm{?}`$ remains one of the main tools in applying Quantum Chromodynamics to hadron physics. Since its birth in 1979, the sum rule method has become more and more advanced not only technically, but also conceptually. In this talk, we give a short review of QCD sum rule results for $`B`$ and $`D`$ mesons and $`\mathrm{\Lambda }_Q`$ and $`\mathrm{\Sigma }_Q`$ baryons. We focus mainly on recent developments concerning semileptonic $`B\pi `$ and $`D\pi `$ transitions$`^\mathrm{?}`$, including a new approach$`^\mathrm{?}`$, which will be discussed in detail, and new results on the $`f^0`$ form factor $`^\mathrm{?}`$. Furthermore we will address pion couplings to $`B`$ and $`D`$ mesons and to $`\mathrm{\Lambda }_Q`$ and $`\mathrm{\Sigma }_Q`$ baryons, meson decay constants and corresponding matrix element for baryons, an estimate of the $`b`$ quark mass from a baryonic sum rule, and finally a recent extraction of the pion distribution amplitude from CLEO data$`^\mathrm{?}`$.
## 2 Pion distribution amplitude from CLEO data
We start with the pion distribution amplitudes, which serve as input in the QCD sum rule method and allow the calculation of heavy-to-light form factors (e.g., $`f_{B\pi }^+`$ and $`f_{D\pi }^+`$) and hadronic coupling constants (e.g., $`g_{B^{}B\pi }`$ and $`g_{D^{}D\pi }`$). Recently, the CLEO collaboration has measured the $`\gamma \gamma ^{}\pi ^0`$ form factor. In this experiment $`^\mathrm{?}`$ <sup>b</sup><sup>b</sup>bThere also exist older results from the CELLO collaboration $`^\mathrm{?}`$., one of the photons is nearly on-shell and the other one is highly off-shell, with a virtuality in the range $`1.5`$ GeV<sup>2</sup>$`9.2`$ GeV<sup>2</sup>. The possibility of extracting the twist-2 pion distribution amplitude from the CLEO data has been studied in the papers $`^{\mathrm{?},\mathrm{?}}`$. There, the light-cone sum rule (LCSR) method has been used to calculate the relevant form factor and to compare the calculation with the measurement of $`\gamma \gamma ^{}\pi ^0`$.
In order to sketch the basic idea we begin with the correlator of two vector currents $`j_\mu =(\frac{2}{3}\overline{u}\gamma _\mu u\frac{1}{3}\overline{d}\gamma _\mu d)`$:
$`{\displaystyle d^4xe^{iq_1x}\pi ^0(0)|T\{j_\mu (x)j_\nu (0)\}|0}=iϵ_{\mu \nu \alpha \beta }q_1^\alpha q_2^\beta F^{\pi \gamma ^{}\gamma ^{}}(s_1,s_2),`$ (1)
where $`q_1,q_2`$ are the momenta of the photons, and $`s_1=q_1^2`$, $`s_2=q_2^2`$ are the virtualities. In the CLEO data, one of the virtualities is small, i.e. $`s_20`$. Since a straightforward OPE calculation is impossible, we have to use analyticity and duality arguments. One can write the form factor as a dispersion relation in $`s_2`$:
$$F^{\pi \gamma ^{}\gamma ^{}}(s_1,s_2)=\frac{\sqrt{2}f_\rho F^{\rho \pi }(s_1)}{m_\rho ^2s_2}+\underset{s_0}{\overset{\mathrm{}}{}}𝑑s\frac{\rho ^h(s_1,s)}{ss_2}.$$
(2)
For the physical ground states $`\rho `$ and $`\omega `$ we take $`m_\rho m_\omega `$; $`\frac{1}{3}\pi ^0(p)|j_\mu |\omega (q_2)\pi ^0(p)|j_\mu |\rho ^0(q_2)=\frac{1}{m_\rho }ϵ_{\mu \nu \alpha \beta }e^\nu q_1^\alpha q_2^\beta F^{\rho \pi }(s_1)`$; $`3\omega |j_\nu |0\rho ^0|j_\nu |0=\frac{f_\rho }{\sqrt{2}}m_\rho e_\nu ^{}`$, $`e_\nu `$ being the polarization vector of the $`\rho `$ meson and $`f_\rho `$ being the decay constant. The spectral density of the higher energy states $`\rho ^h(s_1,s)`$ is derived from the expression for $`F_{QCD}^{\pi \gamma ^{}\gamma ^{}}(s_1,s)`$ calculated in QCD, assuming semi-local quark-hadron duality for $`s>s_0`$. Equating the dispersion relation (2) with the QCD expression at large $`s_2`$, and performing a Borel transformation in $`s_2`$, one gets the LCSR:
$$\sqrt{2}f_\rho F^{\rho \pi }(s_1)=\frac{1}{\pi }\underset{0}{\overset{s_0}{}}𝑑s\text{Im}F_{QCD}^{\pi \gamma ^{}\gamma ^{}}(s_1,s)\text{e}^{\frac{m_\rho ^2s}{M^2}},$$
(3)
where $`M`$ is the Borel parameter. Substituting (3) into (2) and taking $`s_20`$ one finally obtains
$$F^{\pi \gamma \gamma ^{}}(s_1)=\frac{1}{\pi m_\rho ^2}\underset{0}{\overset{s_0}{}}𝑑s\text{Im}F_{QCD}^{\pi \gamma ^{}\gamma ^{}}(s_1,s)\text{e}^{\frac{m_\rho ^2s}{M^2}}+\frac{1}{\pi }\underset{s_0}{\overset{\mathrm{}}{}}\frac{ds}{s}\text{Im}F_{QCD}^{\pi \gamma ^{}\gamma ^{}}(s_1,s).$$
(4)
This expression is the basic sum rule used for the numerical analysis. The calculation of the spectral density of the twist-2 operator including the $`𝒪(\alpha _S)`$ radiative correction gives $`^\mathrm{?}`$
$`{\displaystyle \frac{1}{\pi }}`$ $`\text{Im}_{s_2}F^{\pi \gamma ^{}\gamma ^{}}(s_1,s_2)={\displaystyle \frac{2\sqrt{2}f_\pi s_2s_1}{\left(s_2s_1\right)^3}}(1+{\displaystyle \frac{\alpha _s(\mu )C_F}{12\pi }}`$
$``$ $`(15+\pi ^23\mathrm{log}^2({\displaystyle \frac{s_2}{s_1}}))+a_2(\mu )A_2(s_2,s_1)+a_4(\mu )A_4(s_2,s_1)).`$
As usual, the distribution amplitude of twist 2 is expanded in Gegenbauer polynomials, keeping only the first three terms: $`\phi _\pi =6u(1u)\left(1+a_2C_2^{3/2}+a_4C_4^{3/2}\right).`$ The coefficients $`A_{2,4}`$ in (2) are too complicated to be given here. They can be found in $`^\mathrm{?}`$.
We then combined the twist-2 contribution at NLO with the higher twist contributions up to twist 4, calculated in $`^\mathrm{?}`$, and analyzed the LCSR for the form factor of the process $`\gamma \gamma ^{}\pi ^0`$ numerically.
Details of the analysis are given in $`^\mathrm{?}`$. The coefficients $`a_2`$ and $`a_4`$ of the twist-2 distribution amplitude can be determined by comparing the sum rule (4) with CLEO data $`^\mathrm{?}`$. We find that the deviation of the pion distribution amplitude from the asymptotic form is small. More definitely, putting $`a_4=0`$, we get$`^\mathrm{?}`$
$`a_2(\mu )=0.12\pm 0.03\text{at}\mu =2.4\text{GeV}.`$ (6)
This result agrees well with a recent analysis of the electromagnetic pion form factor $`^\mathrm{?}`$. Fig. 1 shows the form factor $`Q^2F^{\gamma \gamma ^{}\pi }(Q^2)`$ calculated with different distribution amplitudes: Braun-Filyanov$`^\mathrm{?}`$ (dashed lines), Chernyak-Zhitnitsky$`^\mathrm{?}`$ (dotted lines) and (6). In principle, one can also extract the coefficient $`a_4`$. Unfortunately, the present data is not good enough to fix the values of $`a_2`$ and $`a_4`$ simultaneously. The ranges of $`a_2`$ and $`a_4`$ favored by CLEO data are shown in Fig. 1. Obviously, these are in qualitative agreement with (6) and also with the results derived in $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$, where it has also been claimed that the pion distribution amplitude is very close to the asymptotic form.
## 3 Coupling constants $`g_{B^{}B\pi }`$ and $`g_{D^{}D\pi }`$
The hadronic $`B^{}B\pi `$ coupling is defined by the on-shell matrix element
$$\overline{B}^0(p)\pi ^{}(q)B^{}(p+q)=g_{B^{}B\pi }(qϵ),$$
(7)
where the meson four-momenta are given in brackets and $`ϵ_\mu `$ is the polarization vector of the $`B^{}`$. An analogous definition holds for the $`D^{}D\pi `$ coupling. These couplings play an important role in $`B`$ and $`D`$ physics. For example, they determine the magnitude of the weak $`B\pi `$ and $`D\pi `$ form factors at zero pion recoil. Moreover, the coupling constant $`g_{D^{}D\pi }`$ is directly related to the decay width of $`D^{}D\pi `$. The decay $`B^{}B\pi `$ is kinematically forbidden. Theoretically, the $`B^{}B\pi `$ and $`D^{}D\pi `$ couplings have been studied using a variety of methods<sup>c</sup><sup>c</sup>cAn overview is given, e.g., in Tab. 1 of ref. $`^\mathrm{?}`$.. Among these, QCD light-cone sum rules (LCSR) have proved particularly powerful. The LCSR calculations of $`g_{B^{}B\pi }`$ and $`g_{D^{}D\pi }`$ including perturbative QCD effects in LO and NLO were reported in $`^{\mathrm{?},\mathrm{?}}`$. The final LCSR reads
$`f_Bf_B^{}g_{B^{}B\pi }={\displaystyle \frac{m_b^2f_\pi }{m_B^2m_B^{}}}e^{\frac{m_B^2+m_B^{}^2}{2M^2}}[M^2(e^{\frac{m_b^2}{M^2}}e^{\frac{s_0^B}{M^2}})\phi _\pi (1/2,\mu )`$
$`+{\displaystyle \frac{\alpha _sC_F}{4\pi }}{\displaystyle \underset{2m_b^2}{\overset{2s_0^B}{}}}f({\displaystyle \frac{s}{m_b^2}}2)e^{\frac{s}{2M^2}}ds+F^{(3,4)}(M^2,m_b^2,s_0^B,\mu )]`$ (8)
with
$`f(x)`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{4}}+3\mathrm{ln}\left({\displaystyle \frac{x}{2}}\right)\mathrm{ln}\left(1+{\displaystyle \frac{x}{2}}\right){\displaystyle \frac{3(3x^3+22x^2+40x+24)}{2(2+x)^3}}\mathrm{ln}\left({\displaystyle \frac{x}{2}}\right)`$ (9)
$`+6\text{Li}_2\left({\displaystyle \frac{x}{2}}\right)3\text{Li}_2(x)3\text{Li}_2(x1)3\mathrm{ln}(1+x)\mathrm{ln}(2+x)`$
$`{\displaystyle \frac{3(3x^2+20x+20)}{4(2+x)^2}}+{\displaystyle \frac{6x(1+x)\mathrm{ln}(1+x)}{(2+x)^3}}.`$
In (8), we have added the contributions $`F^{(3,4)}`$ from the pion distribution amplitudes of twist 3 and 4, which can be found in $`^\mathrm{?}`$.
For the $`b`$-quark mass and the corresponding continuum threshold we use $`m_b=4.7\pm 0.1`$ GeV and $`s_0^B=352`$GeV<sup>2</sup> respectively. The running coupling constant is taken in the two-loop approximation with $`N_f=4`$ and $`\mathrm{\Lambda }_{\overline{MS}}^{(4)}=315`$ MeV corresponding to $`\alpha _s(m_Z)=0.118`$ $`^\mathrm{?}`$. In the charm case, the corresponding parameters are given by $`m_c=1.3\pm 0.1`$ GeV, $`s_0^D=6`$ GeV<sup>2</sup>, and $`\mathrm{\Lambda }_{\overline{MS}}^{(3)}=380`$ MeV <sup>d</sup><sup>d</sup>dThe meson masses are (in GeV) $`m_B=5.279,m_B^{}=5.325,m_D=1.87,m_D^{}=2.01,`$ and $`f_\pi =132`$ MeV. . Finally, for the pion distribution amplitude $`\phi _\pi (u,\mu )`$ at $`u=0.5`$ and $`\mu =2.4`$ GeV we have $`\phi _\pi (1/2,\mu )=1.23`$ $`^\mathrm{?}`$. The decay constants and resulting coupling constants are summarized in the two tables in Fig. 2. We will make use of them in the next section. Here we just mention that from $`g_{D^{}D\pi }`$ given in Table 2 one obtains $`\mathrm{\Gamma }(D^+D^0\pi ^+)=23\pm 13\text{keV}`$. The current experimental limit $`^\mathrm{?}`$ $`\mathrm{\Gamma }(D^+D^0\pi ^+)<89\text{keV}`$ is still too high to challenge the theoretical prediction.
## 4 The scalar form factor $`f^0`$
In general, the hadronic matrix element of the $`B\pi `$ transition is determined by two independent form factors, $`f^+`$ and $`f^{}`$:
$`\pi (q)|\overline{u}\gamma _\mu b|B(p+q)=2f^+(p^2)q_\mu +\left(f^+(p^2)+f^{}(p^2)\right)p_\mu ,`$ (10)
where $`(p+q)`$ and $`q`$ denote the initial and final state four-momenta and $`\overline{u}\gamma _\mu b`$ is the relevant weak current. The form factor $`f^0`$ is usually defined through the matrix element
$`p^\mu \pi (q)|\overline{u}\gamma _\mu b|B(p+q)=f^0(p^2)(m_B^2m_\pi ^2),`$ (11)
yielding together with (10) $`f^0=f^++\frac{p^2}{m_B^2m_\pi ^2}f^{}.`$ In order to determine $`f^0`$ from sum rules it is advantageous to consider $`f^+`$ and $`f^++f^{}`$. The sum rule for $`f^+`$ has been analysed in $`^{\mathrm{?},\mathrm{?}}`$. The sum rule for the sum of form factors is given by
$`f^++f^{}={\displaystyle \frac{m_bf_\pi }{\pi m_B^2f_B}}{\displaystyle \underset{m_b^2}{\overset{s_0}{}}}𝑑s{\displaystyle \underset{0}{\overset{1}{}}}𝑑u\mathrm{exp}\left({\displaystyle \frac{sm_B^2}{M^2}}\right)\phi _\pi (u)\text{Im}\stackrel{~}{T}_{QCD},`$ (12)
where $`M`$ again denotes the Borel mass. The expression of the hard amplitude $`\stackrel{~}{T}_{QCD}(p^2,s,u,\mu )`$ in LO and NLO can be found in $`^\mathrm{?}`$ and in$`^\mathrm{?}`$, respectively. The leading twist-2 contribution to the imaginary part of the hard amplitude is given by $`^\mathrm{?}`$
$`{\displaystyle \frac{1}{\pi }}\text{Im}\stackrel{~}{T}_{QCD}(s_1,s_2,u,\mu )=\left({\displaystyle \frac{C_F\alpha _s(\mu )}{2\pi }}\right)\mathrm{\Theta }(s_2m_b^2){\displaystyle \frac{m_b}{s_2s_1}}`$
$`\{\mathrm{\Theta }(uu_0)[{\displaystyle \frac{(1u)(uu_0)(s_2s_1)^2}{2u\rho ^2}}{\displaystyle \frac{1}{u(1u)}}({\displaystyle \frac{m^2}{\rho }}1)]`$
$`+`$ $`\delta (uu_0){\displaystyle \frac{1}{2u}}[{\displaystyle \frac{(s_1m_b^2)^2}{s_1^2}}\mathrm{ln}(1{\displaystyle \frac{s_1}{m_b^2}})+{\displaystyle \frac{m_b^2}{s_1}}1]{\displaystyle \frac{1}{1u}}(1{\displaystyle \frac{m_b^2}{s_2}})\}`$
with $`u_0=\frac{m_b^2s_1}{s_2s_1}`$. In Fig. 3, the form factor $`f_{B\pi }^0`$ is plotted together with the UKQCD lattice results $`^\mathrm{?}`$. We see that the radiative contributions improve the agreement between the lattice and the LCSR calculations. Also shown in Fig. 3 are the LCSR results for the form factor $`f_{D\pi }^0`$. We note that $`f_{D\pi }^0(0)=0.66`$.
## 5 New method of calculating $`f^+`$
In this section we review a new method suggested in $`^\mathrm{?}`$ for calculating heavy-to-light form factors. The method is based on first principles. It is an extension of LCSR, but it has a much wider range of applicability, including the intermediate momentum region, where most of the lattice results are located, and even the region near zero recoil. The main idea is to use the operator product expansion with a combination of double and single dispersion relations. The resulting new sum rule has a term which corresponds to the ground-state, as well as contributions which account for all possible physical intermediate states. We start from the usual correlation function
$`F_\mu (p,q)`$ $`=`$ $`i{\displaystyle 𝑑xe^{ipx}\pi (q)|T\{\overline{u}(x)\gamma _\mu b(x),m_b\overline{b}(0)i\gamma _5d(0)\}|0}`$ (14)
$`=`$ $`F(p^2,(p+q)^2)q_\mu +\stackrel{~}{F}(p^2,(p+q)^2)p_\mu ,`$
focusing on the invariant amplitude $`F(p^2,(p+q)^2)`$. In the following, we use the definitions
$`\sigma (p^2,s_2)={\displaystyle \frac{1}{\pi }}\text{Im}_{s_2}F(p^2,s_2),\rho (s_1,s_2)={\displaystyle \frac{1}{\pi ^2}}\text{Im}_{s_1}\text{Im}_{s_2}F(s_1,s_2).`$ (15)
The standard sum rule for the form factor $`f^+(p^2)`$ is obtained by writing a single dispersion relation for $`F(p^2,(p+q)^2)`$ in the $`(p+q)^2`$-channel, inserting the hadronic representation for $`\sigma (p^2,s_2)`$ and Borelizing in $`(p+q)^2`$:
$`_{(p+q)^2}F`$ $`=`$ $`_{(p+q)^2}\left({\displaystyle \frac{2m_B^2f_Bf^+(p^2)}{m_B^2(p+q)^2}}+{\displaystyle \underset{s_2>s_0}{}}𝑑s_2{\displaystyle \frac{\sigma ^{hadr}(p^2,s_2)}{s_2(p+q)^2}}\right).`$ (16)
Note that any subtraction terms which might appear vanish after Borelization. Similarly, the standard light-cone sum rule for the coupling $`g_{B^{}B\pi }`$ is obtained from a double dispersion relation:
$`_{p^2}_{(p+q)^2}F`$ $`=`$ $`_{p^2}_{(p+q)^2}({\displaystyle \frac{m_B^2m_B^{}f_Bf_B^{}g_{B^{}B\pi }}{(p^2m_B^{}^2)((p+q)^2m_B^2)}}`$ (17)
$`+{\displaystyle \underset{\mathrm{\Sigma }}{}}ds_1ds_2{\displaystyle \frac{\rho ^{hadr}(s_1,s_2)}{(s_1p^2)(s_2(p+q)^2)}}),`$
where $`\mathrm{\Sigma }`$ denotes the integration region defined by $`s_1>s_0`$, $`s_2>m_b^2`$ and $`s_1>m_b^2`$, $`s_2>s_0`$.
In contrast to the above procedure we suggest to use a dispersion relation for $`\sigma (p^2,s_2)/(p^2)^l`$ in the $`p^2`$-channel (with $`l`$ being an integer):
$`\sigma (p^2,s_2)`$ $`=`$ $`{\displaystyle \frac{1}{(l1)!}}\left(p^2\right)^l{\displaystyle \frac{d^{l1}}{ds_1^{l1}}}{\displaystyle \frac{\sigma (s_1,s_2)}{s_1p^2}}|_{s_1=0}+{\displaystyle \underset{s_1>m_b^2}{}}𝑑s_1{\displaystyle \frac{(p^2)^l}{s_1^l}}{\displaystyle \frac{\rho (s_1,s_2)}{s_1p^2}},`$ (18)
and to replace $`\sigma (p^2,s_2)`$ in (16) by the r.h.s of (18) <sup>e</sup><sup>e</sup>eBy choosing $`l`$ large enough the dispersion relation (18) will be convergent.. Then, writing a double dispersion relation for $`F(p^2,(p+q)^2)/(p^2)^l`$ and comparing it with the previous result, we obtain the sum rule
$`f^+(p^2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{(p^2)^l}{(m_B^{}^2)^l}}{\displaystyle \frac{f_B^{}g_{B^{}B\pi }}{m_B^{}\left(1\frac{p^2}{m_B^{}^2}\right)}}{\displaystyle \frac{1}{(l1)!}}\left(p^2\right)^l{\displaystyle \frac{d^{l1}}{ds_1^{l1}}}{\displaystyle \frac{f^+(s_1)}{s_1p^2}}|_{s_1=0}`$ (19)
$`+{\displaystyle \frac{1}{2m_B^2f_B}}{\displaystyle \underset{\mathrm{\Sigma }^{}}{}}𝑑s_1𝑑s_2{\displaystyle \frac{(p^2)^l}{s_1^l}}{\displaystyle \frac{\rho (s_1,s_2)}{s_1p^2}}e^{\frac{s_2m_B^2}{M^2}},`$
where the integration region $`\mathrm{\Sigma }^{}`$ is defined by $`s_1>s_0`$ and $`m_b^2<s_2<s_0`$. This sum rule is valid in the whole kinematical range of $`p^2`$. As input we need the first $`(l1)`$ terms of the Taylor expansion of $`f^+(p^2)`$ around $`p^2=0`$. These parameters can be obtained numerically from the standard sum rule for $`f^+(p^2)`$:
$`f^+(p^2)`$ $`=`$ $`{\displaystyle \frac{1}{2m_B^2f_B}}{\displaystyle \underset{m_b^2}{\overset{s_0}{}}}\sigma ^{QCD}(p^2,s_2)e^{\frac{s_2m_B^2}{M^2}}`$ (20)
following from (16). We further need the residue at the pole $`p^2=m_B^{}^2`$, which can be obtained from the sum rule (17), as discussed in the previous section (see (8)).
The case $`l=0`$ has a very transparent physical meaning. The first term represents the contribution of the ground state resonance with mass $`m_B^{}`$, while the second term corresponds to the contributions of all other physical states in this channel. As shown explicitly in$`^\mathrm{?}`$ for twist 2, 3 and 4, the last term is of $`O(\alpha _s)`$ only. This provides an explanation of the empirical fact that the single pole model describes many form factors quite well. In addition, we are now able to quantify the deviation from the pole model, in a model-independent way, applying QCD and light-cone OPE. It should be noted that the parameter $`l`$ plays a similar role as the Borel parameter $`M^2`$. There is a lower limit on $`l`$ such that the dispersion relation (18) converges. Going to higher values of $`l`$ will improve the convergence of the dispersion relations and will suppress higher resonances in the $`B^{}`$-channel. But there is also an upper limit on $`l`$. The higher the value of $`l`$, the more derivatives of $`f^+(p^2)`$ at $`p^2=0`$ enter. At some point, one starts probing the region $`p^2>m_b^22\chi m_b`$, where the standard sum rule (20) breaks down.
Details on the numerical analysis of the new sum rule can be found in $`^\mathrm{?}`$. This analysis nicely supports the qualitative results obtained in $`^\mathrm{?}`$. Using the convenient parameterization$`^\mathrm{?}`$
$$f_{B\pi }^+(p^2)=\frac{f_{B\pi }^+(0)}{(1p^2/m_B^{}^2)(1\alpha _{B\pi }p^2/m_B^{}^2)},$$
(21)
and $`f_{B\pi }^+(0)=0.28\pm 0.05`$$`^{\mathrm{?},\mathrm{?}}`$, we get $`\alpha _{B\pi }=0.4\pm 0.04`$ in remarkable agreement with $`\alpha _{B\pi }=0.32\pm _{0.07}^{0.21}`$ derived in $`^\mathrm{?}`$. Fig. 4 shows a comparison of (21) with recent lattice results $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$. The agreement within uncertainties is very satisfactory. Finally, the LCSR prediction also obeys the constraints derived from sum rules for the inclusive semileptonic decay width in the heavy quark limit $`^\mathrm{?}`$. This is also demonstrated in Fig. 4.
The above results on $`f_{B\pi }^+`$ can be used to calculate the width of the semileptonic decay $`B\pi \overline{l}\nu _l`$ with $`l=e,\mu `$. For the integrated width, one obtains$`^\mathrm{?}`$
$$\mathrm{\Gamma }=\frac{G^2|V_{ub}|^2}{24\pi ^3}𝑑p^2(E_\pi ^2m_\pi ^2)^{3/2}\left[f_{B\pi }^+(p^2)\right]^2=(7.3\pm 2.5)|V_{ub}|^2\text{ps}^1.$$
(22)
Experimentally, combining the branching ratio $`BR(B^0\pi ^{}l^+\nu _l)=(1.8\pm 0.6)10^4`$ with the $`B^0`$ lifetime $`\tau _{B^0}=1.54\pm 0.03`$ ps $`^\mathrm{?}`$, one gets $`\mathrm{\Gamma }(B^0\pi ^{}l^+\nu _l)=(1.17\pm 0.39)10^4\text{ps}^1.`$ From that and (22) one can then determine the quark mixing parameter $`|V_{ub}|`$. The result is
$$|V_{ub}|=(4.0\pm 0.7\pm 0.7)10^3$$
(23)
with the experimental error and theoretical uncertainty given in this order. For the $`D\pi `$ transition and using (21) analogously one obtains $`^\mathrm{?}`$ $`\alpha _{D\pi }=0.01_{0.07}^{+0.11}`$ and $`f_{D\pi }^+(0)=0.65\pm 0.11,`$ which nicely agrees with lattice estimates, for example, the world average $`^\mathrm{?}`$ $`f_{D\pi }^+(0)=0.65\pm 0.10,`$ or the most recent APE result $`^\mathrm{?}`$, $`f_{D\pi }^+(0)=0.64\pm 0.05_{.07}^{+.00}`$. For more details one should consult$`^\mathrm{?}`$.
## 6 Heavy baryons
The study of heavy baryons such as $`\mathrm{\Lambda }_b,\mathrm{\Lambda }_c`$ and $`\mathrm{\Sigma }_b,\mathrm{\Sigma }_c`$, is more complicated. Two and three point QCD sum rules have been investigated in $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$, and have been applied to the heavy-to-light baryon transitions$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. However, we are not aware of applications of LCSR to heavy-to-light baryon transitions. In the following we collect the results available at present (see also $`^\mathrm{?}`$). One has estimated the binding energies, $`\overline{\mathrm{\Lambda }}=Mm_Q`$, of the ground state baryons and the residues of the baryonic currents, $`B|J_B|0=F_Bu_B`$, at NLO. The results are$`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$
$$\overline{\mathrm{\Lambda }}(\mathrm{\Lambda }_Q)=0.77\pm 0.05\text{GeV}\text{and}|F_{\mathrm{\Lambda }_Q}|=0.027\pm 0.001\text{GeV}^3,$$
(24)
$$\overline{\mathrm{\Lambda }}(\mathrm{\Sigma }_Q)=0.94\pm 0.05\text{GeV}\text{and}|F_{\mathrm{\Sigma }_Q}|=0.038\pm 0.003\text{GeV}^3.$$
(25)
Using the experimental value for the mass of the $`\mathrm{\Lambda }_b`$ baryon, $`m(\mathrm{\Lambda }_b)=5.642\pm 0.05`$ GeV $`^\mathrm{?}`$ one finds for the pole mass of the $`b`$ quark: $`m_b=4.88\pm 0.1`$, and for the related $`\overline{MS}`$ mass: $`\overline{m}(\overline{m})=4.25\pm 0.1\text{GeV}.`$ These values are in good agreement with the mass estimates in the meson case$`^{\mathrm{?},\mathrm{?}}`$. Coupling constants have been derived from sum rules in the external axial field $`^\mathrm{?}`$ with the result: $`g_{\mathrm{\Sigma }^{}\mathrm{\Sigma }\pi }=0.83\pm 0.3`$ and $`g_{\mathrm{\Sigma }^{}\mathrm{\Lambda }\pi }=0.58\pm 0.2`$. The semileptonic transition $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ has been studied$`^\mathrm{?}`$ using sum rule techniques. The matrix elements of this weak transition are determined by the Isgur-Wise function
$$\mathrm{\Lambda }_b|\overline{c}\mathrm{\Gamma }b|\mathrm{\Lambda }_c=\xi (w)\overline{u}_c\mathrm{\Gamma }u_b$$
(26)
where $`w=v_bv_c`$ and $`\mathrm{\Gamma }`$ is the Dirac matrix. The baryonic Isgur-Wise function has been estimated by QCD sum rules in $`^\mathrm{?}`$. The slope of this function at $`w=1`$ is found to be $`\rho ^2=1.15\pm 0.2`$ and the shape is fitted very well by $`\xi (w)=\frac{2}{1+w}exp\left((2\rho ^21)\frac{w1}{w+1}\right).`$ Taking into account $`1/m`$ corrections$`^\mathrm{?}`$, one obtains the width $`\mathrm{\Gamma }(\mathrm{\Lambda }_b\mathrm{\Lambda }_ce\nu )=|\frac{V_{cb}}{0.04}|^2610^{14}\text{GeV}.`$
Conclusion: We have given a short review of selected topics concerning QCD sum rules for heavy hadrons. One main development during the past few years has been the NLO improvement. A second important development is the ongoing update of the pion wave functions $`^\mathrm{?}`$. Very recently a new sum rule for $`B\pi `$ has been suggested$`^\mathrm{?}`$. Among the remaining problems, we have mentioned the application of LCSR to baryons requiring the knowledge of baryonic distribution amplitudes.
Acknowledgments: O.Y. acknowledges support from the US Department of Energy. R.R. acknowledges support from the Bundesministerium für Bildung und Forschung under contract number 05HT9WWA9.
## References
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# Covariant and locally Lorentz-invariant varying speed of light theories
## I Introduction
The varying speed of light (VSL) theory provides an elegant solution to the cosmological problems - the horizon, flatness, and Lambda problems of Big-Bang cosmology. The theory has appeared in several guises ()), but in the formulation proposed by Albrecht and Magueijo (see also ) one finds the most direct mechanism for converting the Einstein deSitter model into a cosmological attractor. Unfortunately the foundations of such a theory are far from solid. Covariance and local Lorentz invariance are explicitly broken, and are not replaced by similar far-reaching principles. The difficulty in applying the theory to situations other than cosmology (eg. black holes) stems directly from this deficiency.
This paper is an attempt to remedy this shortcoming. This may be achieved in various different ways, some of which inevitably rather radical. We note that nothing prevents the construction of a theory satisfying the principle of relativity, while still allowing for space-time variations in $`c`$. Such a theory would in general not be Lorentz invariant, but it could still be relativistic. Indeed, Lorentz invariance follows from two independent postulates: the principle of relativity and the principle of constancy of the speed of light. Dropping the latter while keeping the former leads to a new invariance, known as Fock-Lorentz symmetry . This invariance does not distinguish between inertial frames (and therefore satisfies the principle of relativity) but it allows for a varying $`c`$; indeed it allows for a non-invariant $`c`$.
A possible approach is therefore to set up a theory of gravitation based upon a gauged Fock-Lorentz symmetry. However we note that such an enterprise accommodates more than is required by VSL theories: it allows the speed of light at a given point to depend on the observer’s speed. Also the speed of light in the Fock-Lorentz space is anisotropic. Clearly, certain aspects of the second postulate of Einstein’s relativity theory may be kept in the simplest VSL theories, namely that the speed of light at a given point be independent of its color, direction, or the speeds of either emitter or observer.
In this paper we shall be as conservative as possible and preserve all aspects of the second postulate of special relativity consistent with allowing space-time variations of $`c`$. In Section II we show that such a reformulation gleans from the second postulate of relativity all that is operationally meaningful, in the sense that the aspects of the second postulate which we preserve are exactly those which can be the outcome of experiment (such as the Michelson-Morley experiment). The constancy of $`c`$ in space-time, on the other hand, amounts to nothing more than a definition of a system of units. In Section II and III we show that such a theory is locally Lorentz invariant and generally covariant, subject to a minimal generalization of these concepts. In Section IV we summarise the overall structure of such theories, and the basic reasons for adopting it.
We then discuss Lagrangians governing the dynamics of these theories. The main practical drawback of explicit lack of covariance is that it makes an action principle formulation rather awkward (see ). The VSL theories proposed in this paper, on the contrary, are easily amenable to an action principle formulation. However we shall try to borrow some features from earlier models, such as lack of energy conservation.
In Section V we first consider the matter action. We show how it is always possible to define the matter Lagrangian so that $`c`$ does not appear explicitly. Such a principle fixes a large number of scaling laws for other “constants” as a function of $`c`$. It also leads to simpler dynamical equations for $`c`$.
Two constants are left undetermined by these considerations: Planck’s constant $`\mathrm{}`$ and Boltzmann’s constant $`k_B`$. These cannot be determined by classical dynamics, and scaling laws $`\mathrm{}(c)`$ and $`k_B(c)`$ should be postulated. In Section VI we consider the implications of various $`\mathrm{}(c)`$ for quantization. We identify situations in which a varying speed of light leads to quantum particle creation.
Then in Section VII we consider Lagrangians for gravitational dynamics (we include the dynamics of $`c`$ into this discussion - as the field $`c`$ can be seen as an extra gravitational field). We identify the actions which lead to nothing but a change of units in a standard Brans-Dicke theory; all other actions are intrinsically different theories.
The rest of the paper is devoted to the simplest applications of these theories. In Section IX we discuss empty space solutions. We find a variation in $`c`$ and a global space-time which is very similar to those found in Fock-Lorentz space. We also show how Fock-Lorentz space is nothing but a change of units applied to Minkowski space-time. However $`t=\mathrm{}`$ is brought to a finite time in the varying $`c`$ representation. We show that the space is actually extendable beyond this finite time; into what in the fixed $`c`$ representation would be a trans-eternal region.
Another flat-space solution to our theory is a soliton string, close to which the speed of light is much larger. We label it a fast-track. A spaceship moving along a fast track could move at non-relativistic speeds, without a twin paradox effect, and still cover enormous intergalactic distances. These solutions are not dissimilar to gravitational wormholes; and indeed they are mapped into wormhole-like structures in fixed $`c`$ units.
We finally discuss general features of cosmological solutions and black holes in these theories. These will be developed further in two publications currently under preparation . Concerning black holes the main novelty is that for some regions of the couplings the speed of light may go to zero at the horizon. This effectively prevents any observer from entering the horizon, and its interior should therefore be excised from the manifold. We relabel this boundary an “edge”, and comment on the implication of this effect for a generalized cosmic censorship principle.
## II Generalized Lorentz invariance
From an operational point of view all laws of physics should be invariant under global and local changes of units . Indeed measurements are always ratios to standard units, and therefore represent essentially dimensionless quantities. Physics should therefore be dimensionless or unit-invariant. However, this far-reaching principle is rarely incorporated into theoretical constructions, because a concrete choice of units usually simplifies the statement of laws. While this practical consideration should be recognized, it is important to realize that some theoretical constructions are tautological, and amount to nothing more than the specification of a system of units.
An example is the second postulate of special relativity: the constancy of $`c`$. Clearly the postulate is invariant under unit transformations when it states that light of different colors travels at the same speed - as it makes a statement about the ratio of two speeds at a given point, which is a dimensionless quantity. The postulate is also unit-independent when it incorporates the result of the Michelson-Morley experiment: light emitted by sources moving at different speeds travels at the same speed. Again it makes use of ratios of speeds: the ratio of the sources’ speeds, and the ratios of the different light rays’ speeds. However the postulate looses its meaning when it refers to light speed at different points, or even to the speed of light moving in different directions at a given point.
Hence Lorentz invariance in its usual definition is not a unit-independent concept, and indeed relativity is not a unit-independent construction. For instance relativity is not conformally invariant (a conformal transformation being just a particular type of unit transformation). A unit-independent definition of Lorentz invariance may be inferred by taking a Lorentz invariant theory and subjecting it to the most general unit transformation. The resulting theory retains the unit-invariant aspects of the second postulate, and clearly $`c`$ may now be anisotropic and vary in space-time. Under such circumstances what is the structure which represents Lorentz invariance?
For simplicity we specialize to changes of units which only affect the local value of the modulus of $`c`$. We redefine units of time and space in all inertial systems
$`d\widehat{t}`$ $`=`$ $`dtϵ^\alpha `$ (1)
$`d\widehat{x}^i`$ $`=`$ $`dx^iϵ^\beta `$ (2)
where $`ϵ`$ can be any function, and the metric (in this case the Minkowski metric) is left unchanged. If $`\alpha =\beta `$ we have an active conformal transformation (for a passive conformal transformation $`dx`$ and $`dt`$ are left unchanged, and the metric is multiplied by $`ϵ^{2\alpha }`$). If $`\alpha \beta `$, a Lorentz invariant theory is replaced by a theory in which $`c`$ remains isotropic, color independent, and independent of the speeds of observer and emitter; but it varies like $`\widehat{c}ϵ^{\beta \alpha }`$. A general unit transformation may be decomposed into a conformal transformation plus a VSL transformation with $`\beta =0`$.
It is immediately obvious that local Lorentz transformations in the new units are preserved:
$`d\widehat{t}^{}`$ $`=`$ $`\gamma \left(d\widehat{t}{\displaystyle \frac{\widehat{v}}{\widehat{c}^2}}d\widehat{x}\right)`$ (3)
$`d\widehat{x}^{}`$ $`=`$ $`\gamma \left(d\widehat{x}\widehat{v}d\widehat{t}\right)`$ (4)
with
$$\gamma =\frac{1}{\sqrt{1\left(\frac{\widehat{v}}{\widehat{c}}\right)^2}}$$
(5)
Hence the standard definition remains unmodified, if one employs the local value of $`c`$ in the transformation.
A novelty arises because changes of space-time units do not generally produce new coodinate patches because (1) needs not be holonomic: one may have eg. $`d^2\widehat{t}0`$. Hence there would not be a global $`\widehat{t}`$ time coordinate: the new “coordinate elements” would not be differentials of any coordinates. Even if in one frame the transformation (1) were holonomic, in a boosted frame it would not be. Some oddities pertaining to the new units follow. Partial derivatives generally do not commute. The change in the “coordinate time” between two points may depend upon the path taken to link the two points.
We have thus identified the structure of a VSL Lorentz invariant theory. The theory is locally Lorentz invariant in the usual way, using in local transformations the value of $`c`$ at that point. However local measurements of time and space are not closed forms, and therefore cannot be made into coordinates. Integrating factors can always be found, so that $`d\widehat{t}/ϵ^\alpha `$ and $`d\widehat{x}/ϵ^\beta `$ are closed forms, and $`\widehat{c}=ϵ^{\beta \alpha }`$.
Although a time coordinate does not generally exist, in many important cases it may be defined. If $`_\mu c^\mu c<0`$ then local coordinates exist so that $`c`$ only changes in time. We shall call this the homogeneous frame. Then $`d^2\widehat{t}=0`$, and a $`\widehat{t}`$ coordinate can be defined. Hence if we insist upon using a time coordinate we necessarily pick up a preferred reference frame - thereby violating the principle of relativity. This situation will be true in cosmology (where the preferred frame is the cosmological frame) but not in the context of static solutions, such as black hole solutions. Also a time coordinate may always be defined along a line. In particular for a geodesic, the amount of proper time is always well defined, although the proper time between two points depends on the trajectory (a situation already true in general relativity).
## III Generalized covariance
In order to construct a theory of gravitation we need to discuss general covariance. Covariance is the requirement of invariance under the choice of coordinate chart. This may be trivially adapted to VSL if we only use charts employing an “$`x^0`$” coordinate, with dimensions of length rather than time. Then $`c`$ appears nowhere in the usual definitions of differential geometry, which may therefore still be used. The laws for the transformation of tensors are the same as usual. The metric is dimensionless in all components and does not explicitly depend on $`c`$; it transforms like a rank 2 tensor. The usual Cristoffell connection may be defined from the metric by means of the standard formula, without any extra terms in the gradients of $`c`$ (which only appear if we try to revert to a time type of coordinate). A curvature tensor may still be defined in the usual way, and a Ricci tensor and scalar derived from it. The volume measure does not contain $`c`$. As we will see, many novelties introduced by a varying $`c`$ only emerge when we try to connect the $`x^0`$ coordinate with time.
Whenever applying a unit transformation (1) to a covariant theory, the above remarks apply only to the VSL part of the transformation (that is the component with $`\beta =0`$). For the conformal part of the transformation, with $`ϵ^\alpha =ϵ^\beta =\mathrm{\Omega }`$, the structures of differential geometry transform in the usual way . For instance the Ricci scalar transforms as:
$$\widehat{R}=\frac{R}{\mathrm{\Omega }^2}6\frac{\mathrm{}\mathrm{\Omega }}{\mathrm{\Omega }^3}$$
(6)
for active conformal transformations.
It is not altogether surprising that covariance may be redefined so easily for a theory with such different foundations. It has been pointed out that covariance is an empty requirement (see ). Not only does covariance not imply local Lorentz invariance, but also any theory can be made covariant. An example of a covariant formulation of Newtonian gravity is given in . In this theory the tangent space is not a portion of Minkowski space, rather a portion of Galilean space.
## IV An overview of the underlying structure
What structure represents covariance and local Lorentz invariance when $`c`$ is allowed to vary? We found that it is a unit-invariant redefinition of these concepts, which indeed does not differ much from the usual definitions if we phrase them suitably. Local Lorentz transformations are the same as usual, using the local value of $`c`$. Covariance and the usual constructions of differential geometry remain unchanged as long as a $`x^0`$ coordinate is used, or more generally if all coordinates used have the same dimensions.
What is new, then? The novelty is that locally made time and space measurements produce a set of infinitesimals which are generally not closed forms. Therefore space-time measurements cannot be made into local coordinate patches. This leads to the following modification of the structure of relativity. The underlying structure of general relativity is a manifold, combined with its tangent bundle (where physics actually happens). If $`c`$ varies the underlying structure is a fibre bundle. The base manifold has the same structure as usual, but the fibres in which local measurements happen are not the tangent bundle. The fibers are vector spaces obtained by means of a non-holonomic transformation over the tangent bundle.
It may seem rather contorted to adopt the above structure when we know that a unit transformation would transform it into standard covariance and local Lorentz invariance. However such a structure has the merit that it only incorporates those elements of the original structure which are unit-independent, and can therefore be the outcome of experiment. Moreover such structure allows for a varying speed of light within a covariant framework, which is precluded by the standard framework. What one may gain from such extra freedom is a simplified description of any given physical situation, when all fine structure coupling constants are allowed to vary, in what looks like a contrived fashion, if we use units such that $`c`$ is constant.
We wish to propose a theory which permits space-time variations in all coupling constants; more specifically in generalized fine structure constants $`\alpha _i=g_i^2/(\mathrm{}c)`$ \- where $`g_i`$ are the various charges corresponding to all interactions apart from gravitation. This purpose draws inspiration from the findings of . However we restrict such variations so that the ratios between the $`\alpha _i`$ remain constant. This suggests that attributing the variations in the $`\alpha _i`$ to changes in $`c`$ or $`\mathrm{}`$ might lead to a simpler picture. In suitable units we could regard our theory as a “generalized ” Bekenstein changing $`e`$ theory, but in this system of units the picture is rather contrived.
We will see, in Section VIII, that a natural dynamics will emerge in this theory which becomes unnecessarily complicated when the theory is reformulated in fixed $`c`$ units. Whatever the system of units chosen the general theory we will consider is not a dilaton theory. Some important geometrical aspects (such as inaccessible regions of space-time to be studied in Section IX) are missed altogether in the fixed $`c`$ system.
## V Matter fields subjected to VSL
Before embarking on an investigation of the dynamics of $`c`$ and of gravitation, we first undertake a careful examination of the effects of a varying $`c`$ upon the matter fields. The key point here is that it is always possible to write the matter Lagrangian so that is does not depend explicitly on $`c`$. We may break this rule, if we wish to, but this is not necessary. This remark is highly non-trivial, and relies heavily on using an $`x^0`$ coordinate (as opposed to time). The introduction of a time coordinate would not only introduce non-covariant elements in expressions like $`_\mu \varphi =(_t\varphi /c,_i\varphi )`$, but would also force the matter Lagrangian to depend explicitly on $`c`$, via kinetic terms.
Using an $`x^0`$ coordinate the situation is rather different. For instance, for a massless scalar field with no interactions we have:
$$_m=\frac{1}{2}(_\mu \varphi )(^\mu \varphi )$$
(7)
which does not depend on $`c`$. Similarly for a spin 1/2 free massless field we have
$$_m=i\overline{\chi }\gamma ^\mu _\mu \chi $$
(8)
The above expressions, in particular the latter, are sometimes multiplied by $`\mathrm{}c`$ (see for example Mandl and Shaw ). If $`c`$ and $`\mathrm{}`$ are constant this operation has no effects, other than modifying the dimensions of the fields. However in a minimal VSL theory such an operation should be banned. All dynamical fields should be defined with dimensions such that the kinetic terms have no explicit dependence on either $`c`$ or $`\mathrm{}`$. This forces all matter fields to have dimensions of $`\sqrt{E/L}`$.
The only chance for $`_m`$ to depend upon $`c`$ therefore comes from mass and interaction terms. These may always be defined so that no explicit dependence on $`c`$ is present. By dimensional analysis this requirement fully defines how masses, charges, and coupling constants scale with $`c`$, provided we know how $`\mathrm{}`$ scales with $`c`$. This issue will be discussed further in the next Section.
Let us first consider mass terms. For a scalar field we have:
$$_m=\frac{1}{2}\left(_\mu \varphi ^\mu \varphi +\frac{1}{\lambda _\varphi ^2}\varphi ^2\right)$$
(9)
Hence in a minimal VSL theory the Compton wavelength $`\lambda _\varphi `$ of the particle should not depend on $`c`$. For a massive spin 1/2 particle we have
$$_m=i\overline{\chi }\gamma ^\mu _\mu \chi \frac{1}{\lambda _\chi }\overline{\chi }\chi $$
(10)
with a similar requirement. More generally we find that $`c`$ does not appear in mass terms if all particles’ masses are proportional to $`\mathrm{}/c`$ (or their rest energies proportional to $`\mathrm{}c`$).
If we now consider fields coupled to electromagnetism we find that the electric charge $`e`$ should scale like $`\mathrm{}c`$, if explicit dependence on $`c`$ is to be avoided. Consider for instance a $`U(1)`$ gauged complex scalar field. Its action may be written as
$$_m=(D^\mu \varphi )^{}D_\mu \varphi \frac{|\varphi |^2}{\lambda _\varphi ^2}\frac{1}{4}F_{\mu \nu }F^{\mu \nu }$$
(11)
where the $`U(1)`$ covariant derivative is
$$D_\mu =_\mu +i\frac{e}{\mathrm{}c}A_\mu $$
(12)
and the electromagnetic tensor is
$$F_{\mu \nu }=_\mu A_\nu _\nu A_\mu $$
(13)
Hence $`e`$ should be proportional to $`\mathrm{}c`$. The same holds true for fields of any spin coupled to electromagnetism, since $`e`$ only appears in the definition of the covariant derivative. Notice that the constancy of $`e/(\mathrm{}c)`$ is also required for the gauge-invariant field strength tensor not to receive any corrections. Gauge transformations should take the form:
$$\delta A_\mu =\frac{\mathrm{}c}{e}_\mu f$$
(14)
for $`\delta \varphi =if\varphi `$, where $`f`$ is any function. This is necessary so that $`D_\mu \varphi `$ transforms covariantly: $`\delta (D_\mu \varphi )=ifD_\mu \varphi `$. But then the gauge invariant field strength tensor must be defined as
$$F_{\mu \nu }=\frac{\mathrm{}c}{e}\left(_\mu \left(\frac{e}{\mathrm{}c}A_\nu \right)_\nu \left(\frac{e}{\mathrm{}c}A_\mu \right)\right)$$
(15)
and indeed this receives extra terms if $`e/(\mathrm{}c)`$ is not constant.
Inspection of the electroweak and strong interaction Lagrangians reveals that their coupling charges $`g`$ should also scale like $`\mathrm{}c`$. This is indeed a general feature for any interaction, and follows from dimensional analysis. It is always the combination $`g/(\mathrm{}c)`$ that appears in covariant derivatives and, in non-Abelian theories, in the gauge field strength tensor.
Next we discuss two important cases to be used later in this paper: a field undergoing spontaneous symmetry breaking, and a matter cosmological constant. Consider a U(1) gauge symmetric complex scalar field as above, but with a potential
$$V(\varphi )=\frac{1}{\lambda _\varphi ^2}|\varphi |^2\frac{1}{2\lambda _\varphi ^2\varphi _0^2}|\varphi |^4$$
(16)
Then the Compton wavelength $`\lambda _\varphi `$ and $`\varphi _0`$ should both be independent of $`c`$. If the quartic term is ignored then the vacuum is at $`\varphi =0`$, so that we have a massive complex scalar field (with Compton wavelength $`\lambda _\varphi `$), and a massless gauge boson. The charge is $`e\mathrm{}c`$. If we consider the quartic term, as is well known, we have spontaneous symmetry breaking. The vacuum is now at $`|\varphi |=\varphi _0`$. Expanding around the vacuum we find a real scalar field with Compton wavelength $`\lambda _\varphi `$, and a massive gauge boson with Compton wavelength
$$\frac{1}{\lambda _A}=\frac{e}{\mathrm{}c}\varphi _0$$
(17)
which therefore is independent of $`c`$. Hence the rest energies of all massive particles, regardless of the origin of their mass, scale like $`\mathrm{}c`$. Due to spontaneous symmetry breaking the vacuum energy decreases by
$$\mathrm{\Delta }V=\frac{\varphi _0^2}{2\lambda ^2}$$
(18)
and so this process gives rise to a negative vacuum energy, if the original vacuum energy is zero. We shall label it by $`\mathrm{\Lambda }_m=\mathrm{\Delta }V`$, and call it the matter cosmological constant. It adds a term to the matter Lagrangian
$$_m=\mathrm{\Lambda }_m$$
(19)
Under minimal coupling $`\mathrm{\Lambda }_m`$ does not depend on $`c`$. However we could also allow $`\varphi _0`$, and therefore $`\mathrm{\Lambda }_m`$, to depend on $`c`$ (as we shall do in ).
Finally we consider an example of a classical Lagrangian, that of a charged particle in a field:
$`_m(x^\gamma )=`$ (20)
$`{\displaystyle 𝑑\lambda \left[\frac{E_0}{2}\frac{dy^\mu }{d\lambda }\frac{dy_\mu }{d\lambda }+eA_\mu \frac{dy^\mu }{d\lambda }\right]\frac{\delta ^{(4)}(x^\gamma y^\gamma )}{\sqrt{g}}}`$ (21)
Here the affine parameter is $`d\lambda =cd\tau `$, where $`\tau `$ is proper time. Note that any of the variations sometimes employed in the literature, eg. using the square root of $`u^2`$ (with $`u=dx/d\lambda `$), should not be used. This is because, as we shall see, $`u^2`$ needs not remain constant. Minimal coupling therefore requires that the particle’s rest energy ($`E_0=m_0c^2`$) and charge $`e`$ be independent of $`c`$.
In non-minimal theories we may consider a direct dependence on $`c`$ in the matter Lagrangian. This is far from new: for instance Bekenstein’s theory allows for a direct coupling between a varying $`e`$ and all forms of matter coupled to electromagnetism.
### A A worked out example
Consider a massive scalar field $`\varphi `$ in flat space-time (metric $`\eta _{\mu \nu }=\mathrm{diag}(1,1,1,1)`$ if we use a $`x^0`$ coordinate) with a variation in $`c`$ such that
$$c=\frac{c_0}{1+\frac{c_0t}{R}}$$
(22)
in suitable coordinates (so that $`c`$ does not vary in space). We have defined $`c_0`$ as the speed of light at time $`t=0`$. At time $`t=R/c_0`$ the speed of light goes to infinity. As time progresses the speed of light decays to zero, as $`t\mathrm{}`$. We shall see that this is indeed the solution corresponding to flat space-time. Then $`\varphi `$ satisfies:
$$\ddot{\varphi }^2\varphi +\frac{1}{\lambda _\varphi ^2}\varphi =0$$
(23)
which may be solved with Fourier series, with amplitudes subject to
$$\ddot{\varphi }_k+\left(k^2+\frac{1}{\lambda _\varphi ^2}\right)\varphi =0$$
(24)
The solution is
$$\varphi _k=\varphi _0(k)e^{i(\pm k^0x^0+𝐤𝐱)}$$
(25)
with a dispersion relation
$$(k^0)^2=k^2+\frac{1}{\lambda _\varphi ^2}$$
(26)
As expected there is nothing new if we use a $`x^0`$ coordinate.
If we insist on using a time coordinate we find that we can only do so in one inertial frame, the one in which the speed of light is homogeneous. By requesting to use a time coordinate, and make contact with physics, we therefore select a preferred reference frame. In this frame:
$$x^0=c𝑑t=R\mathrm{log}\left(1+\frac{c_0t}{R}\right)$$
(27)
and so we find that around a given time $`t=t_0`$ we have the Taylor expansion:
$$k^0x^0=k^0c(t_0)(tt_0)+k^0R\mathrm{log}\left(1+\frac{c_0t_0}{R}\right)$$
(28)
We find that the local frequency changes proportionally to $`c`$:
$$\omega (t)=k^0c(t)=\frac{k^0c_0}{1+\frac{c_0t}{R}}$$
(29)
In addition there is a phase shift with value
$$\mathrm{\Phi }_0=k^0R\mathrm{log}\left(1+\frac{c_0t_0}{R}\right)$$
(30)
As we approach the initial singularity the wave suffers infinite blueshift. As time flows it redshifts progressively. The similarities between this effect and the cosmological redshift have been pointed out in . However the effect presented here is not due to gravity (expansion) but is due purely to the varying speed of light.
Naturally the above identification of a local frequency is only valid if $`\omega |\dot{c}/c|`$. This amounts to requiring:
$$k^0\frac{1}{R\left(1+\frac{c_0t}{R}\right)}$$
(31)
Hence any plane-wave approximation breaks down near the initial singularity; an interesting result.
## VI Quantization
Unfortunately the requirement that $`c`$ does not appear explicitly in $`_m`$ does not fix the scaling with $`c`$ of all “constants”: Planck’s and Boltzmann’s constants, $`\mathrm{}`$ and $`k_B`$, are left unfixed. Furthermore these two “constants” cannot be fixed by the classical dynamics, i.e. by adding to the action dynamical terms in two scalar fields $`\mathrm{}`$ and $`k_B`$ (as we shall do with $`c`$). Instead these have to be provided as a function of $`c`$, by means of scaling laws $`\mathrm{}(c)`$ and $`k_B(c)`$. These scaling laws should be regarded as postulates of the theory.
Here we explore the implications of various $`\mathrm{}(c)`$ laws. Let us consider first the simple case of a non-relativistic linear harmonic oscillator. Its Lagrangian is given by:
$$L=\frac{1}{2}m\dot{x}^2\frac{1}{2}m\omega ^2x^2$$
(32)
and we assume that $`m`$ and $`\omega `$ are independent of $`c`$, and therefore constant. Its Hamiltonian is
$$H=\frac{p^2}{2m}+\frac{m\omega ^2x^2}{2}$$
(33)
and is time-independent. We postulate that quantization produces an expression of the form
$$\widehat{H}=\mathrm{}\omega (\widehat{N}+1/2)$$
(34)
where $`\widehat{N}`$ is the particle number operator, ie: the classical energy of the oscillator is in quanta of energy $`\mathrm{}\omega `$. Hence
$$\frac{d}{dt}\mathrm{}(\widehat{N}+1/2)=0$$
(35)
This implies that should $`\mathrm{}`$ drop, the particle number would increase, which is hardly surprising. Indeed the amplitude $`A`$ and frequency $`\omega `$ of the classical oscillations remain constant, and therefore so does their total energy $`E=m\omega ^2A^2/2`$. However the quantum particles contained in the oscillator have energies $`\mathrm{}\omega `$ which vary like $`\mathrm{}`$. To reconcile these two facts the number of particles has to vary, proportionally to $`1/\mathrm{}`$. Such a phenomenon has a clear experimental meaning, since the number of particles does not depend on the units being used.
Furthermore if the oscillator is initially in the vacuum state, a drop in $`\mathrm{}`$ suppresses the zero-point energy. Particles should therefore be produced so that $`\widehat{H}`$ remains constant. We have both particle multiplication and particle production (a phenomenon noticed before in VSL theories by ).
Since creation and anihilation operators satisfy a time-independent algebra $`[a,a^{}]=1`$, the best way to express the variability of $`\widehat{N}`$ is by means of a Bogolubov-type transformation. A short calculation shows that:
$$a(t^{})=\alpha a(t)+\beta ^{}a^{}(t)$$
(36)
with
$`|\alpha |^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{}(t)+\mathrm{}(t^{})}{\mathrm{}(t^{})}}`$ (37)
$`|\beta |^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{}(t)\mathrm{}(t^{})}{\mathrm{}(t^{})}}`$ (38)
enforces that the expectation values of $`\widehat{N}(t)=a^{}(t)a(t)`$ satisfy (35).
This discussion generalizes to relativistic quantum field theory, with $`\mathrm{}c`$ replacing $`\mathrm{}`$. Now we should have:
$$\widehat{H}=\underset{𝐤}{}\mathrm{}\omega (\widehat{N}+1/2)$$
(39)
with $`\omega =k^0cc`$. Hence now
$$\frac{d}{dx^0}\mathrm{}c(\widehat{N}+1/2)=0$$
(40)
The time dependence in $`\widehat{N}`$ can now be expressed in the form of a Bogolubov transformation
$$a(k^0,x^{\overline{\mu }})=\alpha a(k^0,x^\mu )+\beta ^{}a^{}(k^0,x^\mu )$$
(41)
with
$`|\alpha |^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{}(x^\mu )c(x^\mu )+\mathrm{}(x^{\overline{\mu }})c(x^{\overline{\mu }})}{\mathrm{}(x^\mu )c(x^\mu )}}`$ (42)
$`|\beta |^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{}(x^\mu )c(x^\mu )\mathrm{}(x^{\overline{\mu }})c(x^{\overline{\mu }})}{\mathrm{}(x^\mu )c(x^\mu )}}`$ (43)
Recalling that $`g_i/(\mathrm{}c)`$ is a constant, we have particle production at a rate proportional to $`1/\alpha _i`$, where $`i`$ labels the various interactions. We shall parameterize $`\mathrm{}(c)`$ by means of an exponent $`q`$ such that
$$\alpha _ig_i\mathrm{}cc^q$$
(44)
## VII Gravitational dynamics
We now set up some possibilities for actions governing the evolution of the metric and speed of light. Only a small class of these actions may be transformed into a dilaton action, by means of a unit transformation. In Appendix we describe a somewhat orthogonal approach.
We shall take as our starting point the action of General Relativity:
$$S=d^4x\sqrt{g}\left(R2\mathrm{\Lambda }+\frac{16\pi G}{c_0^4}_m\right)$$
(45)
where $`R`$ is the Ricci scalar, and $`\mathrm{\Lambda }`$ is the geometrical cosmological constant (as defined in ) and $`_m`$ is the Lagrangian of all the matter fields (including the above mentioned matter cosmological constant).
A changing $`G`$ theory was proposed by Brans and Dicke , and we shall work in analogy to this generalization of General Relativity in what follows, albeit with a couple of crucial differences. The idea in this paper (in ) is to replace $`c`$ ($`G`$) by a field, wherever it appears in (45). In addition one should add a term to the Lagrangian describing the dynamics of $`c`$ ($`G`$). An ambiguity appears because (45) may be divided by any power of $`c`$ ($`G`$), before the replacement is performed. Brans and Dicke avoided commenting on this ambiguity, and cunningly performed the necessary division by $`G`$ which led to a theory with energy conservation. We shall not be hampered by this restriction; indeed we expect violations of energy conservation in VSL. Hence we consider actions in which the replacement is made after the most general division by $`c`$ is made. In the simplest case we define a scalar field
$$\psi =\mathrm{log}\left(\frac{c}{c_0}\right)$$
(46)
so that $`c=c_0e^\psi `$, and take
$$S=d^4x\sqrt{g}(e^{a\psi }(R2\mathrm{\Lambda }+_\psi )+\frac{16\pi G}{c_0^4}e^{b\psi }_m)$$
(47)
The simplest dynamics for $`\psi `$ derives from:
$$_\psi =\kappa (\psi )_\mu \psi ^\mu \psi $$
(48)
where $`\kappa (\psi )`$ is a dimensionless coupling function (to be taken as a constant in most of what follows). We shall impose $`ab=4`$, although this is not necessary.
Notice that $`a=4`$, $`b=0`$, is nothing but a unit transformation applied to Brans-Dicke theory, with
$`\varphi _{bd}`$ $`=`$ $`{\displaystyle \frac{e^{4\psi }}{G}}`$ (49)
$`\kappa (\psi )`$ $`=`$ $`16\omega _{bd}(\varphi _{bd})`$ (50)
This shall be proved in Section VIII, where we identify the full set of cases which are a mere unit transformation applied to existing theories. Among the theories which are truly new, $`a=0`$, $`b=4`$ is particularly simple and we shall call it minimal VSL.
We can trivially generalize this construction, by complicating the dynamics encoded in $`_\psi `$, for instance by adding a potential $`V(\psi )`$ to it. We can also take for $`\psi `$ a complex, vector, or spinor field, with the speed of light deriving from a scalar associated with $`\psi `$ (eg. $`\overline{\psi }\psi `$ for a spinor field). A nice example (developed further in Section X) is a theory in which $`\psi `$ is a complex field, with
$$c=c_0e^{|\psi |^2}$$
(51)
and with a Mexican hat potential added to $`_\psi `$.
Another important novelty of our theory, not included in Brans-Dicke theory (but noted by ), is that we allow $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }_m`$ to depend on $`c`$. It seems fair to allow $`\mathrm{\Lambda }`$, like $`\mathrm{}`$ or $`k_B`$, to depend on $`c`$. After all $`\mathrm{\Lambda }`$ is a much less fundamental constant. On the contrary if $`\mathrm{\Lambda }_m`$ depends on $`c`$, then so does the vacuum expectation value $`\varphi _0`$, and so we have gone beyond minimal matter coupling. In what follows we shall absorb $`\mathrm{\Lambda }_m`$ into a total geometrical Lambda
$$\overline{\mathrm{\Lambda }}=\mathrm{\Lambda }+\frac{8\pi G}{c^4}\mathrm{\Lambda }_m$$
(52)
In our applications to cosmology we shall assume that
$$\mathrm{\Lambda }(c/c_0)^n=e^{n\psi }$$
(53)
and
$$\mathrm{\Lambda }_m(c/c_0)^m=e^{m\psi }$$
(54)
We will see that allowing $`\mathrm{\Lambda }`$ to depend on $`c`$ leads to interesting cosmological scenarios . In such theories it is the presence of a Lambda problem that drives changes in the speed of light. These in turn solve the cosmological constant and other problems of Big Bang cosmology. In effect Lambda acts as a potential driving $`\psi `$.
### A Gravitational field equations
The field equations in this theory may now be derived by varying the action. Variation with respect to the metric leads to gravitational equations
$`G_{\mu \nu }+\mathrm{\Lambda }g_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{8\pi G}{c^4}}T_{\mu \nu }+\kappa \left(_\mu \psi _\nu \psi {\displaystyle \frac{1}{2}}g_{\mu \nu }_\delta \psi ^\delta \psi \right)`$ (56)
$`+e^{a\psi }(_\mu _\nu e^{a\psi }g_{\mu \nu }\mathrm{}e^{a\psi })`$
where the matter stress energy tensor is defined as usual:
$$T_{\mu \nu }=\frac{2}{\sqrt{g}}\frac{\delta S_m}{\delta g^{\mu \nu }}$$
(57)
These equations are particularly simple for minimal VSL ($`a=0`$ and $`b=4`$).
Variation with respect to $`\psi `$ leads to
$`\mathrm{}\psi `$ $`+a_\mu \psi ^\mu \psi `$ (59)
$`={\displaystyle \frac{8\pi G}{c^4(2\kappa +3a^2)}}(aT2a\rho _\mathrm{\Lambda }2b_m)+{\displaystyle \frac{1}{\kappa }}{\displaystyle \frac{d\overline{\mathrm{\Lambda }}}{d\psi }}`$
Again minimal VSL is particularly simple:
$$\mathrm{}\psi =\frac{32\pi G}{c^4\kappa }\left(_m\left(1\frac{m}{4}\right)\mathrm{\Lambda }_m\right)+\frac{1}{\kappa }n\mathrm{\Lambda }$$
(60)
As announced above, in general either a matter or a geometrical Lambda drive changes in $`c`$. The total matter Lagrangian $`_m`$ also drives changes in $`c`$, if $`b0`$. Ambiguities in writing $`_m`$ (total divergences) are therefore relevant for $`c`$, as indeed for the matter field equations under VSL (see below).
### B Impact upon matter field equations
Bianchi identities applied to (56) and (59) imply
$$_\mu (T_\nu ^\mu e^{b\psi })=be^{b\psi }_m_\nu \psi $$
(61)
or equivalently:
$$_\mu T_\nu ^\mu =b(T_\mu ^\nu \delta _\mu ^\nu _m)_\nu \psi $$
(62)
Therefore we only have energy conservation if $`a=4,`$ $`b=0`$. In all other cases a varying $`c`$ creates or destroys energy; indeed beyond the naive expectation (the term in $`_m`$ is far from expected). This fact merely reflects the interaction between the matter fields and the gravitational field $`\psi `$, present due to the coupling $`e^{b\psi }_m`$. This interaction affects the field equations for matter, beyond what was descrobed in Section V (which is only strictly correct if $`b=0`$). Indeed taking the variation with respect to matter fields, in every situation where it is usual to neglect a full divergence, a new term in $`^\mu \psi `$ now appears. For instance scalar fields satisfy a modified Klein-Gordon equation:
$$\left(\mathrm{}\frac{1}{\lambda _\varphi ^2}\right)\varphi =b_\mu \varphi ^\mu \psi $$
(63)
with gradients of $`\psi `$ driving the field $`\varphi `$ and therefore changing its energy balance. All field equations will be similarly affected, with a net result that energy conservation is violated according to (61).
To give a concrete example, the plane wave solution studied in Section V A is now subject to:
$$\ddot{\varphi }_k+\left(k^2+\frac{1}{\lambda _\varphi ^2}\right)\varphi =b\frac{\dot{\varphi }_k}{R}$$
(64)
A solution is
$$\varphi _k=[\varphi _0(k)e^{\frac{bx^0}{R}}]e^{i(\pm k^0x^0+𝐤𝐱)}$$
(65)
subject to the same dispersion relation. Hence, in addition to the effects studied in Section V A, the amplitude of the plane waves is now proportional to $`c^b`$. If $`R>0`$, and $`b>0`$, we not only have a “redshift effect” (affecting the energy of the field quanta), but the classical energy of the field also dissipates.
Finally note that we may also take on board terms which are usually neglected in minimal theories because they are full divergences. If $`b0`$ these terms affect the matter field equations; indeed they drive changes in $`c`$. For instance one could consider electromagnetism based on
$$_m=\frac{1}{4}(F_{\mu \nu }F^{\mu \nu }+\zeta F_{\mu \nu }\stackrel{~}{F}^{\mu \nu })$$
(66)
where $`\stackrel{~}{F}`$ is the dual of $`F`$ and $`\zeta `$ is a constant. The second term is usually irrelevant, because it is a full divergence. However we now have Maxwell’s equations:
$$_\mu F^{\mu \nu }+b(F^{\mu \nu }+\zeta \stackrel{~}{F}^{\mu \nu })_\mu \psi =j^\nu $$
(67)
where $`j^\nu `$ is the electric current.
### C Effect upon classical particles
These processes are also reflected in the equations of motion for a point particle. From (20), with $`e=0`$, we can derive the stress energy tensor:
$$T^{\mu \nu }(x^\delta )=mc^2𝑑\lambda \frac{dy^\mu }{d\lambda }\frac{dy^\nu }{d\lambda }\frac{\delta ^{(4)}(x^\delta y^\delta (\lambda ))}{\sqrt{g}}$$
(68)
where we have assumed that $`mc^2`$ is a constant (so that the matter Lagrangian does not depend on $`c`$). From (61) one gets:
$$\frac{d^2x^\mu }{d\lambda ^2}+\mathrm{\Gamma }_{\nu \delta }^\mu \frac{dx^\nu }{d\lambda }\frac{dx^\delta }{d\lambda }=b\left(\frac{dx^\mu }{d\lambda }\frac{dx^\nu }{d\lambda }\frac{1}{2}\frac{dx^\alpha }{d\lambda }\frac{dx_\alpha }{d\lambda }g^{\mu \nu }\right)\psi _{,\nu }$$
(69)
where we recall $`d\lambda =cd\tau `$. Alternatively we may integrate the volume integral in (20) to obtain action:
$$S=\frac{E_0}{2}𝑑\lambda e^{b\psi }g_{\mu \nu }\dot{x}^\mu \dot{x}^\nu $$
(70)
Direct variation of this action is equivalent to equation (69) and may be more practical. An immediate first integral of this action is:
$$u^2=u_0^2(c/c_0)^b$$
(71)
with $`u^\mu =dx^\mu /d\lambda `$. Hence null particles remain null, but time-like lines have a variable $`u^2`$.
We see that matter no longer follows geodesics. However all bodies with the same set of initial conditions fall in the same way. A weak form of the equivalence principle is therefore satisfied. In particular there is no conflict between these theories and the Eotvos experiment. In we shall investigate the impact of these effects upon the standard tests of gravitational light deflection, and the perihelium of Mercury. Here, however, we limit ourselves to integrating the geodesic equation in the local free-falling frame, or in flat space-time. Then (70) produces the Lagrangian
$$=e^{b\psi }(\dot{(x^0)}^2+\dot{x}^2)$$
(72)
where dots represent $`d/d\lambda `$. There are three conserved quantities: $`E=\dot{x}^0e^{b\psi }`$, $`p=\dot{x}e^{b\psi }`$, and $`=1`$, from which we may conclude
$$\frac{v^2}{c^2}=\frac{p^2}{E^2}=1\frac{e^{b\psi }}{E^2}$$
(73)
As a result the particle’s gamma factor
$$\gamma ^2=\frac{1}{1v^2/c^2}c^b$$
(74)
If $`b0`$ the field $`\psi `$ will therefore accelerate or brake particles.
## VIII Fixed speed of light duals
We now identify which of our theories are simply well-known fixed $`c`$ theories subject to a change of units. By doing so we will also expose the undesirable complication of the fixed $`c`$ picture in all other cases.
Let us first rewrite our theories in units in which $`c`$, $`\mathrm{}`$, and $`G`$ are fixed, but the couplings $`g`$ are variable, thereby mapping VSL theories into “Bekenstein” changing charge theories. Recalling that in VSL units we have $`\alpha _ig_i\mathrm{}cc^q`$ (cf. Eqn. (44), we should perform the following change of units:
$`d\widehat{t}`$ $`=`$ $`dte^{(3\frac{q}{2})\psi }`$ (75)
$`d\widehat{x}`$ $`=`$ $`dxe^{(2\frac{q}{2})\psi }`$ (76)
$`d\widehat{E}`$ $`=`$ $`dEe^{(2\frac{q}{2})\psi }`$ (77)
In the new units $`\widehat{c}`$, $`\widehat{\mathrm{}}`$, and $`\widehat{G}`$ are constant, $`\widehat{g}e^{\frac{q}{2}\psi }`$, and indeed $`\widehat{\alpha }=\alpha c^q`$. Subjecting a VSL minimally coupled matter action to this transformation leads to an action very far from minimal coupling. Indeed all matter fields, eg $`\varphi `$, transform like
$$\widehat{\varphi }=\varphi e^{2\psi }$$
(78)
Hence all kinetic terms become rather contorted, since in the new units
$$_\mu \varphi _{\widehat{\mu }}\widehat{\varphi }+2\widehat{\varphi }_{\widehat{\mu }}\psi $$
(79)
This leads to complex additions to mass and interaction terms. For gauged fields we have
$$D_\mu \varphi D_{\widehat{\mu }}\widehat{\varphi }+2\widehat{\varphi }_{\widehat{\mu }}\psi $$
(80)
leading to similar complications, and to breaking of standard gauge invariance. In conclusion we can transform a VSL theory which is minimally coupled to matter (up to the $`b0`$ factor) into a fixed $`c`$, $`\mathrm{}`$ and $`G`$ theory. However the result is a rather unnatural construction, quite distinct from the changing charge theories previously discussed. None of our theories is a standard changing charge theory in disguise; indeed choosing a standard changing $`g_i`$ picture for them is undesirable.
The above may be avoided if we map our VSL theories into theories in which $`c`$ and $`\mathrm{}`$ are constants, but $`G`$ may vary. Then in order to preserve minimal coupling all matter fields should remain unaffected by the unit transformation, eg $`\widehat{\varphi }=\varphi `$ This requires
$`d\widehat{t}`$ $`=`$ $`dte^{(1\frac{q}{2})\psi }`$ (81)
$`d\widehat{x}`$ $`=`$ $`dxe^{\frac{q}{2}\psi }`$ (82)
$`d\widehat{E}`$ $`=`$ $`dEe^{\frac{q}{2}\psi }`$ (83)
and so we have that
$$\frac{\widehat{G}}{G}=e^{4\psi }$$
(84)
While this ensures minimal coupling for all quantum fields, it does not do the job for classical point particles, if $`q0`$. With the above change of units one should make the identification
$$\widehat{\varphi }_{bd}=\frac{1}{\widehat{G}}=e^{4\psi }$$
(85)
and write down the transformed action:
$`\widehat{S}`$ $`=`$ $`{\displaystyle }d^4\widehat{x}\sqrt{\widehat{g}}(\widehat{\varphi }_{bd}^{\frac{a+q}{4}}(\widehat{R}2\widehat{\mathrm{\Lambda }})`$ (87)
$`{\displaystyle \frac{f(\widehat{\varphi }_{bd})}{\widehat{\varphi }_{bd}}}_\mu \widehat{\varphi }_{bd}^\mu \widehat{\varphi }_{bd}+{\displaystyle \frac{16\pi G_0}{c_0^4}}\widehat{\varphi }_{bd}^{\frac{b+q}{4}}\widehat{}_m)`$
We see that only theories for which $`b+q=0`$ are scalar-tensor theories in disguise. If $`q=0`$ all structure “constants” $`\alpha _i`$ are constant, and indeed for $`b=0`$ (and so $`a=4`$) we can recognize in $`\widehat{S}`$ the Brans-Dicke action. However we see that there are also changing $`\alpha `$ theories which are really Brans Dicke theories in unusual units: theories with $`b=q0`$. Such theories are dilaton theories. On the contrary, if $`b+q0`$ we have theories which can never be mapped into dilaton theories.
In addition one may perform conformal transformations upon VSL theories, mapping them into other VSL theories with different $`a`$ and $`b`$. The relevant formulae shall be given in . By means of conformal transformations it is always possible to write action (47) as a scalar-tensor theory, if the matter Lagrangian is homogeneous in the metric. The latter, however, is clearly not true, carrying with it the crucial implication that there is only one frame in which the coupling to matter is of the form $`e^{b\psi }_m`$, with $`_m`$ independent of $`\psi `$. Thus, the much heralded equivalence between conformal frames is broken as soon as matter is added to gravity and $`\psi `$ (a point clearly made by ). One may recognize $`a=4`$, $`b=0`$ as the Jordan’s frame, $`a=0`$, $`b=4`$ as the Einstein’s frame, and $`a=b=1`$ as the tree-level string frame. We should also note that the general class of couplings we have considered is contained within the theories proposed by Damour and Polyakov as representing low-energy limits to string theory, beyond tree-level. More specifically, using the notation of , our theories are those for which $`B_i(\mathrm{\Phi })`$ is the same for all the matter fields.
In spite of these comments, in we shall make use of conformal transformations to isolate the geodesic frame: the frame in which free-falling charge-free particles follow geodesics. This can always be defined because the Lagrangian of these particles is indeed homogeneous in the metric. While this trick simplifies some calculations, one should always bear in mind that the geodesic frame only looks simpler because a lot of gargabe is swept under the carpet by not writing the Lagrangian for all the other matter fields. Minimal coupling to all forms of matter always picks up a preferred frame, which is not the geodesic frame unless $`b=0`$.
## IX Empty space-time
The analogue of Minkowski space-time may be derived by setting $`T_\nu ^\mu =\mathrm{\Lambda }=0`$ in Equation (56). We should also set $`a=0`$ and $`\kappa =0`$ so as to switch off the gravitational effects of $`\psi `$. Then $`g_{\mu \nu }=\eta _{\mu \nu }=\mathrm{diag}(1,1,1,1)`$, using an $`x^0`$ coordinate.
The speed of light can be found from (59), which in coordinates in which $`\psi `$ is homogeneous becomes $`\ddot{\psi }=0`$. This leads to $`\dot{\psi }=\frac{1}{R}`$, where $`R`$ is an integration constant with dimensions of length ($`R`$ can be positive or negative). If we only use coordinates in which $`\psi `$ is homogeneous (or, as we shall see, if we stay close to the origin compared to the distance $`R`$) then a global time coordinate $`t`$ may be defined. In terms of it we have:
$$\frac{1}{c^2}\frac{dc}{dt}=\frac{1}{R}$$
(88)
which integrates to
$$c=\frac{c_0}{1+\frac{c_0t}{R}}$$
(89)
This is nothing but $`c`$ near the origin in Fock-Lorentz space-time, in which
$$c(𝐫,t;𝐧)=\frac{c_0}{1+\frac{c_0t}{R}}\left(𝐧+\frac{𝐫}{R}\right)$$
(90)
Even though global coordinates cannot be generally defined if $`c`$ varies, we find that this case is special. Relations
$`d\widehat{t}`$ $`=`$ $`{\displaystyle \frac{dt}{\left(1+\frac{c_0t}{R}\right)^{N+1}}}`$ (91)
$`d\widehat{r}`$ $`=`$ $`{\displaystyle \frac{d𝐫}{\left(1+\frac{c_0t}{R}\right)^N}}`$ (92)
(corresponding to $`\alpha =N+1`$ and $`\beta =N`$) may be recovered from
$`\widehat{t}`$ $`=`$ $`{\displaystyle \frac{R}{Nc_0}}\left(1{\displaystyle \frac{1}{\left(1+\frac{c_0t}{R}\right)^N}}\right)`$ (93)
$`\widehat{𝐫}`$ $`=`$ $`{\displaystyle \frac{𝐫}{\left(1+\frac{c_0t}{R}\right)^N}}`$ (94)
(with $`N>0`$) as long as $`|𝐫||R|`$. Hence, near the origin there are global varying $`c`$ coordinates $`t`$ and $`𝐫`$. Global transformation laws between inertial frames may be derived for these coordinates by writing global Lorentz transformations for $`\widehat{t}`$ and $`\widehat{𝐫}`$ and then re-expressing them in terms of $`t`$ and $`𝐫`$.
The case $`N=1`$ is particularly simple. It corresponds to $`q=2`$ for a fixed $`G`$ representation ($`\alpha c^2`$), or $`q=2`$ in a minimal varying $`G`$ representation ($`\alpha 1/c^2`$). In these cases we have global Lorentz transformations for coordinates
$`\widehat{t}`$ $`=`$ $`{\displaystyle \frac{t}{1+\frac{c_0t}{R}}}`$ (95)
$`\widehat{𝐫}`$ $`=`$ $`{\displaystyle \frac{𝐫}{1+\frac{c_0t}{R}}}`$ (96)
These can be inverted into transformations for $`t`$ and $`𝐫`$:
$`t^{}`$ $`=`$ $`{\displaystyle \frac{\gamma \left(t\frac{𝐯𝐫}{c_0^2}\right)}{1(\gamma 1)\frac{c_0t}{R}+\gamma \frac{𝐯𝐫}{Rc_0}}}`$ (97)
$`𝐫_{||}^{}`$ $`=`$ $`{\displaystyle \frac{\gamma \left(𝐫_{||}\frac{𝐯}{c_0}t\right)}{1(\gamma 1)\frac{c_0t}{R}+\gamma \frac{𝐯𝐫}{Rc_0}}}`$ (98)
$`𝐫_{}^{}`$ $`=`$ $`{\displaystyle \frac{𝐫_{}}{1(\gamma 1)\frac{c_0t}{R}+\gamma \frac{𝐯𝐫}{Rc_0}}}`$ (99)
where $`v`$ is the velocity between two inertial frames at the origin at $`t=0`$ (the velocity between two inertial frames varies in space and time and is proportional to $`c`$ ). The transformation we have just obtained is the Fock-Lorentz transformation.
This is an interesting result! The Fock-Lorentz transformation was first derived by Vladimir Fock in his textbook as a pedagogic curiosity. Special relativity may be derived from two postulates: the principle of relativity and the principle of constancy of the speed of light. The latter may be replaced by the requirement that the transformation be linear. Fock examined the effects of dropping the second postulate while keeping the first. He thus arrived at a fractional linear transformation identical with the one we have just derived.
We have just produced an alternative derivation, based on our dynamical equations for the field $`\psi `$. The constant $`R`$ in the Fock Lorentz transformation appears as an integration constant in our solution. Some features of the Fock transformation, not accommodated by our theory (such as an anisotropic $`c`$), can be neglected if we stay close to the origin. Similarly some features of our theory not present in Fock’s theory (such as non-integrability of infinitesimals) can be ignored in the same region. Hence it is not surprising that we have arrived at the same construction.
The Fock-Lorentz transformation has a number of interesting properties, and one of them is crucial for understanding VSL theories. If we consider a proper time interval $`\mathrm{\Delta }t_0`$ (referred to the origin) we find that this is seen in the lab frame as
$$\mathrm{\Delta }t=\frac{\mathrm{\Delta }t_0}{(1+c_0\mathrm{\Delta }t_0/R)\gamma c_0\mathrm{\Delta }t_0/R}$$
(100)
which is qualitatively very different from the usual twin paradox expression. In the standard theory the only invariant non-zero time lapse is infinity. In Fock’s theory such a role is played by $`\mathrm{\Delta }t_0=R/c_0`$; in contrast infinity is no longer invariant but can mapped into finite times and vice-versa. Suitable particle life-times may be mapped to infinity (ie: stability) by a Fock-Lorentz transformation.
Closer inspection shows that if we look at these theories from a fixed $`c`$ perspective $`t=R/c_0`$ is indeed mapped into $`\widehat{t}=\mathrm{}`$ for $`R<0`$ (or $`\widehat{t}=\mathrm{}`$ for $`R>0`$). This is obvious from (91) but also true for other values of $`N`$. Given that the two representations are globally very different one must ask which representation is more physical.
### A Interaction clocks and trans-eternal times
Clearly a change of units transforms our construction into plain Minkowski space-time. Then why not use the fixed $`c`$ representation? The point is that the correspondence is only local. We can extend the VSL empty solution beyond $`t=t_{max}=R/c_0`$, for $`R<0`$. Such extension corresponds to extending Minkowski space-time beyond $`t=\mathrm{}`$. The choice between the two representations is therefore dependent on whether this extension is physical or not.
Let us first examine $`tt_{max}`$ in units in which $`c`$ varies. In this picture $`c`$ goes to infinity at $`t_{max}`$; but this has implications on the time-scales of processes mediated by all interactions. Decay times, rates of change, etc, all depend on the $`\alpha _i`$. A typical time scale associated with a given interaction with energy $`Q`$ is
$$\tau =\frac{\mathrm{}}{\alpha ^2Q}$$
(101)
In a minimal VSL theory $`Q\mathrm{}cc^q`$, $`\alpha c^q`$, and so $`\tau 1/c^{2q+1}`$. But our sensation of time flow derives precisely from change, and this is imparted by interactions and their rates. One may therefore argue that a more solid definition of time should be tied to the rates $`\tau `$, and that a more physical clock should be obtained by making it tick to $`\tau `$. Like all other definitions of time, this definition should not affect physics (which is dimensionless); however it may lead to a clearer picture.
In the varying $`c`$ picture, the number of cycles of an interaction clock as $`tt_{max}`$ is
$$^{t_{max}}\frac{dt}{\tau }$$
(102)
which converges if $`q<0`$, that is if all $`\alpha _i`$ go to zero (all interactions switch off). Hence our claim that the space is extendable beyond $`t=t_{max}`$ is physically meaningful, if $`q<0`$.
Let us now examine the same situation in fixed $`c`$ units. Even though $`c`$ and $`\mathrm{}`$ are now fixed, this is not really just Minkowski space-time. At the very least all charges $`g_i`$ must now be variable, to produce the same changing $`\alpha _i`$. If we want to keep all parameters in (101) constant except for the $`\alpha _i`$ we should change units in the following way:
$`d\widehat{t}`$ $`=`$ $`dte^{(12q)\psi }`$ (103)
$`d\widehat{x}`$ $`=`$ $`dxe^{2q\psi }`$ (104)
$`d\widehat{E}`$ $`=`$ $`dEe^{q\psi }`$ (105)
For $`q<0`$ we have that $`t_{max}`$ is indeed mapped into $`\widehat{t}=\mathrm{}`$. However we find that the number of ticks of an interaction clock as we approach $`\widehat{t}=\mathrm{}`$
$$^{\mathrm{}}\frac{d\widehat{t}}{\widehat{\tau }}$$
(106)
converges. Hence the temporal infinity of “Minkowski” space-time in this theory is spurious. Any natural process would slow down as “fixed-$`c`$ time” went on. More and more of this “time” would be required for any interaction process to take place. Given that our sensation of time flowing is attached to these processes, we could claim that conversely we would feel that “fixed-$`c`$” time would start to go faster and faster. The fact that a finite number of physical ticks is required to reach $`\widehat{t}=\mathrm{}`$ means that any observer could in fact flow through eternity. Such Minkowski space-time is physically extendable beyond $`t=\mathrm{}`$.
We have found the first example of a situation in which the fixed $`c`$ representation, while locally equivalent to a varying $`c`$ representation, may be globally misleading. The advantage of varying $`c`$ units in this case is that they locate at a finite time distance what can in fact be reached within a finite number of cycles of an interaction clock.
## X Fast-tracks in VSL flat-space
More fascinating still is the existence of high-$`c`$ lines, which we shall call fast-tracks. These are flat space-time solutions, in theories in which $`\psi `$ is driven by a potential. We first establish the possibility of such solutions. Let $`\psi `$ be a complex scalar field, with a $`U(1)`$ symmetry which may or may not be gauged (we assume it’s gauged in what follows). Let the speed of light be given by $`c=c_0e^{|\psi |^2}`$. With these modifications we also have to modify the terms in $`a`$ and $`b`$ in (47), but not if $`a=b=0`$, as we shall assume. Let us also assume that the field is driven by a potential
$`_\psi `$ $`=`$ $`(D_\mu \psi )^{}(D^\mu \psi )V(\psi )`$ (107)
$`V(\psi )`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _\psi ^2}}(|\psi |^2\psi _0^2)^2`$ (108)
where $`\psi _0`$ is the field’s vacuum expectation value, and $`\lambda _\psi `$ is the Compton wave-length of $`\psi `$.
Let us consider a Nielsen-Olesen vortex solution to this theory, that is a solution with a boundary condition:
$$\psi =\psi _0e^{in\theta }\mathrm{as}r\mathrm{}$$
(109)
Such a solution is topologically stable. In the vortex’s core, $`|\psi |0`$ and so the speed of light is $`c_0`$. The speed of light outside the core (which is $`c_0e^{\psi _0^2}`$) is therefore much smaller. The field $`\psi `$ undergoes spontaneous symmetry breaking and the unbroken phase, realized in the string’s core, displays a much larger speed of light. An approximate solution for $`r\mathrm{}`$ is
$$\psi =(\psi _0+e^{r/\lambda _\psi })e^{in\theta }$$
(110)
Hence the string core has a width of order $`\lambda _\psi `$, which could easily be macroscopic; outside the core variations in the speed of light die off rapidly. The jump in the speed of light is exponential and depends only on $`\psi _0`$. For $`\psi _03`$, say, the speed of light could be ten orders of magnitude faster inside the string’s core. The size of the core, and the jump in $`c`$, are related to independent parameters.
What would happen if an observer travelled along the string, inside its core? Let a cylinder of high-$`c`$ connect two distant galaxies. Then inside the tube $`vc`$ (cf. (73) with $`b=0`$). Let us assume that $`vc`$ so that no relativistic effects are present. Then the observer could move very fast between these two galaxies, returning without any time dilation effects having taken place. There would not be a twin paradox - clearly this situation, if realizable, is just what intergalactic travel is begging for. In practice, to avoid different aging rates between sedentary and the nomadic twins we should keep the aging pace $`\tau `$ fixed, ie: $`q=1/2`$. Furthermore in order for the $`x^0`$ coordinate to track proper-time for all observers we should have $`\alpha =0`$ (this point will be developed further in in connection with radar echo delay experiments).
In a dual representation, in fixed $`c`$ units, fast tracks are wormholes. If $`\tau `$ is to remain unchanged, and if $`c`$ is to be fixed in the new units, then the distance between the galaxies must shrink by a corresponding factor (recall that in (1) $`\alpha =0`$, and $`\beta 0`$). Hence the fixed $`c`$ dual of the VSL theories we have proposed contain wormhole like solutions even without the presence of gravitating matter. This is due to the fact that the gravitational action is indeed very complicated in the dual picture (notice that the required unit transformation is a combination of VSL and conformal transformations).
Elsewhere we shall show how fast-tracks may appear in other theories, eg in the Bekenstein changing $`\alpha `$ theory. In such theories $`\alpha `$ is much smaller inside the string core, but all other couplings remain unchanged. Hence a nomadic twin will age much slower during the trip, since we age electromagnetically . Strong interactions just provide the nuclei for all the atoms we are made of. But we are essentially made of stable nuclei. Hence if all our nuclei aged a million years we would not notice it. Naturally in such theories one cannot avoid different aging rates between nomadic and sedentary twins - the curse of space travel.
## XI Black Holes with an edge
In we shall examine vacuum spherically symmetric solutions to all these theories. They have a common feature which can be illustrated by the well-known solution in Brans-Dicke theory (which is $`a=4`$, $`b=0`$). Using the isotropic form of the metric:
$$ds^2=fdx^{0^2}+g[dr^2+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)]$$
(111)
we have
$`f`$ $`=`$ $`f_0\left({\displaystyle \frac{1\frac{B}{r}}{1+\frac{B}{r}}}\right)^{2/\lambda }`$ (112)
$`g`$ $`=`$ $`g_0(1+B/r)^4\left({\displaystyle \frac{1\frac{B}{r}}{1+\frac{B}{r}}}\right)^{2(\lambda C1)/\lambda }`$ (113)
$`\psi `$ $`=`$ $`{\displaystyle \frac{C}{4\lambda }}\mathrm{log}\left({\displaystyle \frac{1\frac{B}{r}}{1+\frac{B}{r}}}\right)`$ (114)
where $`f_0`$, $`g_0`$, $`B`$, and $`C`$, are constants, with:
$$\lambda =[(C+1)^2C(1\omega _{bd}C/2)]^{1/2}$$
(115)
Expressions for these constants in terms of the black hole mass $`m`$ and coupling $`\omega _{bd}`$ may be found in . As we approach the horizon ($`r_h=B`$) we find that $`c`$ goes to either zero or infinity (like $`(rr_h)^N`$ with $`N`$ related to $`\omega _{bd}`$). The implication is obvious: for some parameters of the theory (in this case requiring $`q0`$<sup>1</sup><sup>1</sup>1In plain Brans-Dicke ($`a=4,b=0,q=0`$) we have that $`c0`$ but $`\tau \mathrm{}`$ in such a way that particles may enter the horizon. Hence the discussion presented here does not apply, as one would expect.) no observer may enter the horizon. The number of cycles of an interaction clock trying to enter the horizon is given by:
$$^{r_h}\frac{dt}{\tau }=^{r_h}\frac{dr}{v\tau }=^{r_h}\frac{dr}{c^{2q+2}}$$
(116)
which diverges for $`2(q+1)N>1`$.
Again this phenomenon may be interpreted variously, depending on which units are used, but all interpretations lead to the same physical conclusion (which is dimensionless): particles are unable to enter the horizon. In VSL units particles cannot enter the horizon because they stop as $`c`$ goes to zero. In fixed-$`c`$ units they cannot enter the horizon because the time rates of all interactions go to zero (as all couplings go to infinite). Old age strikes before anything “has time” to enter the horizon.
Naturally finite sized bodies suffer from further effects, analogous to tidal stresses, since they will probe gradients in $`c`$. Since $`vc`$ they get squashed if $`c0`$, or get stretched otherwise. $`c`$-induced changes of pace also induce gradients of aging across finite-sized bodies.
A pedagogic illustration, studied further in , is a muon produced close to the black hole, moving towards its horizon. Such a set up is useful, for instance, when trying to convince skeptics of the physical validity of time dilation, or Lorentz contraction (eg. the fate of cosmic ray muons entering the atmosphere). To an Earth observer, if time dilation was not a physical effect the muon should never hit the surface of the Earth. From the point of view of the muon, if the atmosphere did not appear Lorentz contracted, it should have decayed before hitting the surface. The same set up will be of assistance here. No matter how close to the horizon the muon is produced, it never reaches it. In VSL units the muon stops as it tries to enter the horizon, because its speed is close to $`c`$, but $`c`$ goes to zero. In fixed-$`c`$ units the muon moves close to the (constant) speed of light, but its lifetime goes to zero as it tries to enter the horizon. From either perspective the muon can never enter the horizon.
“Horizon” is therefore a misnomer, and we relabel it an “edge”: a boundary where $`c`$ goes to zero sufficiently fast that no object may reach it. On physical grounds we should postulate that regions beyond the edge be excised from the manifold. Then VSL manifolds may have an edge.
We arrive at a similar conclusion to Section IX. VSL and fixed-$`c`$ units are locally but not globally equivalent. The VSL picture may be globally more clear (in the case $`a=4`$, only if $`q0`$). It builds into space-time the topology perceived by actual physical processes, in this case excising regions which are physically inaccessible.
The implications for the theory of singularities are quite impressive. Even though we have a singularity at $`r=0`$, it is physically inaccessible. One may be able to prove that all singularities are subject to the same constraint. This situation was discussed in . It looks as if a stronger version of the cosmic censorship principle might apply to these theories.
## XII Conclusion
One must sympathise with the view that VSL theories are rendered objectionable by their outright violation of Lorentz invariance. However, previous attempts to make the Albrecht-Magueijo model “geometrically honest” were no less ugly than the original; and were useless for cosmology. In this paper we proposed a geometrically honest VSL theory, corresponding to a theory in which all fine structure constants are promoted to dynamical variables. A changing charge interpretation in unnecessarily complicated - so we adopt units in which $`c`$ changes, leading to a simple picture. This should not scandalise anyone.
All these theories are locally Lorentz invariant, and covariant in a sense incorporated by a generalized structure. We find that physics lives on a fibre bundle. Usually physics takes place on the tangent bundle. At each point in space-time there is a tangent space, corresponding to free falling frames in which physics is Minkowskian. We have a similar construction, but in the new units the space is not the tangent space of any coordinate patch in the manifold. It is still a vector space - but it is not the tangent vector space, except in the rare cases where the change of units is holonomic <sup>2</sup><sup>2</sup>2 The situation is more complicated for a gravitation theory based upon Fock-Lorentz space-time. Now the physics’ space at a given point is no longer a vector space, but a projective space. The fibre bundle multiplies the base manifold by a projective space at each point..
Given that a change of units maps these structures to standard covariant and local Lorentz invariant theories, one may wonder why it is worth bothering. To answer this question, throughout this paper we examined these “dual” theories. For them $`c`$ is a constant (as well as $`\mathrm{}`$ and possibly $`G`$), but naturally other quantities must vary. Indeed all couplings must change, at fixed ratios. We therefore have a theory not dissimilar from Bekenstein’s changing $`e`$ theory, but such a picture is horribly misleading for the following reasons.
Firstly all charges, not only $`e`$, will vary. But they vary at constant ratios, so that all changes may be attributed to a change in $`c`$ alone. Hence the dual theory is a theory which promotes coupling constants to dynamical variables, but then only allows rather contrived variations, ie variations which may be absorbed into a changing $`c`$. It seems therefore more natural to consider a changing $`c`$ description, even though the two descriptions are indeed operationally equivalent.
Secondly the minimal dynamics in the two frames is totally different. This results from the fact that the action has units, and therefore changes under a change of units mapping dual theories. The minimal Bekenstein-type of theories does not have the same coupling to gravity as appears in the minimal VSL formulation. Rewriting the Lagrangian of minimal VSL theories in fixed $`c`$ units leads to an unpleasant mess (Section VIII).
Thirdly and more importantly, the correspondence between VSL and its duals is only local. Globally the VSL picture can be more clear. We gave two striking examples. Fock-Lorentz space-time is just a change of units applied to Minkowski space-time; however it contains $`t>\mathrm{}`$ extensions to Minkowski space-time which are physically accessible. The horizon of a black hole may be physically impenetrable, since $`c`$ goes to zero. Calling it an edge, and excising the bit beyond the edge seems reasonable. In the dual picture no warning about the fact that a piece of the manifold is inaccessible is given. It is an afterthought to notice that all interaction strengths force the pace of aging to become very fast; thereby, for all practical purposes, preventing anything from entering the horizon.
Hence the VSL theories we have proposed are changing $`c`$ theories simply because choosing units in which $`c`$ varies leads to a simpler description. Their underlying geometrical structure is that of standard fixed $`c`$ theories subject to a change of units; a fact undeniably placing them at the pinnacle of geometrical honesty. It remains to show that these theories, applied to cosmology, perform as well as the Albrecht and Magueijo model. Such is the purpose of . In any case it is not difficult to guess the overall cosmological picture to emerge in these theories. We see that the presence of a cosmological constant $`\mathrm{\Lambda }`$ generally drives changes in $`c`$, which in turn convert the vacuum energy into radiation leading to a conventional Big Bang. However, such a Big Bang is free from the standard cosmological problems, including the cosmological constant problem. The fact that particle production occurs naturally in these theories ensures that we also solve the entropy problem.
## Acknowledgements
I would like to thank Andy Albrecht, John Barrow, Kim Baskerville, Carlo Contaldi, Tom Kibble, and Kelly Stelle for help in connection with this paper. I am grateful to the Isaac Newton Institute for support and hospitality while part of this work was done.
## Appendix - Bimetric realization of the Albrecht-Magueijo model
A theory which emulates many of the features of the Albrecht and Magueijo model (except for breaking Lorentz invariance) is the following. Let there be two metrics, $`g`$ coupling to gravitation and matter, and $`h`$ coupling to the field $`c`$ only. Then we may take the following action:
$`S`$ $`=`$ $`S_1+S_2`$ (117)
$`S_1`$ $`=`$ $`{\displaystyle d^4x\sqrt{g}\left(R2\mathrm{\Lambda }+\frac{16\pi G}{c_0^4e^{4\psi }}_m\right)}`$ (118)
$`S_2`$ $`=`$ $`{\displaystyle d^4x\sqrt{h}\left(H2\mathrm{\Lambda }_h\kappa h^{\mu \nu }_\mu \psi _\nu \psi \right)}`$ (119)
where $`g_{\mu \nu }`$ and $`h_{\mu \nu }`$ lead to two Einstein tensors $`G_{\mu \nu }`$ and $`H_{\mu \nu }`$, and $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }_h`$ are their respective (geometrical) cosmological cosntants. Varying with respect to $`g`$, $`\psi `$, and $`h`$ leads to equations of motion:
$`G_{\mu \nu }+\mathrm{\Lambda }g_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{8\pi G}{c_0^4e^{4\psi }}}T_{\mu \nu }`$ (120)
$`\mathrm{}_h\psi `$ $`=`$ $`{\displaystyle \frac{32\pi G}{c_0^4e^{4\psi }\kappa }}\sqrt{{\displaystyle \frac{g}{h}}}_m`$ (121)
$`H_{\mu \nu }+\mathrm{\Lambda }_hh_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{\kappa }{2}}\left(_\mu \psi _\nu \psi {\displaystyle \frac{1}{2}}h_{\mu \nu }_\alpha \psi ^\alpha \psi \right)`$ (122)
Hence we may derive from an action principle the property that the field $`\psi `$ does not contribute to the stress-energy tensor which acts as a source to normal space-time curvature. In we shall highlight some curiosities pertaining to these theories.
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# UTHEP-430 Perturbative renormalization factors of Δ𝑆=1 four-quark operators for domain-wall QCD
## I Introduction
$`\mathrm{\Delta }S=1`$ four-quark operators appear in the effective low-energy Hamiltonians for non-leptonic weak decays of kaons, which are relevant for the $`\mathrm{\Delta }I=1/2`$ rule and the $`CP`$-violation parameters $`ϵ^{}/ϵ`$. Although comparison of their experimental results with theoretical predictions based on the standard model is supposed to give a good testing ground for the model, it is not accomplished without any reliable non-perturbative estimates for the hadronic matrix elements of $`\mathrm{\Delta }S=1`$ four-quark operators.
Lattice QCD calculation allow us to evaluate the $`\mathrm{\Delta }S=1`$ hadronic matrix elements from first principles in QCD. While in the past decade the Wilson and the Kogut-Susskind quark actions are used to calculate the four-quark hadronic matrix elements, these quark actions have inherent defects: the explicit chiral symmetry breaking in the Wilson quark action, which causes the non-trivial operator mixing among different chiralities, and the mixture of space-time and flavor symmetries in the Kogut-Susskind quark action, which annoys us by demanding non-trivial matching between continuum and lattice operators. The domain-wall quark model is a five-dimensional Wilson fermion with free boundaries imposed on the fifth direction of $`N_s`$ layers. This quark formulation is expected to have superior features in the large $`N_s`$ limit: no fine tuning for the chiral limit, $`O(a^2)`$ scaling violation and no mixing between four-quark operators with different chiralities. Perturbative calculations have shown that these expectations are fulfilled at the one-loop level in the limit $`N_s\mathrm{}`$, and numerical studies support these features non-perturbatively. These advantageous features over other quark formulations urge us to apply the domain-wall quark to a calculation of the $`\mathrm{\Delta }S=1`$ hadronic matrix elements.
Since the Wilson coefficients in the effective low-energy Hamiltonians for non-leptonic weak decays of kaons are calculated in some continuum renormalization scheme($`e.g.,\overline{\mathrm{MS}}`$), the corresponding $`\mathrm{\Delta }S=1`$ hadronic matrix elements should be given in the same continuum scheme, which requires us to convert the matrix elements obtained by lattice simulations to those defined in continuum renormalization scheme. A necessary ingredient for this transformation is the renormalization factors matching the lattice operators to the continuum ones.
In this article we make a perturbative calculation of the renormalization factors for the $`\mathrm{\Delta }S=1`$ four-quark operators consisting of physical quark fields in the domain-wall QCD(DWQCD). The one-loop corrections of $`\mathrm{\Delta }S=1`$ four-quark operators contain the gluon exchange diagram and the penguin diagram, which leads to the following form for the relation between the lattice and continuum operators:
$`Q_i^{\mathrm{cont}}={\displaystyle \underset{j}{}}𝒵_{ij}^gQ_j^{\mathrm{latt}}+𝒵_i^{\mathrm{pen}}Q_{\mathrm{pen}}^{\mathrm{latt}}+O(g^4),`$ (1)
where $`Q_i`$ denote a set of $`\mathrm{\Delta }S=1`$ four-quark operators. The gluon exchange contributions $`𝒵_{ij}^g`$ are already calculated in our previous paper, whose results are applicable to $`\mathrm{\Delta }S=2`$ four-quark operators. In this work we evaluate the additional contribution denoted by $`𝒵_i^{\mathrm{pen}}Q_{\mathrm{pen}}^{\mathrm{latt}}`$ which originates from the penguin diagram in the $`\mathrm{\Delta }S=1`$ case. With the aid of the results in Ref., we obtain complete expressions of the renormalization factors for the $`\mathrm{\Delta }S=1`$ four-quark operators.
This paper is organized as follows. In Sec.II we present the calculations of the penguin diagram in the continuum and on the lattice. Full expressions of the renormalization factors for the $`\mathrm{\Delta }S=1`$ operators are given in Sec.III. Our conclusions are summarized in Sec.IV In Appendix we give the DWQCD action and the Feynman rule relevant for the present calculation to make this paper self-contained.
The physical quantities are expressed in lattice units and the lattice spacing $`a`$ is suppressed unless necessary. We take SU($`N`$) gauge group with the gauge coupling $`g`$, while $`N=3`$ is specified in the numerical calculations. Hereafter, repeated indices are to be summed over unless otherwise indicated.
## II Penguin diagrams
### A $`\mathrm{\Delta }S`$=1 four-quark operators
We consider the following $`\mathrm{\Delta }S=1`$ four-quark operators,
$`Q_1`$ $`=`$ $`(\overline{s}_au_b)_{VA}(\overline{u}_bd_a)_{VA},`$ (2)
$`Q_2`$ $`=`$ $`(\overline{s}_au_a)_{VA}(\overline{u}_bd_b)_{VA},`$ (3)
$`Q_3`$ $`=`$ $`(\overline{s}_ad_a)_{VA}{\displaystyle \underset{q}{}}(\overline{q}_bq_b)_{VA},`$ (4)
$`Q_4`$ $`=`$ $`(\overline{s}_ad_b)_{VA}{\displaystyle \underset{q}{}}(\overline{q}_bq_a)_{VA},`$ (5)
$`Q_5`$ $`=`$ $`(\overline{s}_ad_a)_{VA}{\displaystyle \underset{q}{}}(\overline{q}_bq_b)_{V+A},`$ (6)
$`Q_6`$ $`=`$ $`(\overline{s}_ad_b)_{VA}{\displaystyle \underset{q}{}}(\overline{q}_bq_a)_{V+A},`$ (7)
$`Q_7`$ $`=`$ $`{\displaystyle \frac{3}{2}}(\overline{s}_ad_a)_{VA}{\displaystyle \underset{q}{}}e^q(\overline{q}_bq_b)_{V+A},`$ (8)
$`Q_8`$ $`=`$ $`{\displaystyle \frac{3}{2}}(\overline{s}_ad_b)_{VA}{\displaystyle \underset{q}{}}e^q(\overline{q}_bq_a)_{V+A},`$ (9)
$`Q_9`$ $`=`$ $`{\displaystyle \frac{3}{2}}(\overline{s}_ad_a)_{VA}{\displaystyle \underset{q}{}}e^q(\overline{q}_bq_b)_{VA},`$ (10)
$`Q_{10}`$ $`=`$ $`{\displaystyle \frac{3}{2}}(\overline{s}_ad_b)_{VA}{\displaystyle \underset{q}{}}e^q(\overline{q}_bq_a)_{VA},`$ (11)
where $`V\pm A`$ indicates the chiral structures $`\gamma _\mu (1\pm \gamma _5)`$. $`a,b`$ denote color indices and $`e^q`$ are quark charges: $`2/3`$ for the up-like quarks and $`1/3`$ for the down-like quarks. The flavor sum is intended over those which are active below the renormalization scale of the operators. This set of operator basis closes under QCD and QED renormalization. We classify these 10 operators into the following four types:
$`Q_1`$ $`=`$ $`(\overline{s}_a\mathrm{\Gamma }_X^\mu u_b)(\overline{u}_b\mathrm{\Gamma }_X^\mu d_a),`$ (12)
$`Q_2`$ $`=`$ $`(\overline{s}_a\mathrm{\Gamma }_X^\mu u_a)(\overline{u}_b\mathrm{\Gamma }_X^\mu d_b),`$ (13)
$`Q_i^e`$ $`=`$ $`(\overline{s}_a\mathrm{\Gamma }_X^\mu d_b){\displaystyle \underset{q}{}}\alpha _i^q(\overline{q}_b\mathrm{\Gamma }_Y^\mu q_a)i=4,6,8,10,`$ (14)
$`Q_i^o`$ $`=`$ $`(\overline{s}_a\mathrm{\Gamma }_X^\mu d_a){\displaystyle \underset{q}{}}\alpha _i^q(\overline{q}_b\mathrm{\Gamma }_Y^\mu q_b)i=3,5,7,9,`$ (15)
where $`\mathrm{\Gamma }_X^\mu =\gamma _\mu (1\gamma _5)`$ and $`\mathrm{\Gamma }_Y^\mu =\gamma _\mu (1\pm \gamma _5)`$. $`\alpha _i^q`$ are given by
$`\alpha _i^q`$ $`=`$ $`1i=3,4,5,6,`$ (16)
$`\alpha _i^q`$ $`=`$ $`{\displaystyle \frac{3}{2}}e^qi=7,8,9,10.`$ (17)
### B Continuum calculation
It is instructive to first show the calculation of the penguin diagram in the continuum regularization schemes. For the present calculation we employ the Naive Dimensional Regularization(NDR) and the Dimensional Reduction(DRED) as the ultraviolet regularization scheme, in each of which the loop momenta of the Feynman integrals are defined in $`D`$ dimensions parameterized by $`D=42ϵ`$ ($`ϵ>0`$). The major difference between both regularization schemes are found in the definitions of the Dirac matrices: NDR scheme defines the Dirac matrices in $`D`$ dimension, while DRED scheme retains them in four dimensions.
We calculate the following Green function with massless quarks:
$$Q_i_{ab;cd}Q_is_a(p_1)\overline{d}_b(p_2)q_c(p_3)\overline{q}_d(p_4),$$
(18)
where $`a,b,c,d`$ label color indices, while spinor indices are suppressed. Truncating the external quark propagators from $`Q_i`$, where we multiply $`ip/_j`$ ($`j=1,2,3,4`$) on $`Q_i`$, we obtain the vertex functions, which is written in the following form up to the one-loop level
$$\left[\mathrm{\Lambda }_i(p)\right]_{ab;cd}=\left[\mathrm{\Lambda }_i^{(0)}(p)+\mathrm{\Lambda }_i^{(1g)}(p)+\mathrm{\Lambda }_i^{(1p)}(p)\right]_{ab;cd},$$
(19)
where the momentum $`p`$ is defined by $`p=p_1p_2=p_4p_3`$. The superscript $`(i)`$ refers to the $`i`$-th loop level. At the one-loop level $`\mathrm{\Lambda }_i^{(1g)}`$ and $`\mathrm{\Lambda }_i^{(1p)}`$ represent the gluon exchange diagram and the penguin diagram, respectively. In this paper, however, we do not consider $`\mathrm{\Lambda }_i^{(1g)}`$ because the gluon exchange contributions to the renormalization factors are already calculated in our previous paper.
The tree-level vertex functions $`\mathrm{\Lambda }_i^{(0)}`$ are given by
$`\left[\mathrm{\Lambda }_1^{(0)}\right]_{ab;cd}`$ $`=`$ $`\left[1\stackrel{~}{}1\right]_{ab;cd}\mathrm{\Gamma }_X^\mu \mathrm{\Gamma }_X^\mu ,`$ (20)
$`\left[\mathrm{\Lambda }_2^{(0)}\right]_{ab;cd}`$ $`=`$ $`\left[1\stackrel{~}{}1\right]_{ab;cd}\mathrm{\Gamma }_X^\mu \mathrm{\Gamma }_X^\mu ,`$ (21)
$`\left[\mathrm{\Lambda }_i^{e(0)}\right]_{ab;cd}`$ $`=`$ $`\alpha _i^u\left[1\stackrel{~}{}1\right]_{ab;cd}\mathrm{\Gamma }_X^\mu \mathrm{\Gamma }_Y^\mu ,`$ (22)
$`\left[\mathrm{\Lambda }_i^{o(0)}\right]_{ab;cd}`$ $`=`$ $`\alpha _i^u\left[1\stackrel{~}{}1\right]_{ab;cd}\mathrm{\Gamma }_X^\mu \mathrm{\Gamma }_Y^\mu ,`$ (23)
where $``$ represent the tensor structures in the Dirac spinor space and $`\stackrel{~}{}`$ and $`\stackrel{~}{}`$ act on the color space as
$`\left[1\stackrel{~}{}1\right]_{ab;cd}`$ $``$ $`\delta _{ad}\delta _{cb},`$ (24)
$`\left[1\stackrel{~}{}1\right]_{ab;cd}`$ $``$ $`\delta _{ab}\delta _{cd}.`$ (25)
We take $`q=u`$ in eq.(18) for convenience. This special choice, however, does not affect the final results for the renormalization factors.
The penguin diagram is illustrated in Fig. 1, which yields the one-loop level vertex functions $`\mathrm{\Lambda }_i^{(1p)}`$,
$`\left[\mathrm{\Lambda }_1^{(1p)}\right]_{ab;cd}`$ $`=`$ $`0,`$ (26)
$`\left[\mathrm{\Lambda }_2^{(1p)}\right]_{ab;cd}`$ $`=`$ $`g^2\left[T^A\stackrel{~}{}T^B\right]_{ab;cd}G_{\mu \nu }^{AB}(p)L_{\alpha \beta }(p)\mathrm{\Gamma }_X^\delta \gamma _\alpha \gamma _\mu \gamma _\beta \mathrm{\Gamma }_X^\delta \gamma _\nu ,`$ (27)
$`\left[\mathrm{\Lambda }_i^{e(1p)}\right]_{ab;cd}`$ $`=`$ $`+g^2\left[T^A\stackrel{~}{}T^B\right]_{ab;cd}G_{\mu \nu }^{AB}(p)L_{\alpha \beta }(p){\displaystyle \underset{q}{}}\alpha _i^q\mathrm{tr}(\mathrm{\Gamma }_Y^\delta \gamma _\alpha \gamma _\mu \gamma _\beta )\mathrm{\Gamma }_X^\delta \gamma _\nu ,`$ (28)
$`\left[\mathrm{\Lambda }_i^{o(1p)}\right]_{ab;cd}`$ $`=`$ $`g^2\left[T^A\stackrel{~}{}T^B\right]_{ab;cd}G_{\mu \nu }^{AB}(p)L_{\alpha \beta }(p)\left[\alpha _i^d\mathrm{\Gamma }_X^\delta \gamma _\alpha \gamma _\mu \gamma _\beta \mathrm{\Gamma }_Y^\delta +\alpha _i^s\mathrm{\Gamma }_Y^\delta \gamma _\alpha \gamma _\mu \gamma _\beta \mathrm{\Gamma }_X^\delta \right]\gamma _\nu ,`$ (29)
where $`T^A`$ ($`A=1,\mathrm{},N^21`$) are generators of color SU($`N`$). $`G_{\mu \nu }^{AB}`$ is the gauge propagator given by
$`G_{\mu \nu }^{AB}(p)=\{\begin{array}{cc}\delta _{AB}\delta _{\mu \nu }{\displaystyle \frac{1}{p^2}}\hfill & \text{NDR},\hfill \\ \delta _{AB}\stackrel{~}{\delta }_{\mu \nu }{\displaystyle \frac{1}{p^2}}\hfill & \text{DRED},\hfill \end{array}`$ (32)
with $`D`$-dimensional metric tensor $`\delta _{\mu \nu }`$ and the four-dimensional one $`\stackrel{~}{\delta }_{\mu \nu }`$. $`L_{\alpha \beta }(p)`$ denote the integral of the quark loop, whose result is
$`L_{\alpha _\beta }(p)`$ $`=`$ $`{\displaystyle \frac{1}{6(4\pi )^2}}\left[\left({\displaystyle \frac{1}{\overline{ϵ}}}+\mathrm{log}(\mu ^2/p^2)+{\displaystyle \frac{5}{3}}+1\right){\displaystyle \frac{p^2}{2}}\delta _{\alpha \beta }+\left({\displaystyle \frac{1}{\overline{ϵ}}}+\mathrm{log}(\mu ^2/p^2)+{\displaystyle \frac{5}{3}}\right)p_\alpha p_\beta \right],`$ (33)
where $`1/\overline{ϵ}=1/ϵ\gamma +\mathrm{ln}|4\pi |`$ and $`\mu `$ is the renormalization scale. We remark that the metric tensor $`\delta _{\alpha \beta }`$ is defined in $`D`$ dimension both for the NDR and DRED schemes. The products of the Dirac matrices are reduced as follows:
$`\mathrm{tr}\left[\mathrm{\Gamma }_Y^\nu \gamma _\alpha \gamma _\mu \gamma _\beta \right]\delta _{\alpha \beta }`$ $`=`$ $`4\times \{\begin{array}{cc}(2D)\delta _{\mu \nu }\hfill & \text{NDR},\hfill \\ 2\delta _{\mu \nu }D\stackrel{~}{\delta }_{\mu \nu }\hfill & \text{DRED},\hfill \end{array}`$ (36)
$`\mathrm{tr}\left[\mathrm{\Gamma }_Y^\nu \gamma _\alpha \gamma _\mu \gamma _\beta \right]p_\alpha p_\beta `$ $`=`$ $`4\times \{\begin{array}{cc}2p_\mu p_\nu \delta _{\mu \nu }p^2\hfill & \text{NDR},\hfill \\ 2p_\mu p_\nu \stackrel{~}{\delta }_{\mu \nu }p^2\hfill & \text{DRED},\hfill \end{array}`$ (39)
$`\mathrm{\Gamma }_X^\nu \gamma _\alpha \gamma _\mu \gamma _\beta \mathrm{\Gamma }_Y^\nu \delta _{\alpha \beta }`$ $`=`$ $`\delta _{XY}\times \{\begin{array}{cc}2(D2)^2\gamma _\mu (1\gamma _5)\hfill & \text{NDR},\hfill \\ 4\left[2\overline{\gamma }_\mu D\gamma _\mu \right](1\gamma _5)\hfill & \text{DRED},\hfill \end{array}`$ (42)
$`\mathrm{\Gamma }_X^\nu \gamma _\alpha \gamma _\mu \gamma _\beta \mathrm{\Gamma }_Y^\nu p_\alpha p_\beta `$ $`=`$ $`\delta _{XY}\times \{\begin{array}{cc}2(D2)[p^2\gamma _\mu 2p_\mu p/](1\gamma _5)\hfill & \text{NDR},\hfill \\ 4[p^2\gamma _\mu 2p_\mu p/](1\gamma _5)\hfill & \text{DRED},\hfill \end{array}`$ (45)
where $`\overline{\gamma }_\mu =\delta _{\mu \nu }\gamma _\nu `$ in the DRED scheme, and we choose $`E_\mu =\overline{\gamma }_\mu \frac{D}{4}\gamma _\mu `$ as an evanescent operator in the DRED scheme. It should be noted that the terms proportional to $`p/`$ are eliminated by imposing on-shell conditions on external quark states.
After some algebra we obtain the following expressions for the vertex functions (26)$``$(29):
$`\left[\mathrm{\Lambda }_2^{(1p)}\right]_{ab;cd}`$ $`=`$ $`{\displaystyle \frac{g^2}{12\pi ^2}}\left[T^A\stackrel{~}{}T^A\right]_{ab;cd}\gamma _\mu (1\gamma _5)\gamma _\mu L_i^{\mathrm{cont}}`$ (46)
$`\left[\mathrm{\Lambda }_i^{e(1p)}\right]_{ab;cd}`$ $`=`$ $`{\displaystyle \frac{g^2}{12\pi ^2}}{\displaystyle \underset{q}{}}\alpha _i^q\left[T^A\stackrel{~}{}T^A\right]_{ab;cd}\gamma _\mu (1\gamma _5)\gamma _\mu L_i^{\mathrm{cont}}`$ (47)
$`\left[\mathrm{\Lambda }_i^{o(1p)}\right]_{ab;cd}`$ $`=`$ $`\delta _{XY}{\displaystyle \frac{g^2}{12\pi ^2}}(\alpha _i^s+\alpha _i^d)\left[T^A\stackrel{~}{}T^A\right]_{ab;cd}\gamma _\mu (1\gamma _5)\gamma _\mu L_i^{\mathrm{cont}}`$ (48)
where
$`L_i^{\mathrm{cont}}`$ $`=`$ $`{\displaystyle \frac{1}{\overline{ϵ}}}+\mathrm{log}(\mu ^2/p^2)+{\displaystyle \frac{5}{3}}+c_i`$ (49)
with
$`c_i=\{\begin{array}{cc}1\hfill & Q_2\text{ and }Q_i^o\text{ in NDR},\hfill \\ 0\hfill & Q_i^e\text{ in NDR},\hfill \\ \frac{1}{4}\hfill & Q_2\text{}Q_i^e\text{ and }Q_i^o\text{ in DRED}.\hfill \end{array}`$ (53)
The pole term $`1/\overline{ϵ}`$ is subtracted from $`L_i^{\mathrm{cont}}`$ in the $`\overline{\mathrm{MS}}`$ scheme.
From the tree-level vertex functions $`\mathrm{\Lambda }_i^{(0)}`$ in eqs.(20)$``$(23) we find
$`\left[T^A\stackrel{~}{}T^A\right]_{ab;cd}\gamma _\mu (1\gamma _5)\gamma _\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[1\stackrel{~}{}1{\displaystyle \frac{1}{N}}1\stackrel{~}{}1\right]_{ab;cd}\gamma _\mu (1\gamma _5){\displaystyle \frac{1}{2}}[\gamma _\mu (1\gamma _5)+\gamma _\mu (1+\gamma _5)]`$ (54)
$`=`$ $`{\displaystyle \frac{1}{4}}\left[\mathrm{\Lambda }_4^{(0)}+\mathrm{\Lambda }_6^{(0)}{\displaystyle \frac{1}{N}}(\mathrm{\Lambda }_3^{(0)}+\mathrm{\Lambda }_5^{(0)})\right]_{ab;cd}.`$ (55)
Using this relation we finally obtain
$`\mathrm{\Lambda }_i^{(1p)}`$ $`=`$ $`{\displaystyle \frac{g^2}{12\pi ^2}}{\displaystyle \frac{1}{4}}C(Q_i)L_i^{\mathrm{cont}}\left[\mathrm{\Lambda }_4^{(0)}+\mathrm{\Lambda }_6^{(0)}{\displaystyle \frac{1}{N}}(\mathrm{\Lambda }_3^{(0)}+\mathrm{\Lambda }_5^{(0)})\right]`$ (56)
with $`C(Q_2)`$=1, $`C(Q_3)=2`$, $`C(Q_4)=C(Q_6)=f_q`$, $`C(Q_8)=f_uf_d/2`$, $`C(Q_9)=1`$, $`C(Q_{10})=f_uf_d/2`$ and $`C(Q_i)=0`$ for other $`i`$, where $`f_q`$, $`f_u`$ and $`f_d`$ denote the number of flavors, the number of charge 2/3 up-like quarks and the number of charge $`1/3`$ down-like quarks, respectively.
### C Lattice calculation
Let us turn to the calculation of the vertex functions on the lattice. In this subsection we use the same notations for quantities defined on the lattice as those for their counterparts in the continuum.
We consider the Green function of eq. (18) on the lattice. At the one-loop level the penguin diagram in Fig. 1 contributes to the vertex functions:
$`\left[\mathrm{\Lambda }_1^{(1p)}\right]_{ab;cd}`$ $`=`$ $`0,`$ (57)
$`\left[\mathrm{\Lambda }_2^{(1p)}\right]_{ab;cd}`$ $`=`$ $`g^2\left[T^A\stackrel{~}{}T^B\right]_{ab;cd}(1w_0^2)^3G_{\mu \nu }^{AB}(p)\mathrm{\Gamma }_X^\delta L_\mu (p)\mathrm{\Gamma }_X^\delta \gamma _\nu ,`$ (58)
$`\left[\mathrm{\Lambda }_i^{e(1p)}\right]_{ab;cd}`$ $`=`$ $`+g^2\left[T^A\stackrel{~}{}T^B\right]_{ab;cd}(1w_0^2)^3G_{\mu \nu }^{AB}(p){\displaystyle \underset{q}{}}\alpha _i^q\mathrm{tr}(\mathrm{\Gamma }_Y^\delta L_\mu (p))\mathrm{\Gamma }_X^\delta \gamma _\nu ,`$ (59)
$`\left[\mathrm{\Lambda }_i^{o(1p)}\right]_{ab;cd}`$ $`=`$ $`g^2\left[T^A\stackrel{~}{}T^B\right]_{ab;cd}(1w_0^2)^3G_{\mu \nu }^{AB}(p)\left[\alpha _i^d\mathrm{\Gamma }_X^\delta L_\mu (p)\mathrm{\Gamma }_Y^\delta +\alpha _i^s\mathrm{\Gamma }_Y^\delta L_\mu (p)\mathrm{\Gamma }_X^\delta \right]\gamma _\nu ,`$ (60)
where $`(1w_0^2)^3`$ is the overlap factor to the four-dimensional quark fields which emerges through the truncation of the external quark legs by multiplying $`ip/_i`$ ($`i=1,\mathrm{},4`$) on the Green function. $`L_\mu `$ denotes the integral of the quark loop, which is expressed by
$`L_\mu (p)`$ $`=`$ $`{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}\{P_RS_F(k)_{1s}+P_LS_F(k)_{N_ss}\}(i)v_\mu (kp/2)_{st}`$ (62)
$`\times \{S_F(kp)_{t1}P_L+S_F(kp)_{tN_s}P_R\}`$
$`=`$ $`{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}I_\mu (k,p),`$ (63)
with
$`I_\mu (k,p)`$ $`=`$ $`c_\mu (kp/2){\displaystyle \frac{1}{1e^{(\alpha +\alpha ^{})}}}\gamma _\mu c_\mu (kp/2){\displaystyle \frac{1}{FF^{}(1e^{(\alpha +\alpha ^{})})}}s/\gamma _\mu s/^{}`$ (65)
$`+irs_\mu (kp/2){\displaystyle \frac{e^\alpha }{F^{}(1e^{(\alpha +\alpha ^{})})}}is/^{}+irs_\mu (kp/2){\displaystyle \frac{e^\alpha ^{}}{F(1e^{(\alpha +\alpha ^{})})}}is/.`$
where $`c_\mu (k)=\mathrm{cos}k_\mu `$, $`s_\mu (k)=\mathrm{sin}k_\mu `$, $`s_\mu ^{}=s_\mu (kp)`$ $`\alpha =\alpha (k)`$, $`\alpha ^{}=\alpha (kp)`$, $`F=F(k)`$ and $`F^{}=F(kp)`$. The expressions of $`v_\mu `$, $`S_F`$, $`\alpha `$ and $`F`$ are given in Appendix.
Since the function $`L_\mu (p)`$ has an infrared divergence for $`p^20`$, we consider extracting its divergent part by employing an analytically integrable expression $`\stackrel{~}{I}_\mu (k,p)`$ which has the same infrared behavior as $`I_\mu (k,p)`$,
$`L_\mu (p)`$ $`=`$ $`{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}\stackrel{~}{I}_\mu (k,p)+{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}\left\{I_\mu (k,p)\stackrel{~}{I}_\mu (k,p)\right\},`$ (66)
where we choose
$$\stackrel{~}{I}_\mu (k,p)=\theta (\mathrm{\Lambda }^2k^2)(1w_0^2)\frac{1}{ik/}\gamma _\mu \frac{1}{i(k/p/)},$$
(67)
with $`\mathrm{\Lambda }`$ $`(\pi )`$ a cut-off. Now the second term in the right hand side of eq.(66) is regular in terms of $`p`$, we can make a Taylor expansion,
$`L_\mu (p)={\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}\stackrel{~}{I}_\mu (k,p)`$ $`+`$ $`{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}\left\{I_\mu (k,0)\stackrel{~}{I}_\mu (k,0)\right\}`$ (68)
$`+`$ $`{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}p_\rho {\displaystyle \frac{}{p_\rho }}\left\{I_\mu (k,p)\stackrel{~}{I}_\mu (k,p)\right\}|_{p=0}`$ (69)
$`+`$ $`{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}{\displaystyle \frac{1}{2}}p_\rho p_\lambda {\displaystyle \frac{^2}{p_\rho p_\lambda }}\left\{I_\mu (k,p)\stackrel{~}{I}_\mu (k,p)\right\}|_{p=0},`$ (70)
where $`\rho ,\lambda =1,2,3,4`$. Before evaluating each term of the right hand side, let us consider the transformation properties of $`L_\mu `$ under various discrete symmetries, from which we can predict the types of terms to be allowed in $`L_\mu `$. The same kind of discussion is found in the Wilson quark case.
Under the operation of charge conjugation matrix $`C=\gamma _4\gamma _2`$ we find
$`Cv_\mu (k)_{st}C^1`$ $`=`$ $`v_\mu ^T(k)_{N_s+1t,N_s+1s},`$ (71)
$`CS_F(k)_{st}C^1`$ $`=`$ $`S_F^T(k)_{N_s+1t,N_s+1s},`$ (72)
$`CL_\mu (p)C^1`$ $`=`$ $`L_\mu ^T(p).`$ (73)
Parity transformation gives
$`\gamma _4v_4(k_4,\stackrel{}{k})_{st}\gamma _4`$ $`=`$ $`v_4(k_4,\stackrel{}{k})_{N_s+1s,N_s+1t},`$ (74)
$`\gamma _4v_i(k_4,\stackrel{}{k})_{st}\gamma _4`$ $`=`$ $`v_i(k_4,\stackrel{}{k})_{N_s+1s,N_s+1t},`$ (75)
$`\gamma _4S_F(k_4,\stackrel{}{k})_{st}\gamma _4`$ $`=`$ $`S_F(k_4,\stackrel{}{k})_{N_s+1s,N_s+1t},`$ (76)
$`\gamma _4L_4(p_4,\stackrel{}{p})\gamma _4`$ $`=`$ $`L_4(p_4,\stackrel{}{p}),`$ (77)
$`\gamma _4L_i(p_4,\stackrel{}{p})\gamma _4`$ $`=`$ $`L_i(p_4,\stackrel{}{p}),`$ (78)
where $`i=1,2,3`$. These discrete symmetries restrict the form of $`L_\mu `$ as
$$L_\mu (p)=\frac{1}{a^2}c_0\gamma _\mu +\frac{1}{a}c_1i\sigma _{\mu \nu }p_\nu +c_{2a}\gamma _\mu p^2+c_{2b}p_\mu p/+c_{2c}\gamma _\mu p_\mu ^2+O(a),$$
(79)
with $`\sigma _{\mu \nu }=[\gamma _\mu ,\gamma _\nu ]/2`$, where the massless quark case is considered. We can eliminate further terms by the Ward-Takahashi identity:
$$2\mathrm{s}\mathrm{i}\mathrm{n}(p_\mu a/2)v_\mu (kp/2)_{st}=[S_F(k)^1]_{st}[S_F(kp)^1]_{st},$$
(80)
which leads to
$`2\mathrm{s}\mathrm{i}\mathrm{n}(p_\mu a/2)L_\mu (p)`$ $`=`$ $`i{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}[P_R\{S_F(k)_{11}S_F(kp)_{11}\}P_L`$ (84)
$`+P_R\left\{S_F(k)_{1N_s}S_F(kp)_{1N_s}\right\}P_R`$
$`+P_L\left\{S_F(k)_{N_s1}S_F(kp)_{N_s1}\right\}P_L`$
$`+P_L\{S_F(k)_{N_sN_s}S_F(kp)_{N_sN_s}\}P_R]`$
$`=`$ $`0`$ (85)
where the periodicity of $`S_F`$ is used. Under the requirement that the terms of the left hand side must vanish order by order in terms of the lattice spacing, we find $`c_0=c_{2c}=0`$ and $`c_{2a}=c_{2b}`$ in eq.(79), with which the expression of $`L_\mu `$ is simplified as
$$L_\mu (p)=\frac{1}{a}c_1i\sigma _{\mu \nu }p_\nu +c_{2a}(\gamma _\mu p^2p_\mu p/)+O(a).$$
(86)
Here we should note that
$`\gamma _\delta \sigma _{\mu \nu }\gamma _\delta =0,`$ (87)
$`\mathrm{tr}\left(\sigma _{\mu \nu }\gamma _\delta (1\pm \gamma _5)\right)=0,`$ (88)
which means $`\sigma _{\mu \nu }p_\nu `$ term in eq.(86) does not contribute to the wave functions (57)$``$(60).
Now we see that $`L_\mu (0)=0`$ and linear term in $`p`$ is irrelevant in the expansion (70), we focus on the remaining terms. Performing the integration of $`\stackrel{~}{I}_\mu (k,p)\stackrel{~}{I}_\mu (k,0)`$ analytically, we obtain
$`{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}\left\{\stackrel{~}{I}_\mu (k,p)\stackrel{~}{I}_\mu (k,0)\right\}`$ $`=`$ $`(1w_0^2){\displaystyle \frac{1}{16\pi ^2}}[(p^2\gamma _\mu +p_\mu p/){\displaystyle \frac{1}{3}}(\mathrm{log}(\mathrm{\Lambda }^2/p^2)+5/6){\displaystyle \frac{p^2}{6}}\gamma _\mu ].`$ (89)
As for the last term in the expansion of eq. (70), some tedious algebra leads to the following expression for the integrand:
$`p_\rho p_\lambda {\displaystyle \frac{^2}{p_\rho p_\lambda }}\left\{I_\mu (p)\stackrel{~}{I}_\mu (p)\right\}|_{p=0}`$ (90)
$`=p^2\gamma _\mu \left(R_{\mu \nu }^bR_{\mu \nu \nu }^d\right)+p_\mu p/\left(2R_{\mu \nu }^c+R_{\mu \mu \nu }^d+R_{\mu \nu \mu }^d\right)`$ (91)
$`+p_\mu ^2\gamma _\mu \left(R_\mu ^a+R_{\mu \mu }^bR_{\mu \nu }^b+2R_{\mu \mu }^c2R_{\mu \nu }^c+R_{\mu \mu \mu }^d+R_{\mu \nu \nu }^dR_{\mu \mu \nu }^dR_{\mu \nu \mu }^d\right)`$ (92)
$`+\gamma _\nu \gamma _\alpha p_\nu p_\alpha \gamma _\mu (R_{\mu \alpha \nu }^dR_{\mu \nu \alpha }^d),`$ (93)
where $`R_\mu ^a`$, $`R_{\mu \nu }^b`$, $`R_{\mu \nu }^c`$ and $`R_{\mu \nu \alpha }^d`$ are given by
$`R_\mu ^a`$ $`=`$ $`e^\alpha {\displaystyle \frac{c_\mu }{4(e^\alpha e^\alpha )}}+{\displaystyle \frac{s_\mu ^2f_\mu }{(e^\alpha e^\alpha )^2}}`$ (98)
$`+{\displaystyle \frac{c_\mu (10s_\mu ^2s^2)}{4F^2(1e^{2\alpha })}}{\displaystyle \frac{s_\mu ^2s^2F_\mu }{F}}{\displaystyle \frac{4c_\mu s_\mu ^4G_{\mu \mu }}{F}}`$
$`+r{\displaystyle \frac{5s_\mu ^2}{4F(e^\alpha e^\alpha )}}2re^\alpha s_\mu ^4G_{\mu \mu }`$
$`+{\displaystyle \frac{rs_\mu ^2}{4F(e^\alpha e^\alpha )}}+{\displaystyle \frac{2rs_\mu ^4e^\alpha }{F(e^\alpha e^\alpha )^3}}\left[f_\mu ^2(e^\alpha +e^\alpha )g_{\mu \mu }(e^\alpha e^\alpha )\right]`$
$`32(1w_0^2)\theta (\mathrm{\Lambda }^2k^2){\displaystyle \frac{k_\mu ^4}{(k^2)^4}},`$
$`R_{\mu \nu }^b`$ $`=`$ $`{\displaystyle \frac{c_\mu }{(e^\alpha e^\alpha )^3}}\left[s_\nu ^2f_\nu ^2(e^\alpha +e^\alpha )(h_\nu +s_\nu ^2g_{\nu \nu })(e^\alpha e^\alpha )\right]`$ (103)
$`{\displaystyle \frac{c_\mu s_\nu ^2}{F^2(1e^{2\alpha })}}+c_\mu {\displaystyle \frac{(2s_\mu ^2s^2)(H_\nu +s_\nu ^2G_{\nu \nu })}{F}}`$
$`+re^\alpha s_\mu ^2(s_\nu ^2G_{\nu \nu }+H_\nu )`$
$`+re^\alpha {\displaystyle \frac{s_\mu ^2}{F(e^\alpha e^\alpha )^3}}\left[(s_\nu ^2g_{\nu \nu }+h_\nu )(e^\alpha e^\alpha )s_\nu ^2f_\nu ^2(e^\alpha +e^\alpha )\right]`$
$`(1w_0^2)\theta (\mathrm{\Lambda }^2k^2){\displaystyle \frac{2(k^22k_\mu ^2)(4k_\nu ^2k^2)}{(k^2)^4}},`$
$`R_{\mu \nu }^c`$ $`=`$ $`e^\alpha {\displaystyle \frac{c_\nu s_\mu ^2}{2F^2(e^\alpha e^\alpha )}}+s_\nu ^2F_\nu {\displaystyle \frac{s_\mu ^2}{F}}+{\displaystyle \frac{2c_\mu s_\mu ^2s_\nu ^2G_{\mu \nu }}{F}}`$ (107)
$`r{\displaystyle \frac{c_\mu c_\nu }{2F(e^\alpha e^\alpha )}}+re^\alpha \left[s_\mu ^2s_\nu ^2G_{\mu \nu }s_\mu ^2c_\nu F_\mu c_\mu s_\nu ^2F_\nu /2\right]`$
$`+re^\alpha {\displaystyle \frac{c_\mu s_\nu ^2f_\nu }{2F(e^\alpha e^\alpha )^2}}+re^\alpha {\displaystyle \frac{s_\mu ^2s_\nu ^2}{F(e^\alpha e^\alpha )^3}}\left[g_{\mu \nu }(e^\alpha e^\alpha )f_\mu f_\nu (e^\alpha +e^\alpha )\right]`$
$`+(1w_0^2)\theta (\mathrm{\Lambda }^2k^2)16k_\nu ^2{\displaystyle \frac{k_\mu ^2}{(k^2)^4}},`$
$`R_{\mu \nu \alpha }^d`$ $`=`$ $`2{\displaystyle \frac{c_\mu s_\nu ^2F_\nu c_\alpha }{F}}4(1w_0^2)\theta (\mathrm{\Lambda }^2k^2){\displaystyle \frac{k_\nu ^2}{(k^2)^3}},`$ (108)
with
$`f_\mu `$ $`=`$ $`{\displaystyle \frac{rW+c_\mu +r\mathrm{cosh}\alpha }{W\mathrm{sinh}\alpha }},`$ (109)
$`h_\mu `$ $`=`$ $`{\displaystyle \frac{c_\mu rW+c_\mu ^2s_\mu ^2+rc_\mu \mathrm{cosh}\alpha }{W\mathrm{sinh}\alpha }},`$ (110)
$`g_{\mu \nu }`$ $`=`$ $`g_{\nu \mu }={\displaystyle \frac{f_\mu f_\nu }{\mathrm{tanh}\alpha }}+{\displaystyle \frac{r^2}{W\mathrm{sinh}\alpha }}+r{\displaystyle \frac{f_\mu +f_\nu }{W}},`$ (111)
$`F_\mu `$ $`=`$ $`e^{2\alpha }{\displaystyle \frac{r+Wf_\mu }{F^2(e^\alpha e^\alpha )}}{\displaystyle \frac{f_\mu }{F(e^\alpha e^\alpha )^2}},`$ (112)
$`H_\mu `$ $`=`$ $`e^{2\alpha }{\displaystyle \frac{rc_\mu Wh_\mu }{F^2(e^\alpha e^\alpha )}}+{\displaystyle \frac{h_\mu }{F(e^\alpha e^\alpha )^2}},`$ (113)
$`G_{\mu \nu }`$ $`=`$ $`G_{\nu \mu }={\displaystyle \frac{1}{F}}\left[{\displaystyle \frac{g_{\mu \nu }}{(e^\alpha e^\alpha )^2}}f_\mu f_\nu {\displaystyle \frac{e^\alpha +e^\alpha }{(e^\alpha e^\alpha )^3}}\right]+e^\alpha {\displaystyle \frac{f_\mu (r+Wf_\nu )+f_\nu (r+Wf_\mu )}{F^2(e^\alpha e^\alpha )^2}}`$ (114)
$``$ $`{\displaystyle \frac{e^\alpha }{F^3(e^\alpha e^\alpha )}}\left[2e^{2\alpha }(r+Wf_\mu )(r+Wf_\nu )+e^\alpha F\{r(f_\mu +f_\nu )+Wf_\mu f_\nu +Wg_{\mu \nu }\}\right].`$ (115)
We choose $`r=1`$ in this calculation. Using the results of eqs.(89) and (93), $`L_\mu (p)`$ is expressed as
$`L_\mu (p)`$ $`=`$ $`{\displaystyle \frac{(1w_0^2)}{16\pi ^2}}[(p_\mu p/p^2\gamma _\mu ){\displaystyle \frac{1}{3}}(\mathrm{log}(\mathrm{\Lambda }^2/p^2)+5/6)`$ (117)
$`+Ap^2\gamma _\mu +Bp_\mu p/+\mathrm{\Delta }Rp_\mu ^2\gamma _\mu +\delta R\gamma _\nu \gamma _\alpha p_\nu p_\alpha \gamma _\mu ],`$
where
$`A`$ $`=`$ $`{\displaystyle \frac{1}{6}}+{\displaystyle \frac{16\pi ^2}{1w_0^2}}{\displaystyle \frac{1}{2}}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}\left\{R_{\mu \nu }^bR_{\mu \nu \nu }^d\right\},`$ (118)
$`B`$ $`=`$ $`{\displaystyle \frac{16\pi ^2}{1w_0^2}}{\displaystyle \frac{1}{2}}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}\left\{2R_{\mu \nu }^c+R_{\mu \mu \nu }^d+R_{\mu \nu \mu }^d\right\},`$ (119)
$`\mathrm{\Delta }R`$ $`=`$ $`{\displaystyle \frac{16\pi ^2}{1w_0^2}}{\displaystyle \frac{1}{2}}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}\left\{R_\mu ^a+R_{\mu \mu }^bR_{\mu \nu }^b+2(R_{\mu \mu }^cR_{\mu \nu }^c)+R_{\mu \mu \mu }^d+R_{\mu \nu \nu }^dR_{\mu \mu \nu }^dR_{\mu \nu \mu }^d\right\},`$ (120)
$`\delta R`$ $`=`$ $`{\displaystyle \frac{16\pi ^2}{1w_0^2}}{\displaystyle \frac{1}{2}}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}\left\{R_{\mu \alpha \nu }^dR_{\mu \nu \alpha }^d\right\}=0,`$ (121)
with $`\mu \nu \alpha `$. We find that $`\delta R=0`$ from the symmetry of the integrand. The integrals are estimated by a mode sum for a periodic box of a size $`L^4`$ with $`L=64`$ after transforming the momentum variable through $`p_\mu =q_\mu \mathrm{sin}q_\mu `$. We choose $`\mathrm{\Lambda }=\pi `$ for the cut-off in the integrand $`\stackrel{~}{I}_\mu `$. The numerical results for $`A`$ are $`B`$ are presented in Table I as a function of $`M`$. We observe that the expected relation $`A=B`$ is well satisfied. We also checked that $`\mathrm{\Delta }R`$ is consistent with zero as expected. Finally we obtain
$`L_\mu (p)`$ $`=`$ $`{\displaystyle \frac{(1w_0^2)}{16\pi ^2}}{\displaystyle \frac{L^{\mathrm{latt}}}{3}}[p_\mu p/p^2\gamma _\mu ],`$ (122)
where
$`L^{\mathrm{latt}}`$ $`=`$ $`\mathrm{log}(\pi ^2/p^2)+{\displaystyle \frac{5}{6}}+3B.`$ (123)
Substituting the above expression for $`L_\mu (p)`$ in eqs.(57)$``$(60), we have
$`\left[\mathrm{\Lambda }_1^{(1p)}\right]_{ab;cd}`$ $`=`$ $`0,`$ (124)
$`\left[\mathrm{\Lambda }_2^{(1p)}\right]_{ab;cd}`$ $`=`$ $`{\displaystyle \frac{g^2}{12\pi ^2}}(1w_0^2)^4L^{\mathrm{latt}}\left[T^A\stackrel{~}{}T^A\right]_{ab;cd}\gamma _\mu (1\gamma _5)\gamma _\mu ,`$ (125)
$`\left[\mathrm{\Lambda }_i^{e(1p)}\right]_{ab;cd}`$ $`=`$ $`{\displaystyle \frac{g^2}{12\pi ^2}}(1w_0^2)^4{\displaystyle \underset{q}{}}\alpha _i^qL^{\mathrm{latt}}\left[T^A\stackrel{~}{}T^A\right]_{ab;cd}\gamma _\mu (1\gamma _5)\gamma _\mu ,`$ (126)
$`\left[\mathrm{\Lambda }_i^{o(1p)}\right]_{ab;cd}`$ $`=`$ $`\delta _{XY}{\displaystyle \frac{g^2}{12\pi ^2}}(1w_0^2)^4(\alpha _i^s+\alpha _i^d)L^{\mathrm{latt}}\left[T^A\stackrel{~}{}T^A\right]_{ab;cd}\gamma _\mu (1\gamma _5)\gamma _\mu ,`$ (127)
where we used the on-shell conditions for the external quark states to eliminate the terms proportional to $`p/`$. In the same way to derive eq.(56) in the continuum calculation, the lattice wave functions are written in the compact form as follows,
$`\mathrm{\Lambda }_i^{(1p)}`$ $`=`$ $`{\displaystyle \frac{g^2}{12\pi ^2}}{\displaystyle \frac{1}{4}}C(Q_i)L^{\mathrm{latt}}\left[\mathrm{\Lambda }_4^{(0)}+\mathrm{\Lambda }_6^{(0)}{\displaystyle \frac{1}{N}}(\mathrm{\Lambda }_3^{(0)}+\mathrm{\Lambda }_5^{(0)})\right],`$ (128)
where $`C(Q_i)`$ are already obtained in the previous subsection.
Comparing the continuum vertex functions in eq. (56) and lattice ones in eq. (128), we find that the penguin diagram contributions to the renormalization factor in eq.(1) are written as
$$𝒵_i^{\mathrm{pen}}Q_{\mathrm{pen}}^{\mathrm{latt}}=\frac{1}{(1w_0^2)^2}Z_i^{\mathrm{pen}}[Q_4+Q_6\frac{1}{N}(Q_3+Q_5)]^{\mathrm{latt}},$$
(129)
where the penguin operator and its coefficient are given by
$`Q_{\mathrm{pen}}^{\mathrm{latt}}`$ $`=`$ $`\left[Q_4+Q_6{\displaystyle \frac{1}{N}}\left(Q_3+Q_5\right)\right]^{\mathrm{latt}}`$ (130)
$`Z_i^{\mathrm{pen}}`$ $`=`$ $`{\displaystyle \frac{g^2}{12\pi ^2}}{\displaystyle \frac{C(Q_i)}{4}}(L^{\mathrm{latt}}L_i^{\mathrm{cont}})={\displaystyle \frac{g^2}{16\pi ^2}}{\displaystyle \frac{C(Q_i)}{3}}[\mathrm{log}(\mu a)^2+z_i^{\mathrm{pen}}]`$ (131)
with $`z_i^{\mathrm{pen}}=\mathrm{log}(\pi ^2)+3B{\displaystyle \frac{5}{6}}c_i`$. Numerical values of $`z_i^{\mathrm{pen}}`$ with the DRED scheme are presented in Table I as a function of $`M`$.
## III Full renormalization factors
In order to write down the complete expressions for the renormalization factors of the $`\mathrm{\Delta }S=1`$ operators $`Q_i`$, we still need to know the contributions from the gluon exchange diagrams, which is denoted by $`𝒵_{ij}^gQ_j^{\mathrm{latt}}`$ in eq. (1). Fortunately, they can be obtained from the combinations of the results in our previous paper, where we calculated the gluon exchange diagrams for the following four-quark operators,
$`𝒪_\pm (q_1,q_2,q_3,q_4)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left[(\overline{q}_1q_2)_{VA}(\overline{q}_3q_4)_{VA}\pm (\overline{q}_1q_4)_{VA}(\overline{q}_3q_2)_{VA}\right],`$ (132)
$`𝒪_1(q_1,q_2,q_3,q_4)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[C_F(\overline{q}_1q_2)_{VA}(\overline{q}_3q_4)_{V+A}+(\overline{q}_1T^Aq_2)_{VA}(\overline{q}_3T^Aq_4)_{V+A}\right],`$ (133)
$`𝒪_2(q_1,q_2,q_3,q_4)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[{\displaystyle \frac{1}{2N}}(\overline{q}_1q_2)_{VA}(\overline{q}_3q_4)_{V+A}+(\overline{q}_1T^Aq_2)_{VA}(\overline{q}_3T^Aq_4)_{V+A}\right],`$ (134)
with $`C_F=(N^21)/(2N)`$ the second Casimir of SU($`N`$) group. The tensor structures in the color space are $`1\stackrel{~}{}1`$ and $`T^A\stackrel{~}{}T^A`$.
The operators $`Q_i`$ are related to $`𝒪_{\pm ,1,2}`$ as
$`Q_2\pm Q_1`$ $`=`$ $`8𝒪_\pm (s,u,u,d),`$ (135)
$`Q_3\pm Q_4`$ $`=`$ $`8{\displaystyle \underset{q}{}}𝒪_\pm (s,d,q,q),`$ (136)
$`Q_9\pm Q_{10}`$ $`=`$ $`8{\displaystyle \underset{q}{}}{\displaystyle \frac{3e^q}{2}}𝒪_\pm (s,d,q,q),`$ (137)
$`Q_6`$ $`=`$ $`8{\displaystyle \underset{q}{}}𝒪_2(s,d,q,q),`$ (138)
$`Q_8`$ $`=`$ $`8{\displaystyle \underset{q}{}}{\displaystyle \frac{3e^q}{2}}𝒪_2(s,d,q,q),`$ (139)
$`NQ_5+Q_6`$ $`=`$ $`8{\displaystyle \underset{q}{}}𝒪_1(s,d,q,q),`$ (140)
$`NQ_7+Q_8`$ $`=`$ $`8{\displaystyle \underset{q}{}}{\displaystyle \frac{3e^q}{2}}𝒪_1(s,d,q,q).`$ (141)
From these relations the renormalized operators, which include the contributions of both the gluon exchange and penguin diagrams, are expressed in terms of the lattice operators as follows:
$`Q_i^{\mathrm{cont}}`$ $`=`$ $`{\displaystyle \frac{1}{(1w_0^2)^2Z_w^2}}\left[Z_{ij}^gQ_j^{\mathrm{latt}}+Z_i^{\mathrm{pen}}Q_{\mathrm{pen}}^{\mathrm{latt}}\right]`$ (142)
where
$`Z_{ii}^g`$ $`=`$ $`\{\begin{array}{cc}1+{\displaystyle \frac{g^2}{16\pi ^2}}[{\displaystyle \frac{3}{N}}\mathrm{log}(\mu a)^2+{\displaystyle \frac{z_++z_{}}{2}}]\hfill & i=1,2,3,4,9,10,\hfill \\ 1+{\displaystyle \frac{g^2}{16\pi ^2}}[{\displaystyle \frac{3}{N}}\mathrm{log}(\mu a)^2+z_1v_{21}]\hfill & i=5,7,\hfill \\ 1+{\displaystyle \frac{g^2}{16\pi ^2}}[{\displaystyle \frac{3(N^21)}{N}}\mathrm{log}(\mu a)^2+z_2+v_{21}]\hfill & i=6,8,\hfill \end{array}`$ (146)
$`Z_{ij}^g`$ $`=`$ $`\{\begin{array}{cc}=Z_{ji}^g={\displaystyle \frac{g^2}{16\pi ^2}}[3\mathrm{log}(\mu a)^2+{\displaystyle \frac{z_+z_{}}{2}}]\hfill & (i,j)=(1,2),(3,4),(9,10),\hfill \\ {\displaystyle \frac{g^2}{16\pi ^2}}[3\mathrm{log}(\mu a)^2+{\displaystyle \frac{z_2z_1+v_{21}v_{12}}{N}}]\hfill & (i,j)=(5,6),(7,8),\hfill \\ {\displaystyle \frac{g^2}{16\pi ^2}}Nv_{21}\hfill & (i,j)=(6,5),(8,7),\hfill \\ 0\hfill & \text{others}\hfill \end{array}`$ (151)
with $`z_{\pm ,1,2}`$, $`v_{12}`$, and $`v_{21}`$ presented in our previous paper. The penguin operator $`Q_{\mathrm{pen}}^{\mathrm{latt}}`$ and its coefficient $`Z_i^{\mathrm{pen}}`$ are given in eqs. (130) and (131), respectively. For the sake of convenience, we list the differences between the NDR and DRED schemes in the finite part of the renormalization constant,
$`\left({\displaystyle \frac{z_++z_{}}{2}}\right)^{\mathrm{NDR}}`$ $`=`$ $`\left({\displaystyle \frac{z_++z_{}}{2}}\right)^{\mathrm{DRED}}{\displaystyle \frac{N^26}{2N}},`$ (152)
$`(z_1v_{21})^{\mathrm{NDR}}`$ $`=`$ $`(z_1v_{21})^{\mathrm{DRED}}{\displaystyle \frac{N^28}{2N}},`$ (153)
$`(z_2+v_{21})^{\mathrm{NDR}}`$ $`=`$ $`(z_2+v_{21})^{\mathrm{DRED}}{\displaystyle \frac{N^24}{N}},`$ (154)
$`\left({\displaystyle \frac{z_+z_{}}{2}}\right)^{\mathrm{NDR}}`$ $`=`$ $`\left({\displaystyle \frac{z_+z_{}}{2}}\right)^{\mathrm{DRED}}{\displaystyle \frac{5}{2}},`$ (155)
$`\left({\displaystyle \frac{z_2z_1+v_{21}v_{12}}{N}}\right)^{\mathrm{NDR}}`$ $`=`$ $`\left({\displaystyle \frac{z_2z_1+v_{21}v_{12}}{N}}\right)^{\mathrm{DRED}}{\displaystyle \frac{7}{2}},`$ (156)
$`N(v_{21})^{\mathrm{NDR}}`$ $`=`$ $`N(v_{21})^{\mathrm{DRED}}3,`$ (157)
$`(z_i^{\mathrm{pen}})^{\mathrm{NDR}}`$ $`=`$ $`(z_i^{\mathrm{pen}})^{\mathrm{DRED}}+\{\begin{array}{cc}{\displaystyle \frac{5}{4}}\hfill & i=2,3,5,7,9,\hfill \\ {\displaystyle \frac{1}{4}}\hfill & i=4,6,8,10,\hfill \end{array}`$ (160)
where $`v_{12}=v_{21}=0`$ in the DRED scheme.
We are interested in the magnitude of the renormalization factors with the currently accessible coupling constant in numerical simulations. Let us estimate it taking $`\beta =6.0`$ with $`M=1.8`$ as a representative case. All the necessary ingredients in this analysis are given in our previous papers and Table I of this report. With the use of the mean field improvement, we have
$`(1w_0(\stackrel{~}{M})^2)Z_w^{\mathrm{MF}}(\stackrel{~}{M})=0.9085,`$ (161)
$`{\displaystyle \frac{z_+^{\mathrm{MF}}(\stackrel{~}{M})+z_{}^{\mathrm{MF}}(\stackrel{~}{M})}{2}}=12.848,`$ (162)
$`{\displaystyle \frac{z_+^{\mathrm{MF}}(\stackrel{~}{M})z_{}^{\mathrm{MF}}(\stackrel{~}{M})}{2}}=2.491,`$ (163)
$`z_1^{\mathrm{MF}}(\stackrel{~}{M})=11.187,`$ (164)
$`z_2^{\mathrm{MF}}(\stackrel{~}{M})=18.6,`$ (165)
$`z_2^{\mathrm{MF}}(\stackrel{~}{M})z_1^{\mathrm{MF}}(\stackrel{~}{M})=7.4,`$ (166)
$`z_i^{\mathrm{pen}}(\stackrel{~}{M})=2.328,`$ (167)
in the DRED scheme, where we employ
$`g_{\overline{MS}}^2(1/a)`$ $`=`$ $`\left[P{\displaystyle \frac{\beta }{6}}0.13486\right]^1=2.1792,`$ (168)
$`\stackrel{~}{M}`$ $`=`$ $`M+4(u1)=1.3112,`$ (169)
with $`P=u^4=0.59374`$ the plaquette value at $`\beta =6.0`$ in the quenched approximation. Combining these results we obtain
$`\left({\displaystyle \frac{Z_{ii}^g}{(1w_0^2)^2Z_w^2}}\right)^{\mathrm{MF}}`$ $`=`$ $`\{\begin{array}{cc}u^2\times 0.9968\hfill & i=1,2,3,4,9,10,\hfill \\ u^2\times 1.0245\hfill & i=5,7,\hfill \\ u^2\times 0.9006\hfill & i=6,8,\hfill \end{array}`$ (173)
$`\left({\displaystyle \frac{Z_{ij}^g}{(1w_0^2)^2Z_w^2}}\right)^{\mathrm{MF}}`$ $`=`$ $`\{\begin{array}{cc}=Z_{ji}^g=0.0416\hfill & (i,j)=(1,2),(3,4),(9,10),\hfill \\ 0.0412\hfill & (i,j)=(5,6),(7,8),\hfill \\ 0\hfill & \text{others}\hfill \end{array}`$ (177)
$`\left({\displaystyle \frac{Z_i^{\mathrm{pen}}}{(1w_0^2)^2Z_w^2}}\right)^{\mathrm{MF}}`$ $`=`$ $`C(Q_i)\times 0.0130,`$ (178)
at $`\mu a=1`$ in the DRED scheme, where we factor out the tadpole contributions. We find that the penguin diagram contributions to the renormalization factor are quite small.
## IV Conclusion
In this paper we have presented the full expressions for the renormalization factors of the $`\mathrm{\Delta }S=1`$ four-quark operators up to the one-loop level including the contributions of both the gluon exchange and penguin diagrams. Our calculation, however, does not include mixing with lower dimensional operators $`\overline{s}d`$ and $`\overline{s}\gamma _5d`$. The coefficients of these operators diverge as inverse powers of the lattice spacing, due to which it is practically inappropriate to subtract these lower dimensional operators by the perturbation theory. We are instead forced to use the non-perturbative methods such as those based on the chiral perturbation theory and the Schrödinger functional scheme. We leave it to future work.
## Acknowledgments
This work is supported in part by the Grants-in-Aid for Scientific Research from the Ministry of Education, Science and Culture (Nos. 12640253, 12014202). Y. K. is supported by Japan Society for Promotion of Science.
## Appendix
In this appendix we explain the domain-wall fermion action and its Feynman rules relevant for the present calculation. The domain-wall fermion action is written as,
$`S_{\mathrm{DW}}`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \underset{s=1}{\overset{N_s}{}}}[{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu }{}}(\overline{\psi }(n)_s(r+\gamma _\mu )U_\mu (n)\psi (n+\widehat{\mu })_s+\overline{\psi }(n)_s(r\gamma _\mu )U_\mu ^{}(n\widehat{\mu })\psi (n\widehat{\mu })_s)`$ (180)
$`+{\displaystyle \frac{1}{2}}(\overline{\psi }(n)_s(1+\gamma _5)\psi (n)_{s+1}+\overline{\psi }(n)_s(1\gamma _5)\psi (n)_{s1})+(M1+4r)\overline{\psi }(n)_s\psi (n)_s]`$
$`+`$ $`m{\displaystyle \underset{n}{}}\left(\overline{\psi }(n)_{N_s}P_R\psi (n)_1+\overline{\psi }(n)_1P_L\psi (n)_{N_s}\right),`$ (181)
where $`U_\mu `$ is the link variable of the SU($`N`$) gauge group and the Wilson parameter is set to $`r=1`$. Four dimensional space-time coordinate is labeled by $`n`$ and $`s`$ is an extra fifth dimensional index which runs from $`1`$ to $`N_s`$. Since we impose no gauge interaction along the fifth dimension, it is also possible to consider that the index $`s`$ labels the $`N_s`$ “flavor” space. In our one-loop calculation we take $`N_s\mathrm{}`$ to avoid complications arising from the finite $`N_s`$ such as mixing among the four-quark operators with different chiralities. The parameter $`m`$ denotes the physical quark mass and at $`m=0`$ one chiral zero mode is supposed to appear under the condition $`0<M<2`$ for the Dirac “mass” $`M`$. $`P_{R/L}`$ are projection operators defined by $`P_{R/L}=(1\pm \gamma _5)/2`$. For the gauge part we employ a standard four dimensional Wilson plaquette action.
The quark fields on the four-dimensional space-time are given by the combinations of the fermion fields at the boundaries,
$`q(n)=P_R\psi (n)_1+P_L\psi (n)_{N_s},`$ (182)
$`\overline{q}(n)=\overline{\psi }(n)_{N_s}P_R+\overline{\psi }(n)_1P_L.`$ (183)
We consider the composite operators constructed with these “physical” quark fields, and our renormalization procedure is based on the Green functions consisting of only these quark fields.
In order to obtain the Feynman rules we perform the weak coupling expansion of the quark and gauge actions. The fermion propagator with momentum $`p`$ is obtained by inverting the domain-wall Dirac operator in eq.(181), which is expressed by $`S_F(p)_{st}`$ as an $`N_s\times N_s`$ matrix in the “flavor” space. In the present one-loop calculation, however, we do not need the whole matrix elements because we consider the Green functions consisting of the physical quark fields. The relevant fermion propagators are restricted to following three types:
$`q(p)\overline{q}(p)`$ $`=`$ $`P_RS_F(p)_{11}P_L+P_LS_F(p)_{N_sN_s}P_R+P_LS_F(p)_{N_s1}P_L+P_RS_F(p)_{1N_s}P_R`$ (184)
$`=`$ $`{\displaystyle \frac{i\gamma _\mu \mathrm{sin}p_\mu +\left(1We^\alpha \right)m}{\left(1e^\alpha W\right)+m^2(1We^\alpha )}},`$ (185)
$`q(p)\overline{\psi }_s(p)`$ $`=`$ $`P_RS_F(p)_{1s}+P_LS_F(p)_{N_ss}`$ (186)
$`=`$ $`{\displaystyle \frac{1}{F}}\left(i\gamma _\mu \mathrm{sin}p_\mu m\left(1We^\alpha \right)\right)\left(e^{\alpha (N_ss)}P_R+e^{\alpha (s1)}P_L\right)`$ (188)
$`+{\displaystyle \frac{1}{F}}\left[m\left(i\gamma _\mu \mathrm{sin}p_\mu m\left(1We^\alpha \right)\right)F\right]e^\alpha \left(e^{\alpha (s1)}P_R+e^{\alpha (N_ss)}P_L\right),`$
$`\psi _s(p)\overline{q}(p)`$ $`=`$ $`S_F(p)_{s1}P_L+S_F(p)_{sN_s}P_R`$ (189)
$`=`$ $`{\displaystyle \frac{1}{F}}\left(e^{\alpha (N_ss)}P_L+e^{\alpha (s1)}P_R\right)\left(i\gamma _\mu \mathrm{sin}p_\mu m\left(1We^\alpha \right)\right)`$ (191)
$`+{\displaystyle \frac{1}{F}}\left(e^{\alpha (s1)}P_L+e^{\alpha (N_ss)}P_R\right)e^\alpha \left[m\left(i\gamma _\mu \mathrm{sin}p_\mu m\left(1We^\alpha \right)\right)F\right]`$
with
$`W`$ $`=`$ $`1Mr{\displaystyle \underset{\mu }{}}(1\mathrm{cos}p_\mu ),`$ (192)
$`\mathrm{cosh}(\alpha )`$ $`=`$ $`{\displaystyle \frac{1+W^2+_\mu \mathrm{sin}^2p_\mu }{2|W|}},`$ (193)
$`F`$ $`=`$ $`1e^\alpha Wm^2\left(1We^\alpha \right),`$ (194)
where the argument $`p`$ in the factors $`\alpha `$ and $`W`$ is suppressed.
In the perturbative calculation of Green functions we assume that the external quark momenta and masses are much smaller than the lattice cut-off. In this case the external quark propagators can be expanded in terms of them. We have the following expressions as the leading term of the expansion:
$`q\overline{q}(p)`$ $`=`$ $`{\displaystyle \frac{1w_0^2}{ip/+(1w_0^2)m}},`$ (195)
$`q\overline{\psi }_s(p)`$ $`=`$ $`q\overline{q}(p)\left(w_0^{s1}P_L+w_0^{N_ss}P_R\right),`$ (196)
$`\psi _s\overline{q}(p)`$ $`=`$ $`\left(w_0^{s1}P_R+w_0^{N_ss}P_L\right)q\overline{q}(p),`$ (197)
where $`w_0=1M`$. The form of $`q\overline{q}(p)`$ tells us that $`\sqrt{1w_0^2}`$ is the overlap factor to the four-dimensional quark fields at the tree-level.
The one-gluon fermion vertex with outgoing and incoming fermion momenta $`p`$ and $`q`$, respectively, is given by
$`V_{1\mu }^A(q,p)_{st}=gT^Av_\mu \left({\displaystyle \frac{q_\mu +p_\mu }{2}}\right)_{st}`$ $`=`$ $`gT^A\{i\gamma _\mu \mathrm{cos}\left({\displaystyle \frac{q_\mu +p_\mu }{2}}\right)+r\mathrm{sin}\left({\displaystyle \frac{q_\mu +p_\mu }{2}}\right)\}\delta _{st},`$ (198)
where $`T^A`$ $`(A=1,\mathrm{},N^21)`$ are generators of color SU($`N`$). We do not need two-gluon fermion vertex in the present calculation.
The gluon propagator for a gluon of momentum $`p`$ is written as
$`G_{\mu \nu }^{AB}(p)={\displaystyle \frac{\delta _{AB}}{4\mathrm{sin}^2(p/2)}}\left[\delta _{\mu \nu }(1\alpha ){\displaystyle \frac{4\mathrm{sin}(p_\mu /2)\mathrm{sin}(p_\nu /2)}{4\mathrm{sin}^2(p/2)}}\right],`$ (199)
where $`\mathrm{sin}^2(p/2)=_\mu \mathrm{sin}^2(p_\mu /2)`$ and we choose the Feynman gauge ($`\alpha =1`$) in our calculation.
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# Joint measurements via quantum cloning
## I Introduction
The first scheme for the joint measurement of noncommuting observables performed on a single quantum system was introduced by Arthurs-Kelly . The problem of evaluating the minimum added noise in the joint measurement of position and momentum, and more generally of a pair of observables whose commutator is not a c-number was then solved by Yuen . A similar approach to the problem has been followed in Ref. . In the case of two quadratures of one mode of the electromagnetic field the problem can be phrased in terms of a coherent POVM whose Naimark extension introduces an additional mode of the field. This kind of measurement can be realised by means of a heterodyne detector .
The case of the angular momentum of a quantum system is more difficult, and no measurement scheme has appeared in the literature so far. Spin coherent states can be introduced and interpreted as continuous (overcomplete) POVM, but the corresponding Naimark extension is unknown. It was shown that the spin coherent POVM minimises suitably defined quantities that represent the precision and the disturbance of the measurement , but explicit realisations of such a POVM are not known . The joint measurement of the three components $`J_x`$, $`J_y`$ and $`J_z`$ of the angular momentum could also be studied with discrete spectrum, rather than continuous. This problem does not have a solution yet. Joint measurements are a crucial ingredient in general quantum teleportation schemes , and are essential in connecting the quantum with the classical meaning of the angular momentum itself. Therefore, it is of great interest to find schemes that realise them.
The idea of this paper is to use quantum cloning to achieve joint measurements. It is well known that perfect cloning of unknown quantum systems is forbidden by the laws of quantum mechanics . The first universal cloning machine for spin 1/2 systems was proposed in Ref. , and later proved to be optimal in . More general universal transformations were then proposed in and proved to be optimal in Refs. (in Ref. the CP map of the optimal cloning transformations for finite dimensional systems was derived). However, the complete unitary transformation achieving the optimal cloning is not known (in Ref. some matrix elements of the unitary transformation are given for the case of qubits).
If we want to use quantum cloning to realise joint measurements, we may need to optimise it for a reduced covariance group, depending on the kind of the desired joint measurement. The cloning transformations mentioned above were optimised by imposing total covariance, i.e. for all possible unitary transformations. In general a restriction of the covariance group leads to a higher fidelity of the cloning transformation, as for example in the case of phase covariant cloning , where, however, only the bounds for the fidelity of the optimal cloning are given, but not the form of the optimal map.
For infinite dimensional systems it is not clear how to find the universal transformations for cloning. The extension to infinite dimension of the maps given in Ref. needs a regularization procedure, an example of which is given here in section IV. The infinite dimensional $`12`$ cloning machine proposed in is universal for coherent states, with resulting fidelity equal to $`2/3`$. In this paper we show that the cloning transformations proposed in are optimal for the joint measurement of orthogonal quadratures, and the joint measurement can be generalised to any angle between two noncommuting quadratures.
In the case of finite dimensional systems we will study the joint measurement of the three components of spin 1/2 states by operating a $`13`$ universal optimal cloning transformation on the original state and then performing independent measurements of $`\sigma _x`$, $`\sigma _y`$ and $`\sigma _z`$ on the three output copies. We will show that the resulting POVM is not optimal with respect to the added noise, but it gives a pretty good approximation of the optimal one.
The paper is organised as follows. In section II we consider the case of spin 1/2 systems. We first recall the optimal universal $`13`$ cloning transformations and then exploit them to achieve joint spin measurements. In section III we study the case of infinite dimensional systems, first reviewing the optimal $`12`$ transformation of Ref. and then applying it to the joint measurement of two quadratures of one mode of the electromagnetic field. In section IV we present a regularization of the map in Ref. in order to extend it to infinite dimensional Hilbert spaces, and show that the universal cloning does not achieve the optimal joint measurement. We summarise the results in section IV.
## II The finite dimensional case: joint spin measurements
In this section we analyse the case of spin 1/2 systems, by first reviewing the optimal universal $`13`$ cloning transformations which produce three output copies from a single input, and then exploiting this procedure to achieve joint measurements of the spin components. We will show that the joint spin measurement obtained in this way is only an approximation of the spin measurement POVM of Ref. .
### A Optimal $`13`$ cloning
We consider the case of universal cloning, namely transformations whose efficiency does not depend on the form of the input state. General $`NM`$ universal cloning transformations, which act on $`N`$ copies of a pure state $`|\psi `$ and produce $`M`$ output copies as close as possible to the input state, were proposed in Ref. and later proved to be optimal in Refs. . We consider here the form given in Ref. . The output state $`\rho _M`$ of the $`M`$ copies for spin 1/2 systems is given by
$`\rho _M={\displaystyle \frac{N+1}{M+1}}S_M\left(|\psi \psi |^N1𝐥^{(MN)}\right)S_M,`$ (1)
where $`S_M`$ is the the projection operator onto the symmetric subspace of the $`M`$ output copies. The fidelity $`F(N,M)=\psi |\text{Tr}_{M1}[\rho _M]|\psi `$ of each output copy with respect to the initial state $`|\psi `$ is given by
$`F(N,M)={\displaystyle \frac{M(N+1)+N}{M(N+2)}}.`$ (2)
Since the cloning transformation is universal it can be also viewed as a shrinking transformation of the Bloch vector of each copy, described by the shrinking factor $`\eta (N,M)`$ : the density operator describing the state of the $`M`$ output copies is given by $`\rho _{out}=\frac{1}{2}[1𝐥+\eta (N,M)\stackrel{}{s}_{in}\stackrel{}{\sigma }]`$, where $`\stackrel{}{s}_{in}`$ denotes the Bloch vector of the initial state $`|\psi `$ and $`\{\sigma _\alpha ,\alpha =x,y,z\}`$ are the Pauli operators. For the optimal transformations (1) the shrinking factor is
$`\eta (N,M)={\displaystyle \frac{N}{M}}{\displaystyle \frac{M+2}{N+2}}.`$ (3)
In the particular case of the $`13`$ cloning transformation, which we will consider in the following, the above map takes the form
$`\rho _3={\displaystyle \frac{1}{2}}S_3\left(|\psi \psi |1𝐥^2\right)S_3,`$ (4)
where $`S_3`$ is the projector on the space spanned by the vectors $`\{|s_is_i|,i=0÷3\}`$, with $`|s_0=|\mathrm{\hspace{0.17em}000}`$, $`|s_1=1/\sqrt{3}(|\mathrm{\hspace{0.17em}001}+|\mathrm{\hspace{0.17em}010}+|\mathrm{\hspace{0.17em}100})`$, $`|s_2=1/\sqrt{3}(|\mathrm{\hspace{0.17em}011}+|\mathrm{\hspace{0.17em}101}+|\mathrm{\hspace{0.17em}110})`$ and $`|s_3=|\mathrm{\hspace{0.17em}111}`$, where $`\{|\mathrm{\hspace{0.17em}0},|\mathrm{\hspace{0.17em}1}\}`$ is a basis for each spin 1/2 system. The value of the shrinking factor in this case is $`\eta (1,3)=5/9`$.
### B The joint spin measurement via cloning
We will now study a method to measure jointly the three components of a spin 1/2 system by first generating three approximate copies of the input state through an optimal $`13`$ cloning transformation, and then performing independent measurements on the three copies, namely measuring a different spin component on each copy. The POVM corresponding to the usual projection measurement of the $`\alpha `$-component of the Bloch vector on one copy is given by the operator $`[1𝐥+m_\alpha \sigma _\alpha ]/2`$, where $`\alpha =x,y,z`$ and $`m_\alpha =\pm 1`$ corresponds to the outcome of the measurement. The POVM $`\mathrm{\Omega }(\stackrel{}{m})`$ describing the measurement of the three components, each performed on a different copy, is then given by
$`\mathrm{\Omega }(\stackrel{}{m})={\displaystyle \frac{1}{8}}(1𝐥+m_x\sigma _x)(1𝐥+m_y\sigma _y)(1𝐥+m_z\sigma _z),`$ (5)
where the triplet $`\{m_x,m_y,m_z\}`$ represents the outcomes of the measurement. We will now consider the sequence of the $`13`$ cloning transformation followed by the measurement of a spin component on each of the three copies as a joint measurement on the initial input state of the original copy. In order to derive the corresponding POVM we first compute the probability distribution $`p(\stackrel{}{m})`$ as a function of the vector $`\stackrel{}{m}=\{m_x,m_y,m_z\}`$ of the outcomes
$`p(\stackrel{}{m})`$ $`=`$ $`\text{Tr}\left[\mathrm{\Omega }(\stackrel{}{m}){\displaystyle \frac{1}{2}}S_3\left(|\psi \psi |1𝐥^2\right)S_3\right]`$ (6)
$`=`$ $`\text{Tr}_1\left[|\psi \psi |{\displaystyle \frac{1}{2}}\text{Tr}_{2,3}[S_3\mathrm{\Omega }(\stackrel{}{m})S_3]\right],`$ (7)
where Tr<sub>i</sub> denotes the partial trace over the $`i`$th clone. This measurement, viewed as a joint measurement on the original copy $`|\psi \psi |`$ can then be described in terms of the POVM $`\mathrm{\Pi }(\stackrel{}{m})`$
$`\mathrm{\Pi }(\stackrel{}{m})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\text{Tr}_{2,3}\left[S_3\mathrm{\Omega }(\stackrel{}{m})S_3\right].`$ (8)
A lengthy and straightforward matrix algebra gives the resulting POVM in the simple form
$`\mathrm{\Pi }(\stackrel{}{m})={\displaystyle \frac{1}{8}}\left[1𝐥+{\displaystyle \frac{5}{9}}\stackrel{}{m}\stackrel{}{\sigma }\right].`$ (9)
Notice that the $`5/9`$ factor in front of the Pauli operators corresponds to the shrinking factor of the optimal $`13`$ cloning transformation. This can be intuitively expected because the average value of the spin components of the three cloned copies that are measured is shrunk by this factor.
We will now compute the accuracy of this joint measurement in order to have a comparison with the coherent POVM given in Ref. . As mentioned above, the POVM (9) leads to the following rescaling between the measured average value $`\sigma _\alpha _m`$ and the theoretical one for all the three spin components
$`\sigma _\alpha _m={\displaystyle \underset{\stackrel{}{m}}{}}m_\alpha \text{Tr}[|\psi \psi |\mathrm{\Pi }(\stackrel{}{m})]={\displaystyle \frac{5}{9}}\psi |\sigma _\alpha |\psi .`$ (10)
Therefore, the unbiased estimate $`\sigma _\alpha _e`$ for the spin components corresponds to rescaling the measured outcome variables to $`m_\alpha =\pm 9/5`$, such that
$`\sigma _\alpha _e={\displaystyle \frac{9}{5}}\sigma _\alpha _m`$ (11)
and the second moment is also rescaled as follows
$`\mathrm{\Delta }\sigma _\alpha ^2_e={\displaystyle \frac{81}{25}}\mathrm{\Delta }\sigma _\alpha ^2_m.`$ (12)
In order to study the uncertainty of this measurement we compute the sum of the variances corresponding to the three spin components $`J_\alpha =\sigma _\alpha /2`$. Since $`\sigma _\alpha ^2_m=1`$ for all the components, the uncertainty in the estimate is given by
$`\mathrm{\Delta }J^2_e`$ $`=`$ $`{\displaystyle \underset{\alpha =x,y,z}{}}J_\alpha ^2_eJ_\alpha _e^2`$ (13)
$`=`$ $`{\displaystyle \frac{1}{4}}\left(3{\displaystyle \frac{81}{25}}1\right)={\displaystyle \frac{109}{50}}.`$ (14)
We will now compute the corresponding accuracy for the coherent measurement . The coherent POVM is given by the projection onto spin coherent states $`|𝐧𝐧|`$ , where $`𝐧=(\mathrm{sin}\theta \mathrm{cos}\varphi ,\mathrm{sin}\theta \mathrm{sin}\varphi ,\mathrm{cos}\theta )`$ is a unit vector and
$`𝐧𝐉|𝐧=j|𝐧.`$ (15)
Let us calculate as an example the uncertainty related to the component $`J_z`$. Since the measurement is unbiased, the measured mean values of the spin components $`J_\alpha _m`$ coincide with the theoretical mean values, and we don’t need to introduce rescaling factors as we did in the previous case. The estimated values therefore coincide with the measured ones. For the component $`J_z`$ one has
$`J_z_m={\displaystyle 𝑑\mu (𝐧)(j+1)\mathrm{cos}\theta |𝐧𝐧|},`$ (16)
where $`d\mu (𝐧)=d𝐧(2j+1)/4\pi `$. The measured mean value of $`J_z^2`$ is given by
$`J_z^2_m={\displaystyle 𝑑\mu (𝐧)(j+1)^2\mathrm{cos}^2\theta |𝐧𝐧|},`$ (17)
that can be written as
$`J_z^2_m={\displaystyle \frac{j+1}{j+3/2}}\left[\psi |J_z^2|\psi +{\displaystyle \frac{1}{2}}(j+1)\right].`$ (18)
The measured mean values related to the $`x`$ and $`y`$ components can be calculated analogously and one has the same relation as Eq. (18) for all components $`\alpha =x,y,z`$. The total uncertainty in the spin measurement then takes the form
$`\mathrm{\Delta }J^2_e`$ $`=`$ $`{\displaystyle \frac{j(j+1)^2}{j+3/2}}+3{\displaystyle \frac{(j+1)^2}{2j+3}}{\displaystyle \underset{\alpha =x,y,z}{}}J_\alpha _e^2`$ (19)
$``$ $`2j+1,`$ (20)
where for $`j=1/2`$ and pure states the bound is achieved, and is equal to 2. This value has to be compared with Eq. (14), obtained by three measurements on the three cloned copies. As we can see, the joint measurement via cloning does not achieve the minimum added noise as the optimal POVM of Eqs. (15) and (16), however it provides a good approximation. Notice that the minimum added noise would be achieved by a discrete POVM of the form $`\mathrm{\Pi }(\stackrel{}{m})=\frac{1}{8}[1𝐥+\stackrel{}{m}\stackrel{}{\sigma }]`$.
## III The infinite dimensional case: joint quadrature measurements
In this section we study the cloning for infinite dimensional systems proposed in Ref. . We review the optimal $`12`$ transformation and then apply it to the joint measurement of two quadratures of one mode of the electromagnetic field. We will show that the cloning transformation is optimal for joint measurements of orthogonal quadratures, and the joint measurement can be generalised to any angle between two noncommuting quadratures by suitably changing the state of the ancilla.
### A Optimal $`12`$ cloning
For the following, it is convenient to introduce the formalism of heterodyne eigenvectors. Consider the heterodyne-current operator $`Z=a+b^{}`$, which satisfies the commutation relation $`[Z,Z^{}]=0`$ and the eigenvalue equation $`Z|z_{ab}=z|z_{ab}`$, with $`z`$. The eigenstates $`|z_{ab}`$ are given by
$`|z_{ab}D_a(z)|0_{ab}=D_b(z^{})|0_{ab},`$ (21)
where $`D_d(z)=e^{zd^{}z^{}d}`$ denotes the displacement operator for mode $`d`$ and $`|0_{ab}\frac{1}{\sqrt{\pi }}_{n=0}^{\mathrm{}}()^n|n_a|n_b`$. The eigenstates $`|z_{ab}`$ are a complete orthogonal set with Dirac-normalization $`{}_{ab}{}^{}z|z^{}_{ab}^{}=\delta ^{(2)}(zz^{})`$, $`\delta ^{(2)}(z)`$ denoting the delta function over the complex plane. For $`z=0`$ the state $`|0_{ab}`$ can be approximated by a physical (normalizable) state, corresponding to the output of a non-degenerate optical parametric amplifier (NOPA)—so-called twin beam—in the limit of infinite gain at the NOPA .
It is also useful to evaluate the expression $`{}_{cb}{}^{}z|z^{}_{ab}^{}`$ which is given by
$`{}_{cb}{}^{}z|z^{}_{ab}^{}={\displaystyle \frac{1}{\pi }}D_a(z^{})𝒯_{ac}D_c^{}(z),`$ (22)
where $`𝒯_{ac}=_n|n_a_cn|`$ denotes the transfer operator satisfying the relation $`𝒯_{ac}|\psi _c=|\psi _a`$ for any state $`|\psi `$. In the following we transpose the main results of the continuous variable cloning of Ref. , according to the formalism just introduced. The input state at the cloning machine can be written
$`|\varphi =|\phi _c{\displaystyle _{}}d^2zf(z,z^{})|z_{ab}`$ (23)
where $`|\phi _c`$ is the initial state to be cloned, belonging to the Hilbert space $`_c`$, whereas $`_a`$ is the Hilbert space pertaining to the cloned state, and $`_b`$ is an ancillary Hilbert space. We do not specify for the moment the explicit form of the function $`f(z,z^{})`$. The cloning transformation is realized by the unitary operator
$`U`$ $`=`$ $`\mathrm{exp}\left[c(a^{}+b)c^{}(a+b^{})\right]`$ (24)
$`=`$ $`\mathrm{exp}\left[2i\left(Y_c\text{Re}ZX_c\text{Im}Z\right)\right]`$ (25)
with $`X_c,Y_c`$ denoting the conjugated quadratures for mode $`c`$, namely $`X_c=(c+c^{})/2`$ and $`Y_c=(cc^{})/2i`$.
The unitary evolution in Eq. (25) can be approached experimentally by means of a network of three NOPA’s under suitable gain conditions . Notice the simple relation $`U|z_{ab}=D_c^{}(z)|z_{ab}`$. The state after the cloning transformation is given by
$`|\varphi _{out}=U|\varphi ={\displaystyle _{}}d^2zf(z,z^{})D_c^{}(z)|\phi _c|z_{ab}.`$ (26)
Let us evaluate the one-mode restricted density matrix $`\varrho _c`$ and $`\varrho _a`$ corresponding to the state $`|\varphi _{out}`$, for the Hilbert spaces $`_c`$ and $`_a`$ supporting the two clones. For $`\varrho _c`$ one has
$`\varrho _c`$ $`=`$ $`\text{Tr}_{ab}[|\varphi _{out}\varphi _{out}|]`$ (27)
$`=`$ $`{\displaystyle _{}}d^2w{\displaystyle _{}}d^2z{\displaystyle _{}}d^2z^{}f(z,z^{})f^{}(z^{},z^{})\times `$ (29)
$`{}_{ab}{}^{}w|D_c^{}(z)|\phi _c|z_{ab}_c\phi |D_c(z^{}){}_{ab}{}^{}z^{}|w_{ab}^{}`$
$`=`$ $`{\displaystyle _{}}d^2z|f(z,z^{})|^2D_c^{}(z)|\phi _c_c\phi |D_c(z),`$ (30)
where we have evaluated the trace by using the completeness and the orthogonality relation of the eigenstates $`|w_{ab}`$. For $`\varrho _a`$, using Eq. (22), one has
$`\varrho _a`$ $`=`$ $`\text{Tr}_{cb}[|\varphi _{out}\varphi _{out}|]`$ (31)
$`=`$ $`{\displaystyle _{}}d^2w{\displaystyle _{}}d^2z{\displaystyle _{}}d^2z^{}f(z,z^{})f^{}(z^{},z^{})\times `$ (33)
$`{}_{cb}{}^{}w|D_c^{}(z)|\phi _c|z_{ab}_c\phi |D_c(z^{}){}_{ab}{}^{}z^{}|w_{cb}^{}`$
$`=`$ $`{\displaystyle _{}}d^2w{\displaystyle _{}}{\displaystyle \frac{d^2z}{\pi }}{\displaystyle _{}}{\displaystyle \frac{d^2z^{}}{\pi }}f(z,z^{})f^{}(z^{},z^{})`$ (34)
$`\times `$ $`D_a(z)𝒯_{ac}[D_c^{}(w)D_c^{}(z)|\phi _c{}_{c}{}^{}\phi |D_c(z^{})D_c(w)]`$ (35)
$`\times `$ $`𝒯_{ca}D_a^{}(z^{})`$ (36)
$`=`$ $`{\displaystyle _{}}d^2w|\stackrel{}{f}(w,w^{})|^2D_a^{}(w)|\phi _a_a\phi |D_a(w),`$ (37)
where $`\stackrel{}{f}(w,w^{})`$ denotes the Fourier transform over the complex plane
$`\stackrel{}{f}(w,w^{})={\displaystyle _{}}{\displaystyle \frac{d^2z}{\pi }}e^{wz^{}w^{}z}f(z,z^{}).`$ (38)
Hence, for $`f(z,z^{})=\stackrel{}{f}(z,z^{})`$ one has $`\varrho _c=\varrho _a`$, namely the two clones are identical. In the following we will show that the choice of the function $`f(z,z^{})`$ determines a criterium of optimality in terms of joint measurement of noncommuting quadratures of the original system through the measurement of separate (commuting) quadratures over the two clones.
### B The joint quadrature measurement via cloning
Quantum cloning allows one to engineer new joint measurements of a quantum system, by suitably measuring the cloned copies. In the case of $`12`$ copies just introduced, measuring two quadratures on the two clones is equivalent to the joint measurement of conjugated quadratures on the original, similarly to a heterodyne measurement. Consider the simplest case
$`f(z,z^{})=\sqrt{{\displaystyle \frac{2}{\pi }}}\mathrm{exp}(|z|^2)`$ (39)
in Eqs. (23), (30) and (37). One obtains $`\varrho _c=\varrho _a`$, namely the two clones are identical, and their state is given by the original state $`|\phi `$ degraded by Gaussian noise. The state preparation $`|\chi `$ pertaining to the Hilbert space $`_a_b`$ is given explicitly by
$`|\chi `$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle _{}}d^2ze^{|z|^2}|z_{ab}`$ (40)
$`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle _{}}d^2ze^{|z|^2}D_a(z)|0_{ab}`$ (41)
$`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle _{}}d^2ze^{\frac{3}{2}|z|^2}{\displaystyle \underset{n,m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!m!}}\times `$ (43)
$`z^n(z^{})^ma^na^n|0_{ab}`$
$`=`$ $`\sqrt{2\pi }{\displaystyle \frac{2}{3}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}({\displaystyle \frac{2}{3}})^na^na^n|0_{ab}`$ (44)
$`=`$ $`\sqrt{2\pi }{\displaystyle \frac{2}{3}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{2}{3}})^n{\displaystyle \frac{1}{n!}}{\displaystyle \frac{(a^{}a)!}{(a^{}an)!}}|0_{ab}`$ (45)
$`=`$ $`\sqrt{2\pi }{\displaystyle \frac{2}{3}}\left({\displaystyle \frac{1}{3}}\right)^{a^{}a}|0_{ab}`$ (46)
$`=`$ $`{\displaystyle \frac{2\sqrt{2}}{3}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{1}{3}}\right)^n|n_a|n_b`$ (47)
$`=`$ $`e^{\text{atanh}\frac{1}{3}(aba^{}b^{})}|\mathrm{\hspace{0.17em}0}.`$ (48)
One recognizes in Eq. (48) the twin-beam state at the output of a NOPA with total number of photons $`N=\chi |a^{}a+b^{}b|\chi =1/4`$, corresponding to a gain $`G=9/8`$ . More generally, notice that
$`\sqrt{{\displaystyle \frac{2}{\pi \mathrm{\Delta }^2}}}{\displaystyle _{}}d^2ze^{\mathrm{\Delta }^2|z|^2}|z`$ (49)
$`=`$ $`\sqrt{1\lambda ^2}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(\lambda )^n|n|n`$ (50)
$`=`$ $`e^{\text{atanh}\lambda (aba^{}b^{})}|\mathrm{\hspace{0.17em}0},`$ (51)
with $`\lambda =(\mathrm{\Delta }^21/2)/(\mathrm{\Delta }^2+1/2)`$.
Now let us evaluate the entangled state $`\varrho `$ at the output of the cloning machine. After tracing over the ancillary mode $`b`$, one has
$`\varrho `$ $`=`$ $`\text{Tr}_b[|\varphi _{out}\varphi _{out}|]`$ (52)
$`=`$ $`{\displaystyle \frac{1}{2}}P_{c,a}(|\phi _c{}_{c}{}^{}\phi |𝟙_𝕒)_{𝕔,𝕒},`$ (53)
where $`P_{c,a}`$ is the projector given by
$`P_{c,a}`$ $`=`$ $`{\displaystyle _{}}d^2z{\displaystyle \frac{2}{\pi }}e^{|z|^2}D_c^{}(z)D_a(z)`$ (54)
$`=`$ $`V\left({\displaystyle _{}}{\displaystyle \frac{d^2z}{\pi }}e^{\frac{1}{2}|z|^2}D_c^{}(z)𝟙_a\right)V^{}`$ (55)
$`=`$ $`V\left({\displaystyle _{}}{\displaystyle \frac{d^2z}{\pi }}e^{zc^{}}e^{z^{}c}𝟙_a\right)V^{}`$ (56)
$`=`$ $`V(|\mathrm{\hspace{0.17em}0}_c{}_{c}{}^{}0|𝟙_a)V^{},`$ (57)
with $`V=\mathrm{exp}[\frac{\pi }{4}(c^{}aca^{})]`$ that realizes the unitary transformation
$`V\left(\begin{array}{c}c\\ a\end{array}\right)V^{}={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}1\hfill & \hfill 1\\ 1\hfill & \hfill 1\end{array}\right)\left(\begin{array}{c}c\\ a\end{array}\right).`$ (64)
In the last line of Eq. (57) a derivation similar to Eq. (48) has been followed. Measuring the quadratures $`X_c`$ and $`Y_a`$ over the two clones is then equivalent to perform a measurement on the original state $`|\phi _c`$, with the measurement described by the following POVM
$`F(x,y)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\text{Tr}_a[P_{c,a}|x_c{}_{c}{}^{}x||y_{a}^{}{}_{a}{}^{}y|P_{c,a}],`$ (65)
where $`|x_c`$ and $`|y_a`$ denote the eigenstates of $`X_c`$ and $`Y_a`$, respectively. From the following relations
$`V^{}|x_c_cx||y_a_ay|V`$ (66)
$`=2|\sqrt{2}(xiy)_{ca}_{ca}\sqrt{2}(xiy)|,`$ (67)
$`{}_{c}{}^{}0|z_{ca}={\displaystyle \frac{1}{\sqrt{\pi }}}|z^{}_a,`$ (68)
$`V|\alpha _c|\beta _a=|(\alpha +\beta )/\sqrt{2}_c|(\beta \alpha )/\sqrt{2}_a,`$ (69)
\[in Eqs. (68) and (69) the single-mode states denote coherent states\] one obtains
$`F(x,y)={\displaystyle \frac{1}{\pi }}|x+iy_c_cx+iy|,`$ (70)
namely the coherent-state POVM, which is the well-known optimal POVM for the joint measurement of the conjugated quadratures $`X_c`$ and $`Y_c`$. In fact, from Eqs. (26), (30) and (37), one has the following relations between the quantum expectation values $`\varphi _{out}|\mathrm{}|\varphi _{out}`$ over the output state $`|\varphi _{out}`$ of Eq. (26) with respect to the values $`\phi |\mathrm{}|\phi `$ over the original input state
$`\varphi _{out}|g(c,c^{})|\varphi _{out}=\text{Tr}_c[\varrho _cg(c,c^{})]`$ (71)
$`=`$ $`{\displaystyle _{}}d^2z|f(z,z^{})|^2\phi |g(cz,c^{}z^{})|\phi ,`$ (72)
$`\varphi _{out}|g(a,a^{})|\varphi _{out}=\text{Tr}_a[\varrho _ag(a,a^{})]`$ (73)
$`=`$ $`{\displaystyle _{}}d^2z|\stackrel{}{f}(z,z^{})|^2\phi |g(cz,c^{}z^{})|\phi ,`$ (74)
which holds for any function $`g`$. In particular, for $`f(z,z^{})`$ given by Eq. (39), one has
$`\varphi _{out}|\mathrm{\Delta }X_c^2|\varphi _{out}=\phi |\mathrm{\Delta }X_c^2|\phi +{\displaystyle \frac{1}{4}},`$ (75)
$`\varphi _{out}|\mathrm{\Delta }Y_a^2|\varphi _{out}=\phi |\mathrm{\Delta }Y_c^2|\phi +{\displaystyle \frac{1}{4}},`$ (76)
namely one achieves the simultaneous measurement of conjugated quadratures over the input state with minimum added noise , thus proving the optimality of the joint measurement.
The condition in order to obtain identical clones $`f(z,z^{})=\stackrel{}{f}(z,z^{})`$ can be satisfied also by a bivariate Gaussian of the form
$`f(z,z^{})=\sqrt{{\displaystyle \frac{2}{\pi }}}\mathrm{exp}\left({\displaystyle \frac{\text{Re}^2z}{\sigma ^2}}\sigma ^2\text{Im}^2z\right).`$ (77)
In the following we will show that in such case the cloning trasformation becomes optimal for the joint measurement of noncommuting quadratures at angles which depend on the parameter $`\sigma `$ in Eq. (77). In fact, Eq. (53) is replaced by
$`\varrho ={\displaystyle \frac{1}{2}}P_{c,a}(\sigma )(|\phi _c{}_{c}{}^{}\phi |𝟙_𝕒)_{𝕔,𝕒}(\sigma ),`$ (78)
where the projector $`P_{c,a}(\sigma )`$ is evaluated as follows
$`P_{c,a}(\sigma )={\displaystyle _{}}d^2z{\displaystyle \frac{2}{\pi }}\mathrm{exp}({\displaystyle \frac{\text{Re}^2z}{\sigma ^2}}\sigma ^2\text{Im}^2z)\times `$ (80)
$`D_c^{}(z)D_a(z)`$
$`=`$ $`V\left[{\displaystyle _{}}{\displaystyle \frac{d^2z}{\pi }}\mathrm{exp}\left({\displaystyle \frac{\text{Re}^2z}{2\sigma ^2}}{\displaystyle \frac{\sigma ^2\text{Im}^2z}{2}}\right)D_c^{}(z)𝟙_a\right]V^{}`$ (81)
$`=`$ $`VS_c(\mathrm{ln}\sigma )\left({\displaystyle _{}}{\displaystyle \frac{d^2z}{\pi }}e^{\frac{1}{2}|z|^2}D_c^{}(z)𝟙_a\right)S_c^{}(\mathrm{ln}\sigma )V^{}`$ (82)
$`=`$ $`VS_c(\mathrm{ln}\sigma )(|\mathrm{\hspace{0.17em}0}_c{}_{c}{}^{}0|𝟙_a)S_c^{}(\mathrm{ln}\sigma )V^{}`$ (83)
$`=`$ $`S_c(\mathrm{ln}\sigma )S_a(\mathrm{ln}\sigma )V(|\mathrm{\hspace{0.17em}0}_c{}_{c}{}^{}0|𝟙_a)V^{}\times `$ (85)
$`S_c^{}(\mathrm{ln}\sigma )S_a^{}(\mathrm{ln}\sigma )`$
$`=`$ $`S_c(\mathrm{ln}\sigma )S_a(\mathrm{ln}\sigma )P_{c,a}S_c^{}(\mathrm{ln}\sigma )S_a^{}(\mathrm{ln}\sigma ),`$ (86)
with $`S_d(r)=\mathrm{exp}[r(d^2d^2)/2]`$ denoting the squeezing operator for mode $`d`$ that realizes the unitary transformation
$`S_d^{}(r)dS_d(r)=(\mathrm{cosh}r)d+(\mathrm{sinh}r)d^{}.`$ (87)
As in Eq. (65), one can evaluate the POVM that is obtained upon measuring the quadratures $`X_c`$ and $`Y_a`$ over the clones. From the relations for the quadrature projectors
$`S_c^{}(\mathrm{ln}\sigma )|x_c_cx|S_c(\mathrm{ln}\sigma )={\displaystyle \frac{1}{\sigma }}|x/\sigma _c_cx/\sigma |,`$ (88)
$`S_a^{}(\mathrm{ln}\sigma )|y_a_ay|S_a(\mathrm{ln}\sigma )=\sigma |x\sigma _c_cx\sigma |,`$ (89)
and Eqs. (67), (68) and (69), one has
$`F_\sigma (x,y)={\displaystyle \frac{1}{2}}\text{Tr}_a[P_{c,a}(\sigma )|x_c{}_{c}{}^{}x||y_{a}^{}{}_{a}{}^{}y|P_{c,a}(\sigma )]`$ (90)
$`=`$ $`{\displaystyle \frac{1}{\pi }}S_c(\mathrm{ln}\sigma )|{\displaystyle \frac{x}{\sigma }}+i\sigma y_c{\displaystyle \frac{}{}}_c{\displaystyle \frac{x}{\sigma }}+i\sigma y|S_c^{}(\mathrm{ln}\sigma )`$ (91)
$`=`$ $`{\displaystyle \frac{1}{\pi }}D_c(x+iy)S_c(\mathrm{ln}\sigma )|\mathrm{\hspace{0.17em}0}_c_c0|S_c^{}(\mathrm{ln}\sigma )D_c^{}(x+iy).`$ (92)
Eq. (91) shows that the POVM is formally a squeezed state. Such kind of POVM is optimal for the joint measurement of the two noncommuting quadrature operators $`X_\varphi ,X_\varphi `$, with $`\varphi =\text{arctg}(\sigma ^2)`$. In fact, one has the relations
$`{\displaystyle 𝑑x𝑑y(x\mathrm{cos}\varphi \pm y\mathrm{sin}\varphi )F_\sigma (x,y)}=X_{\pm \varphi },`$ (94)
$`{\displaystyle 𝑑x𝑑y(x\mathrm{cos}\varphi \pm y\mathrm{sin}\varphi )^2F_\sigma (x,y)}`$
$`=`$ $`X_{\pm \varphi }^2+{\displaystyle \frac{1}{4}}\left|\mathrm{sin}(2\varphi )\right|=X_{\pm \varphi }^2+{\displaystyle \frac{1}{2}}\left|[X_\varphi ,X_\varphi ]\right|,`$ (95)
namely the outcomes $`x\mathrm{cos}\varphi \pm y\mathrm{sin}\varphi `$ trace the expectation values of the observables $`X_{\pm \varphi }`$ respectively, with minimum added noise .
## IV Regularization of the universally covariant cloning
In this section we give a procedure to extend the completely positive (CP) map for the universal cloning of Werner’s paper in the case of infinite dimensional Hilbert space. The procedure is based on a suitable regularization in order to achieve a trace-preserving map. In particular, we will show that the universal $`12`$ cloning does not provide a tool to obtain the joint measurement of noncommuting observables. Hence, we prove that Werner-type cloning and the cloning of Ref. used in the previous section are different, and they are optimal for different purposes.
We rewrite here the CP map for $`NM`$ cloning given in Ref.
$`T(\varrho )={\displaystyle \frac{d[N]}{d[M]}}S_M(\varrho 𝟙^{(𝕄)})𝕊_𝕄,`$ (96)
where $`d[N]=\left(\genfrac{}{}{0pt}{}{d+N1}{N}\right)`$, $`d`$ being the dimension of a single-copy Hilbert space; $`S_M`$ is the projector on the symmetric subspace, as mentioned in Sect. II; and $`\varrho =|\psi \psi |^N`$ is the initial state of $`N`$ identical copies in the state $`|\psi \psi |`$. The projector $`S_M`$ can be written in terms of two-site permutation operators $`\mathrm{\Pi }_{(ij)}`$ (transposition), by using recursively the relation
$`S_M={\displaystyle \frac{1}{M}}\left(𝟙+{\displaystyle \underset{𝕚=\mathrm{𝟙}}{\overset{𝕄\mathrm{𝟙}}{}}}\mathbb{\Pi }_{(𝕚𝕄)}\right)S_{M1}.`$ (97)
The permutation operator $`\mathrm{\Pi }_{(ij)}`$ can be expressed one the Hilbert space $`_i_j`$ as follows
$`\mathrm{\Pi }_{(ij)}={\displaystyle \underset{n}{}}A_nA_n^{},`$ (98)
where $`\{A_n\}`$ are a generic set of operators satisfying the completeness relation
$`B={\displaystyle \underset{n}{}}\text{Tr}[A_n^{}B]A_n.`$ (99)
For example, in the case of $`12`$ cloning for spin $`1/2`$ one has
$`S_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}(𝟙𝟙+{\displaystyle \frac{\mathrm{𝟙}}{\mathrm{𝟚}}}{\displaystyle \underset{𝕚=\mathrm{𝟘}}{\overset{\mathrm{𝟛}}{}}}\sigma _𝕚\sigma _𝕚)`$ (100)
$`=`$ $`{\displaystyle \frac{3}{4}}𝟙𝟙+{\displaystyle \frac{\mathrm{𝟙}}{\mathrm{𝟜}}}{\displaystyle \underset{𝕚=\mathrm{𝟙}}{\overset{\mathrm{𝟛}}{}}}\sigma _𝕚\sigma _𝕚`$ (101)
where $`\sigma _0=𝟙`$ and $`\sigma _i`$ ($`=1,2,3`$) are the customary Pauli matrices.
The map in Eq. (96) can be formally extended to infinite dimensional Hilbert space upon using the transposition operator
$`\stackrel{~}{\mathrm{\Pi }}_{(ij)}={\displaystyle \frac{d^2\alpha }{\pi }D_i(\alpha )D_j^{}(\alpha )},`$ (102)
however the trace-preserving condition on physical CP maps imposes to replace the identity operator in Eq. (96) with a normalizable state. Here we suggest a regularization of $`12`$ cloning in $`_c_a`$ by using Eq. (102) along with the normalizable (thermal) state $`\lambda ^{aa}`$, and then we write
$`\stackrel{~}{T}(\varrho )`$ $`=`$ $`K\stackrel{~}{S}_2\left(\varrho \lambda ^{a^{}a}\right)\stackrel{~}{S}_2,`$ (103)
where $`K`$ is a constant and
$`\stackrel{~}{S}_2={\displaystyle \frac{1}{2}}\left(𝟙_𝕔𝟙_𝕒+\stackrel{~}{\mathbb{\Pi }}_{(𝕔𝕒)}\right).`$ (104)
From the identities
$`\text{Tr}_a[\stackrel{~}{\mathrm{\Pi }}_{(ca)}]=𝟙_𝕔,\text{Tr}_𝕔[\stackrel{~}{\mathbb{\Pi }}_{(𝕔𝕒)}]=𝟙_𝕒,`$ (105)
$`\stackrel{~}{\mathrm{\Pi }}_{(ca)}(AB)=(BA)\stackrel{~}{\mathrm{\Pi }}_{(ca)},`$ (106)
and the trace-preserving condition $`\text{Tr}\stackrel{~}{T}(\varrho )=1`$, one obtains the expression for the factor $`K`$
$`K=2\left\{\text{Tr}\left[(𝟙+\varrho )\lambda ^{𝕔^{}𝕔}\right]\right\}^1.`$ (107)
Notice that the dependence of $`K`$ on $`\varrho `$ makes the transformation in Eq. (103) nonlinear, however such a nonlinear character is vanishing for $`\lambda 1`$. The regularization indeed consists in taking the limit $`\lambda 1`$. In this case the one-site restricted density matrix is given by
$`\text{Tr}_1[\stackrel{~}{T}(\varrho )]=\text{Tr}_2[\stackrel{~}{T}(\varrho )]`$ (108)
$`={\displaystyle \frac{1}{2}}\left(\varrho +{\displaystyle \frac{\lambda ^{c^{}c}}{\text{Tr}[\lambda ^{a^{}a}]}}\right),\lambda 1,`$ (109)
which generalizes the customary depolarizing Pauli channel to the infinite dimensional case.
In the following we will show that, differently from the cloning of section III, our regularization of Werner-type cloning does not allow one to achieve the optimal joint measurement of conjugated quadratures. In fact, similarly to Eq. (65), one can evaluate the POVM that corresponds to separate quadrature measurements over the two clones as follows
$`G(x,y)=\text{Tr}_c[K\lambda ^{a^{}a}\stackrel{~}{S}_2|x_c{}_{c}{}^{}x||y_{a}^{}{}_{a}{}^{}y|\stackrel{~}{S}_2].`$ (110)
Asymptotically, in the limit $`\lambda 1`$, one rewrites
$`G(x,y){\displaystyle \frac{1\lambda }{2}}\times `$ (111)
$`({}_{a}{}^{}y|\lambda ^{a^{}a}|y_{a}^{}|x_c{}_{c}{}^{}x|+{}_{a}{}^{}x|\lambda ^{a^{}a}|x_{a}^{}|y_{c}^{}{}_{c}{}^{}y|+`$ (112)
$`{}_{a}{}^{}x|\lambda ^{a^{}a}|y_{a}^{}|x_c_cy|+{}_{a}{}^{}y|\lambda ^{a^{}a}|x_{a}^{}|y_c_cx|).`$ (113)
Notice that one has
$`{\displaystyle 𝑑x𝑑yxG(x,y)}=`$ (115)
$`{\displaystyle \frac{1}{2}}X_c+{\displaystyle \frac{1\lambda }{2}}\left(\text{Tr}[X_a\lambda ^{a^{}a}]+\lambda ^{c^{}c}X_c+X_c\lambda ^{c^{}c}\right)`$
$``$ $`{\displaystyle \frac{1}{2}}X_c`$ (118)
$`{\displaystyle 𝑑x𝑑yx^2G(x,y)}=`$
$`{\displaystyle \frac{1}{2}}X_c^2+{\displaystyle \frac{1\lambda }{2}}\left(\text{Tr}[X_a^2\lambda ^{a^{}a}]+\lambda ^{c^{}c}X_c^2+X_c^2\lambda ^{c^{}c}\right)`$
$``$ $`{\displaystyle \frac{1}{2}}X_c^2+{\displaystyle \frac{1}{8}}\left(1+{\displaystyle \frac{2\lambda }{1\lambda }}\right),`$ (119)
and analogous expressions for integration on $`y`$. Hence, the average values of the variables $`x`$ and $`y`$ provide the expectation values of the quadratures $`X_c`$ and $`Y_c`$ (apart from the scaling factor $`1/2`$, similar to the shrinking factor of section II). However, one can see that the statistical error for such variables diverges for $`\lambda 1`$ since the the second moment goes to infinity.
The symmetrizer in Eq. (104) can be rewritten as follows
$`\stackrel{~}{S}_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}V[𝟙_𝕔𝟙_𝕒+{\displaystyle \frac{𝕕^\mathrm{𝟚}\alpha }{\pi }𝔻_𝕔(\sqrt{\mathrm{𝟚}}\alpha )𝟙_𝕒}]𝕍^{}`$ (120)
$`=`$ $`{\displaystyle \frac{1}{2}}V[𝟙_𝕔𝟙_𝕒+()^{𝕔^{}𝕔}𝟙_𝕒]𝕍^{}`$ (121)
$`=`$ $`V[{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}|2n_c{}_{c}{}^{}2n|𝟙_𝕒]V^{}.`$ (122)
This expression can be more easily compared with the projector of Eq. (57) that achieves the cloning transformation for the optimal joint measurement. The different action of the two projectors is clear on the basis of coherent states. One has
$`\stackrel{~}{S}_2|\alpha _c|\beta _a|\alpha _c|\beta _a+|\beta _c|\alpha _a`$ (123)
$`P_{c,a}|\alpha _c|\beta _a|(\alpha +\beta )/2_c|(\alpha +\beta )/2_a,`$ (124)
hence the operator $`P_{c,a}`$ indeed projects on a space that is smaller than the symmetric subspace. In fact the cloning map $`𝒯(\varrho )=\frac{1}{2}P_{ca}(\sigma )(\varrho 𝟙_𝕒)_{𝕔𝕒}(\sigma )`$ is not universally covariant, but is covariant only under the group of unitary displacement operators, namely
$`𝒯\left(D(\alpha )\varrho D^{}(\alpha )\right)=D(\alpha )^2𝒯(\varrho )D^{}(\alpha )^2.`$ (125)
## V Conclusions
In this paper we have investigated the possibility of achieving joint measurements of noncommuting observables on a single quantum system by means of quantum cloning.
We have shown that the universally covariant cloning is not optimal for joint measurements, and a suitable non covariant cloning is needed. Different measures of quality should be used for quantum cloning, depending on what final use is to be made of the output copies. This is also indicated by recent studies of different copying machines for information transfer . If we want to use quantum cloning to realise joint measurements, we need to optimise it for a suitable reduced covariance group, depending on the kind of the desired joint measurement. For spin 1/2—a finite dimensional example—the universal cloning optimised by imposing total covariance exhibit added noise in the joint measurement of the spin components. This shows that for finite dimensional systems the completely covariant cloning is not optimised for joint measurements, but in order to achieve optimal joint measurements the cloning transformations should be optimised with some ad hoc procedure. Also in the infinite dimensional case, the suitably regularized universal covariant cloning map does not allow to achieve the ideal joint measurement of noncommuting observables.
A restriction of the covariance group in general leads to a higher fidelity of the cloning transformation, as in the case of phase covariant cloning or for the cloning map of Ref. . The last case indeed provides a tool to perform the ideal joint measurement, as we have shown in section III.
Regarding the experimental feasibility of the schemes of measurement presented in this paper, we want to stress that a way to implement the universal cloning transformations was proposed in , with clones as indistinguishable photons, and the final measurement of the three spin components on the three output copies would correspond to nonlinear observables of radiation, whose measurement is not currently feasible. On the contrary, the infinite dimensional case is much more realistic, since the $`12`$ cloning transformation considered in section III can be achieved experimentally by means of a sequence of parametric amplifiers , and the quadrature measurements are obtained by customary homodyne detectors.
This work was supported in part by the European Union project EQUIP (contract IST-1999-11053) and by Ministero dell’Università e della Ricerca Scientifica e Tecnologica under the project “Quantum information transmission and processing: quantum teleportation and error correction”. One of us (M.F.S.) is also acknowledging support from ESF-QIT program.
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# GREEN’S FUNCTIONS IN AXIAL AND LORENTZ-TYPE GAUGES AND APPLICATION TO THE AXIAL POLE PRESCRIPTION AND THE WILSON LOOP *talk presented at the workshop QFT2000,held at Saha Institute of Nuclear Physics , Calcutta,INDIA from 18 Jan 2000 till 22 Jan 2000. , email address:sdj@iitk.ac.in
## 1 MOTIVATIONS FOR THE WORK
This work summarizes the progress made recently in applying the finite field-dependent BRS \[FFBRS\] transformations to connect the axial and the Lorentz type gauges and the resultant applications to the problems associated with the axial gauges.These transformations preserve the vacuum expectation of gauge-invariant observables by their explicit construction.The summary is based on the works \[1-4\] with Aalok Misra and on and is based on the earlier works .
The known high energy physics is well represented by the standard model ;which is a non-abelian gauge theory. Hence practical calculations (and their methodology) in nonabelian gauge theories assume importance in Particle Physics. Practical calculations in a non-abelian gauge theory require a choice of a gauge and there are many choices available. \[Theories involving gravity also has a gauge-invariance and require a gauge choice.\] There are many distinct \[families of\] gauges, variously useful in different situations. For example, we have:
Lorentz Gauges:
Coulomb Gauge:
Axial Gauges: \[Includes Light-cone,Temporal etc.\]
Planar Gauges:
radial Gauges:
Quadratic Gauges:
in SBGT R-xi Gauges;etc.
The uses of these gauges in a variety of different contexts have been elaborated in reference 8. A priori, we expect from gauge invariance,that the values for physical observables calculated in different gauges are identical. Formal proofs of such equivalence for the S-matrix elements has been given in a given class of gauges, say the Lorentz-type gauges with a variable gauge parameter $`\lambda `$ .Some isolated attempts to connect S-matrix elements in singular \[rather than a class of them\] gauges also have been done. For example, formal equivalence of S-matrix elements in the Coulomb and the Landau gauges \[both singular gauges\] has also been established .Similar formal attempts to connect the \[singular\] temporal gauge with Feynman gauge in the canonical formalism has also been done . It is important to note that, however, the Green’s functions in the gauges such as the Coulomb ,the axial,the planar and the light-cone in the path integral formulation are ambiguous on account of the unphysical singularities in their propagators. Hence, it becomes important to know how to define the Green’s functions in such gauges in such a manner that they are compatible to those in a well-defined covariant gauge such as the Lorentz gauge. A general procedure that connects Green’s functions in the path integral formulation in two classes of gauges, say the Lorentz to the axial,has been lacking until recently. Such comparisons are important not just in a formal sense but also in practice.Precisely because of this, the proper treatment of the 1/$`(\eta .q)^p`$ type poles in axial and light-cone gauges (and also similar questions in the Coulomb gauge ) has occupied a lot of attention and the criterion used for their validation has, in fact, been the comparison with the calculational results in the Lorentz gauges.Such comparative calculations, where possible have to be done by brute force and have been done to O\[g<sup>4</sup>\] generally; thus limiting the scope of their confirmation. At a time, a physically observable anomalous dimension was reported to differ in Lorentz and axial gauges . Such questions motivate us to develop a general path integral formalism that can address all these questions in a wide class of gauges in a single framework. In a purely Feynman diagrammatic approach,we ,of course, have the attempt of Cheng and Tsai\[14 \].
Further, gauges such as axial gauges suffer from the prescription ambiguity for the 1/(n.q)<sup>p</sup> type singularities. Several ad hoc prescriptions have been given; of these the Principal value prescription\[PVP\] and the Leibbrandt-Mandelstam\[M-L\] are the ones used extensively in literature.While much work has been done on these,they however have run inevitably into difficulties of various kinds .Such ambiguities do not exist in Lorentz type gauges. We thus expect that a field transformation from the Lorentz gauges to axial gauges will enable us to derive the correct prescription for such singularities.
## 2 THE PROCEDURE ADOPTED
Such considerations motivate us to seek a field transformation that relates different family of gauges. In this section,we shall outline in brief the procedure we adopt in these works.The details of the procedure will be left to the next section.The procedure consists of :
* Construct a field transformation that transforms the path integral from one gauge choice to another.
* Develop method to correlate the Green’s functions in two sets of gauges formally.
* Establish a relation between propagators. This is expected to give an understanding of how to deal with 1/n.q -type singularities in the axial gauge propagator.\[The prescription problem.\]
We know explicitly the infinitesimal gauge transformation that relates two gauge functions differing infinitesimally; viz.
F<sup>α</sup> –\> F<sup>α</sup> \+ $`\delta `$F<sup>α</sup>
It is given by
$`\delta `$A<sup>α</sup><sub>μ</sub> = D$`{}_{\mu }{}^{}{}_{}{}^{\alpha \beta }`$ $`M_{\beta \gamma }^1\delta `$F<sup>γ</sup>
This is an example of a field-dependent infinitesimal gauge transformation.It is difficult to integrate and explicitly evaluate the finite version of this.For example,even in a simple case such as seeking the transformation from $``$.A = 0 \[Landau\] to A<sub>0</sub>’= 0\[temporal\] gauge,the transformation can be formally solved for ,but is difficult to evaluate explicitly. Here,A’ is given by
A’ = $`\frac{1}{g}`$U\[A\] \[ $``$+gA\]U\[A\]<sup>-1</sup>
With
U\[A\] = T {exp $``$$`{}_{}{}^{t}{}_{_0}{}^{}`$ A<sub>0</sub>(x,t’) dt’}
For such reasons, we instead seek an alternate approach in which we try to integrate an analogue of the BRS transformations.This seems to be more easily manageable. This property arises mainly from the facts that BRS transformations are nilpotent, The “finite” BRS and “infinitesimal” BRS have the same form unlike the gauge transformations<sup>1</sup><sup>1</sup>1 In the context of the usual BRS transformations,where “$`\delta \mathrm{\Lambda }`$” and “$`\mathrm{\Lambda }`$” are merely *constants,*this remark may look trivial because $`\mathrm{\Lambda }^2`$=0;however in connection with the FFBRS transformations,this distinction is not trivial.See the remark below (3.3). .
## 3 FFBRS TRANSFORMATIONS:A GENERALIZATION OF BRS
We note the invariance of the Faddeev-Popov Effective Action \[FPEA\] in a gauge ’F’ is not dependent on the anticommuting parameter “$`\delta \mathrm{\Lambda }`$” being infinitesimal nor on whether it is field-independent or not .In fact transformations such as these:
A’(x) = A(x) + D<sub>μ</sub>c(x) $`\mathrm{\Theta }`$\[A(y),c(y),$`\overline{c}`$ (y)\] (3.1a)
c<sup>α</sup>’(x) = c<sup>α</sup>(x) – 1/2 g f$`{}_{}{}^{\alpha \beta \gamma }c_{}^{\beta }c^\gamma `$ $`\mathrm{\Theta }`$\[A(y),c(y),$`\overline{c}`$ (y)\] (3.1b)
$`\overline{c}`$<sup>α</sup>(x) = $`\overline{c}`$<sup>α</sup>(x) \+ $`\frac{A^\alpha }{\lambda }`$ $`\mathrm{\Theta }`$\[A(y),c(y),$`\overline{c}`$ (y)\] (3.1c)
or in brief,
$`\varphi _i^{}`$=$`\varphi _i`$ +$`\delta _{i,BRS}[\varphi `$\]$`\mathrm{\Theta }`$\[A(y),c(y),$`\overline{c}`$ (y)\] (3.2)
where $`\mathrm{\Theta }`$\[A(y),c(y),$`\overline{c}`$ (y)\] is a field-dependent functional but does not depend on ’x’ are also symmetries of the FPEA.$`\mathrm{\Theta }`$ need not be “infinitesimal”,in that the Green’s functions of $`\mathrm{\Theta }`$ with appropriate external lines can be a finite number even when $`\mathrm{\Theta }`$<sup>2</sup>=0.For examples of these, see e.g. Ref.8 and some of these will also appear in applications below.
It has been shown that in some special cases of these,the corresponding FFBRS can be obtained by integration of an infinitesimal field-dependent BRS \[IFBRS\] transformation:
d$`\varphi `$(x,$`\kappa `$)/d$`\kappa `$ =$`\delta _{BRS}`$ \[$`\varphi `$(x,$`\kappa `$))\]$`\mathrm{\Theta }^{}`$ \[$`\varphi `$(y,$`\kappa `$)\] (3.3)
which stand for the three equations:
dA(x,$`\kappa `$)/d$`\kappa `$=D<sub>μ</sub>c(x,$`\kappa `$) $`\mathrm{\Theta }`$’\[A(y,$`\kappa `$),c(y,$`\kappa `$),$`\overline{c}`$ (y,$`\kappa `$)\] (3.4a)
dc(x,$`\kappa `$)/d$`\kappa `$=- 1/2 g f$`{}_{}{}^{\alpha \beta \gamma }c_{}^{\beta }(x,\kappa )c^\gamma `$ (x,$`\kappa )`$$`\mathrm{\Theta }`$’\[A(y,$`\kappa `$),c(y,$`\kappa `$),$`\overline{c}`$ (y,$`\kappa `$)\] (3.4b)
d$`\overline{c}`$(x,$`\kappa `$)/d$`\kappa `$= $``$$``$A<sup>α</sup>(x,$`\kappa )`$/$`\lambda `$ $`\mathrm{\Theta }`$’\[A(y,$`\kappa `$),c(y,$`\kappa `$),$`\overline{c}`$ (y,$`\kappa `$)\] (3.4c)
In this case, the integral can be evaluated in a closed form.When integrated it, in fact, maintains its BRS form on account of the two properties of BRS mentioned earlier.{For a more transparent derivation,you may also refer to }.This is much simpler than and unlike a gauge transformation where a finite gauge transformation does not retain the simple form of an infinitesimal gauge transformation.The result,thus,exhibits the FFBRS structure of (3.1), with $`\mathrm{\Theta }`$ given by:
$`\mathrm{\Theta }`$\[$`\varphi `$\] = $`\mathrm{\Theta }`$’\[$`\varphi `$\] {exp(f\[$`\varphi `$\]) - 1}/ f\[$`\varphi `$\] (3.5)
with
f =$``$$`__i`$ \[$`\delta `$$`\mathrm{\Theta }`$’/d$`\varphi _i`$\] $`\delta _{i,BRS}`$ \[$`\varphi `$(x,$`\kappa `$))\] (3.6)
The Jacobian for the above non-local transformation is difficult to evaluate directly.However,it turns out that in many cases of interest, the Jacobian for such FFBRS can effectively be cast in the form exp {iS<sub>1</sub>}with a local S<sub>1</sub>;which in fact leads to a total effective action {S<sub>eff</sub>+S<sub>1</sub>} that is an FPEA for another gauge theory .Thus \[in such cases\], the FFBRS converts one path-integral<sup>2</sup><sup>2</sup>2 for a more accurate statement, please see reference 6. into another rather than one action to another;the action being in fact invariant under under the FFBRS.To be explicit,if we express D$`\varphi `$=D$`\varphi `$’J\[$`\varphi `$’\].Then ,in such cases, for any gauge-invariant observable G\[$`\varphi ]`$,we have
\<\<G\[$`\varphi `$\]\>\><sub>L</sub>= $``$D$`\varphi `$ G\[$`\varphi `$\]exp{ i S<sub>eff</sub>$`{}_{}{}^{L}[`$$`\varphi `$$`]`$}=$``$D$`\varphi `$’J\[ $`\varphi `$’\]G\[$`\varphi ^{}`$\]exp{ i S<sub>eff</sub>$`{}_{}{}^{L}[`$$`\varphi ^{}`$$`]`$}
$``$$``$D$`\varphi `$’G\[$`\varphi ^{}`$\]exp{ i S<sub>eff</sub>$`{}_{}{}^{L}[`$$`\varphi ^{}`$$`]`$+iS$`{}_{1}{}^{}[`$$`\varphi ^{}`$\]}=$``$D$`\varphi `$’G\[$`\varphi ^{}`$\]exp{ i S<sub>eff</sub>$`{}_{}{}^{A}[`$$`\varphi ^{}`$$`]`$} (3.7)
Here, S<sub>eff</sub><sup>A</sup> is the effective action of another gauge such as the axial;
In particular,for G\[$`\varphi ]`$= I,we have that $`\varphi `$$``$$`\varphi `$’ transforms vacuum to vacuum amplitude in the two gauges.
We thus outline a general prescription for constructing an FFBRS connecting any two families of Yang-Mills effective actions:
* \[a\] Establish a continuous route of interpolating gauges \[if necessary\] from one family to another;
* \[b\]Postulate an infinitesimal field-dependent BRS transformation.The form of the infinitesimal gauge transformation (where available) serves as a preliminary hint.
* \[c\] Using the form for the interpolating S, if necessary, guess a form for S<sub>1</sub>\[$`\varphi `$,$`\kappa `$\].
* \[d\] Evaluate the Jacobian for an infinitesimal BRS in step \[b\].This is easily evaluated compared to that for an FFBRS.
* \[e\]Impose the condition meant for J$``$exp{ iS<sub>1</sub> }. This condition leads to constraints on $`\mathrm{\Theta }`$’ and the coefficients and the form for S<sub>1</sub> .This condition reads :
#### \<\<$`\frac{i}{J}`$$`\frac{dJ}{d\kappa }`$+$`\frac{dS_1[\varphi (\kappa ),\kappa ]}{d\kappa }`$\><sub>κ</sub>$``$0 (3.8)
#### and involves the Jacobian for the infinitesimal transformation of (3.4)which is easy to evaluate.
### We note that while the above procedure is required while dealing with arbitrary effective actions for Yang-Mills theories, in the special case of arbitrary two Faddeev-Popov Effective Actions \[FPEA\], a simpler proof has also recently been given. As examples of the former class, not included in the latter, we note, for clarity, that the latter class does not include BRS-antiBRS effective action of Baulieu and Thierry-Mieg as well as say the BRS invariant action in planar gauges. For such cases we need to follow the entire process as presented .
## 4 THE RESULTS FOR THE FFBRS FOR LORENTZ $``$AXIAL
We follow the procedure as outlined in section 3.\[Alternately we can also use the result of ref.17 ,now available.\]
First we understand the n.A=0 gauge as the $`\lambda 0`$ limit of the gauge with
S<sub>gf</sub> = -$`\frac{1}{2\lambda }`$ d$`{}_{}{}^{4}x`$ \[n.A\]<sup>2</sup> (4.1)
( \[together with the corresponding ghost term\].
Then we construct an intermediate gauge-fixing term
S<sub>gf</sub> = -$`\frac{1}{2\lambda }`$ d$`{}_{}{}^{4}x`$ \[(1- $`\kappa `$)$``$A +$`\kappa `$ n.A\]<sup>2</sup> (4.2)
together with the corresponding ghost term. From these,we make an ansatz for $`\mathrm{\Theta }`$’ and S<sub>1</sub>; and impose the Jacobian condition.We then obtain as one possible solution the following:
$`\mathrm{\Theta }`$’\[$`\varphi `$($`\kappa `$)\] =$`id^4y`$ $`\overline{c}`$<sup>γ</sup>(y,$`\kappa `$) \[$`.`$A<sup>γ</sup>(y,$`\kappa )`$ -$`\eta `$.A<sup>γ</sup>(y,$`\kappa )]`$ (4.3)
Then,the FFBRS transformation that takes one from the Lotentz-type gauges to the axial-type gauges is given by (3.1),with $`\mathrm{\Theta }`$ given by
$`\mathrm{\Theta }`$\[$`\varphi `$\] = $`\mathrm{\Theta }`$’\[$`\varphi `$\] {exp(f\[$`\varphi `$\]) - 1}/ f (4.4)
and
f\[$`\varphi `$\]=$``$$`__i`$ \[$`\delta `$$`\mathrm{\Theta }`$’\[$`\varphi `$\] /d$`\varphi _i`$\] $`\delta _{i,BRS}`$ \[$`\varphi `$(x)\] (4.5a)
=$`id^4y`$ {$`\overline{c}`$(y)\[$``$-$`\eta `$\] D<sub>μ</sub>c(y) + $`.`$A<sup>γ</sup>(y$`)`$\[$`.`$A<sup>γ</sup>(y) -$`\eta `$.A<sup>γ</sup>(y)\]$`\frac{1}{\lambda }`$} (4.5b)
We mention some of the properties of the results:
\[A\] While the results look complicated, a simple formula that enables one to evaluate the Green’s functions in one set of gauges in terms of the Feynman rules in another set of gauges is possible as seen in the next section.
\[B\]Similar procedure can be followed for other pairs of gauges not otherwise connected.\[See references 6 and 17\].
\[C\]The transformations preserve the vacuum expectation values of gauge invariant operators automatically.e.g.The Wilson Loop \[See section 8\].
## 5 THE RESULTS FOR GREEN’S FUNCTIONS FOR THE TWO GAUGES.
Unlike the path integral and the vacuum expectation value of gauge invariant observables, the relation between Green’s functions in the two sets of gauges requires another, but related transformation . This transformation turns out to be an integral of the IFBRS :
dA(x,$`\kappa `$)/d$`\kappa `$=D<sub>μ</sub>c(x,$`\kappa `$) $`\mathrm{\Theta }`$’\[A(y,$`\kappa `$),c(y,$`\kappa `$),$`\overline{c}`$ (y,$`\kappa `$)\] (5.1a)
dc(x,$`\kappa `$)/d$`\kappa `$=- 1/2 g f$`{}_{}{}^{\alpha \beta \gamma }c_{}^{\beta }(x,\kappa )c^\gamma (x,\kappa `$) $`\mathrm{\Theta }`$’\[A(y,$`\kappa `$),c(y,$`\kappa `$),$`\overline{c}`$ (y,$`\kappa `$)\] (5.1b)
d$`\overline{c}`$(x,$`\kappa `$)/d$`\kappa `$= $`[`$$``$A$`{}_{}{}^{\alpha }(1\kappa )+\kappa \eta .A`$$`{}_{}{}^{\alpha }]`$/$`\lambda `$ $`\mathrm{\Theta }^{}`$\[A(y,$`\kappa `$),c(y,$`\kappa `$),$`\overline{c}`$ (y,$`\kappa `$)\] (5.1c)
We can re-express these in a single equation:
d$`\varphi `$<sub>i</sub>(x,$`\kappa `$)/d$`\kappa `$ $``$$`(\stackrel{~}{\delta _{1i}[\varphi }`$\]+$`\kappa `$$`(\stackrel{~}{\delta _{2i}[\varphi }`$\])$`\mathrm{\Theta }^{}`$\[A(y,$`\kappa `$),c(y,$`\kappa `$),$`\overline{c}`$ (y,$`\kappa `$)\] (5.2)
Let us say that the integral of the above from $`\kappa `$=0 to $`\kappa `$=1 reads
$`\varphi ^{}`$$``$$`\varphi (x,\kappa =1)`$ $`\varphi (x,\kappa =0)`$$`+\mathrm{\Delta }\varphi [`$$`\varphi ]=`$$`\varphi +\mathrm{\Delta }\varphi `$ (5.3)
It is the above solution (5.3) that appears in the relation between the Green’s functions between two gauges.In fact, defining for a gauge ’F’
\<\<O\[$`\varphi `$\]\>\><sub>F</sub>= $``$D$`\varphi `$ O\[$`\varphi `$\]exp{ i S<sub>eff</sub>$`{}_{}{}^{F}[`$$`\varphi `$$`]`$} (5.4)
The result of this is that
\<\<O\[$`\varphi `$\]\>\><sub>A</sub> = \<\< O \[$`\varphi +\mathrm{\Delta }\varphi `$\]\>\><sub>L</sub> (5.5)
$`\mathrm{\Delta }\varphi `$ can in fact be given through the relations :
$`\mathrm{\Delta }\varphi `$ = {$`\stackrel{~}{\delta _1}`$$`[\varphi ]\mathrm{\Theta }_1[\varphi ]`$$`+\stackrel{~}{\delta _2}`$$`[\varphi ]\mathrm{\Theta }_2[\varphi ]`$}$`\mathrm{\Theta }^{}[\varphi ]`$ (5.6)
where as defined in ,
$`\mathrm{\Theta }_{1,2}[\varphi ]`$=$``$$`__0`$<sup>1</sup>d$`\kappa `$(1,$`\kappa `$) exp {$`\kappa f_1[\varphi ]+\frac{\kappa ^2}{2}f_2[\varphi ]\}`$ (5.7)
The result again seems complicated, in principle, to be of use in evaluating Green’s functions; it can, however, be put in tractable form as below <sup>3</sup><sup>3</sup>3 See Equation (5.4) for definitions of $`\stackrel{\textcolor[rgb]{1,0,1}{~}}{\textcolor[rgb]{1,0,1}{\delta }_{\textcolor[rgb]{1,0,1}{1}\textcolor[rgb]{1,0,1}{i}}\textcolor[rgb]{1,0,1}{[}\textcolor[rgb]{1,0,1}{\varphi }}`$\] and$`\stackrel{\textcolor[rgb]{1,0,1}{~}}{\textcolor[rgb]{1,0,1}{\delta }_{\textcolor[rgb]{1,0,1}{2}\textcolor[rgb]{1,0,1}{i}}\textcolor[rgb]{1,0,1}{[}\textcolor[rgb]{1,0,1}{\varphi }}`$\]. :
\<\<O\[$`\varphi `$\]\>\><sub>A</sub>’= \<\<O\[$`\varphi `$\]\>\><sub>L</sub> +i$``$$`__0`$<sup>1</sup>d$`\kappa `$ D$`\varphi `$ exp{ i S<sub>eff</sub>$`{}_{}{}^{M}[`$$`\varphi `$,$`\kappa ]`$}
$``$$`_i(\stackrel{~}{\delta _{1i}[\varphi }`$\]+$`\kappa `$$`(\stackrel{~}{\delta _{2i}[\varphi }`$\])(-i$`\mathrm{\Theta }^{})`$$`\frac{\delta ^LO}{\delta \varphi _i}`$ (5.8)
The proper definition of the Green’s functions in Lorentz gauges requires that we include a term -i$`ϵ`$$``$d$`{}_{}{}^{4}x`$($`\frac{1}{2}`$AA -$`\overline{c}`$ c) in the effective action. Similarly,proper definition of the axial Green’s functions in Axial-type gauges requires that we include an appropriate $`ϵ`$-term.The correct $`ϵ`$-term in axial-type gauges can be obtained by the FFBRS transformation and the relation of the form (5.5)applied appropriately . It can be shown that the effect of this term on the axial gauge Green’s functions is expressed in the simplest form when it is expressed in the relation (5.8) and happens to be simply to add to S<sub>eff</sub>$`{}_{}{}^{M}[`$$`\varphi `$,$`\kappa ]`$ the same -i$`ϵ`$$``$d$`{}_{}{}^{4}x`$($`\frac{1}{2}`$AA -$`\overline{c}`$ c) inside the $`\kappa `$-integral.Thus, taking care of the proper definition of the axial Green’s functions ,these Green’s functions ,compatible with those in Lorentz gauges ,are given by
\<\<O\[$`\varphi `$\]\>\><sub>A</sub>’= \<\<O\[$`\varphi `$\]\>\><sub>L</sub> +i$``$$`__0`$<sup>1</sup>d$`\kappa `$ D$`\varphi `$ exp{ i S<sub>eff</sub>$`{}_{}{}^{M}[`$$`\varphi `$,$`\kappa ]`$ +$`ϵ`$$``$d$`{}_{}{}^{4}x`$($`\frac{1}{2}`$AA -$`\overline{c}`$ c)}
$``$$`_i(\stackrel{~}{\delta _{1i}[\varphi }`$\]+$`\kappa `$$`(\stackrel{~}{\delta _{2i}[\varphi }`$\])(-i$`\mathrm{\Theta }^{})`$$`\frac{\delta ^LO}{\delta \varphi _i}`$ (5.9)
As seen above, a given Green’s function to a given finite order can be evaluated by means of a finite set of diagrams with vertices from the Lorentz gauges and BRS variations AND the propagators from the mixed gauges.A $`\kappa `$-integral is also required to be performed.
An example:
Consider O\[$`\varphi `$\]= A$`{}_{\mu }{}^{\alpha }(x)`$A<sup>β</sup>$`{}_{\nu }{}^{}(y).`$Then, \<\<A$`{}_{\mu }{}^{\alpha }(x)`$A<sup>β</sup>$`{}_{\nu }{}^{}(y)`$\>\><sub>A</sub> = i G<sup>Aαβ</sup><sub>μν</sub>(x-y) for the connected part of the axial gauge propagator.Then, in obvious notations, (5.9) reads,
iG<sup>Aαβ</sup><sub>μν</sub>(x-y)=iG<sup>Lαβ</sup><sub>μν</sub>(x-y)+i$``$$`__0`$<sup>1</sup>d$`\kappa `$ D$`\varphi `$ exp{ i S<sub>eff</sub>$`{}_{}{}^{M}[`$$`\varphi `$,$`\kappa ]`$ +$`ϵ`$$``$d$`{}_{}{}^{4}x`$($`\frac{1}{2}`$AA -$`\overline{c}`$ c)}
$`[`$D<sub>μ</sub>c<sup>α</sup>(x)A<sup>β</sup>$`{}_{\nu }{}^{}(y)`$+A$`{}_{\mu }{}^{\alpha }(x)`$D<sub>μ</sub>c<sup>β</sup>(y)\] $`d^4z`$ $`\overline{c}`$<sup>γ</sup>(z) \[$`.`$A<sup>γ</sup>($`z)`$ -$`\eta `$.A<sup>γ</sup>(z)\] (5.10)
The above relation gives the value of the exact axial propagator compatible to the Green’s functions in Lorentz gauges .The result is exact to all orders.As mentioned earlier, to any finite order in g,the right hand side can be evaluated by a finite sum of Feynman diagrams.
We can consider the above relation to O\[g$`{}_{}{}^{0}]`$.Then it will give the free axial propagator compatible with the Lorentz gauges.This is automatically expected to give information as to how the 1/n.q -type singularities in the axial gauge propagator should be interpreted.\[The prescription problem.\].It reads,
iG<sup>0Aαβ</sup><sub>μν</sub>(x-y)=iG<sup>0Lαβ</sup><sub>μν</sub>(x-y)+i$``$$`__0`$<sup>1</sup>d$`\kappa `$ D$`\varphi `$ exp{ i S<sub>eff</sub>$`{}_{}{}^{M}[`$$`\varphi `$,$`\kappa ]`$ +$`ϵ`$$``$d$`{}_{}{}^{4}x`$($`\frac{1}{2}`$AA -$`\overline{c}`$ c)}
$`[`$$``$<sub>μ</sub>c<sup>α</sup>(x)A<sup>β</sup>$`{}_{\nu }{}^{}(y)`$+A$`{}_{\mu }{}^{\alpha }(x)`$$`_\nu `$c<sup>β</sup>(y)\] $`d^4z`$ $`\overline{c}`$<sup>γ</sup>(z) \[$`.`$A<sup>γ</sup>($`z)`$ -$`\eta `$.A<sup>γ</sup>(z)\] (5.11)
We, of course, need to evaluate the last term to O\[g$`{}_{}{}^{0}]`$. G<sup>0</sup> refers to free propagator.This is discussed further in Sec.7.
## 6 AN ALTERNATE DERIVATION BASED ON BRS:
The result (5.8 )can also be given a direct derivation based on BRS.Since, for many calculations in practice, the latter result is sufficient, we give a glance of the procedure.
We write expectation of O\[$`\varphi `$\] for the mixed gauge:
\<\<O\[$`\varphi `$\]\>\><sub>κ</sub>$`=`$ D$`\varphi `$O\[$`\varphi `$\]exp{ i S<sub>eff</sub>$`{}_{}{}^{M}[`$$`\varphi `$,$`\kappa ]`$} (6.1)
And evaluate d/d$`\kappa `$ of the above.We simplify the result for $`\frac{d}{d\kappa }`$\<\<O\[$`\varphi `$\]\>\><sub>κ</sub> using ,in particular,the BRS WT-identity of the Mixed gauge.(For the complete technical details,please see ). We then integrate the result over $`\kappa `$ from 0 to 1 in the end.This leads us to the equation (5.8).
We note ,however,as remarked earlier,that a complete definition of the Lorentz/Axial gauge Green’s function \<\<O\[$`\varphi `$\]\>\> requires an appropriate $`ϵ`$-term in each case as in equation (5.9 ).The derivation of (5.9 )in reference 1, however,did require the use of FFBRS/IFBRS construction in sections 4 and 5 explicitly.
## 7 APPLICATION TO THE 1/n.q -type SINGULATITIES
The equation (5.11 ), in the momentum space, reads:
G<sup>0A</sup><sub>μν</sub>(k) = G<sup>0L</sup><sub>μν</sub>(k) + i$``$$`__0`$<sup>1</sup>d$`\kappa `$ \[k<sub>μ</sub>$`\stackrel{~}{G^{0M}(k,\kappa )}`$(-ik<sup>σ</sup>-$`\eta ^\sigma )`$$`\stackrel{~}{G_{\sigma \nu }^{0M}(k,\kappa )}`$\+ ($`\mu `$,k) $``$ (v,-k) \] (7.1)
Here,$`\stackrel{~}{G^{0M}(k,\kappa )}`$ and $`\stackrel{~}{G_{\sigma \nu }^{0M}(k,\kappa )}`$are the propagators arising from the mixed gauge action together with O($`ϵ`$) terms viz.S<sub>eff</sub>$`{}_{}{}^{M}[`$$`\varphi `$,$`\kappa ]`$ -i$`ϵ`$$``$d$`{}_{}{}^{4}x`$($`\frac{1}{2}`$AA -$`\overline{c}`$ c).This can be integrated \[ 2,3\] to yield an expression for G<sup>0A</sup><sub>μν</sub>(k).
At first sight the expression for G<sup>0A</sup><sub>μν</sub>(k) appears hopelessly complicated; but as discussed below,a much much simpler effective expression can be arrived at from it.G<sup>0A</sup><sub>μν</sub>(k) reads :
G<sup>0A</sup><sub>μν</sub>(k)=G<sup>0L</sup><sub>μν</sub>(k)+ \[(k<sub>μ</sub>k<sub>ν</sub>$`\mathrm{\Sigma }_1+\eta _\mu k_\nu `$$`\mathrm{\Sigma }`$<sub>2</sub>) ln $`\mathrm{\Sigma }_3`$ \+ (k$``$-k,$`\mu \nu `$) (7.2)
where,
$`\mathrm{\Sigma }`$<sub>1</sub>$``$$`\frac{(k^2i\eta .k)[\frac{\eta .k+i\eta ^2}{k^2i\eta .k}+i\lambda \frac{(1\lambda )\eta .k}{k^2+iϵ}]}{ϵ\mathrm{\Sigma }}`$;
$`\mathrm{\Sigma }`$<sub>2</sub>$``$$`\frac{(k^2i\eta .k)[(\frac{i\eta .k+k^2}{k^2i\eta .k})+1\frac{(1\lambda )iϵ}{k^2+iϵ}]}{ϵ\mathrm{\Sigma }}`$;
$`\mathrm{\Sigma }`$<sub>3</sub>$``$$`\frac{i(\eta .k+ϵ)(k^2+iϵ\lambda )}{(k^2+iϵ)(iϵ\lambda \sqrt{k^4(k^2+iϵ\lambda )[k^2+\frac{(\eta .k)^2+iϵ\eta ^2}{k^2+iϵ}]}}`$;
with
$`\mathrm{\Sigma }`$$``${(1-$`\lambda `$)\[($`\eta `$.k)<sup>2</sup>\+ 2ik$`{}_{}{}^{2}\eta .k]+iϵ`$k$`{}_{}{}^{2}(12\lambda )`$ \+ $`\lambda `$(k$`{}_{}{}^{2}+iϵ)^2+\eta ^2(k^2+iϵ`$)+$`ϵ^2\lambda \}`$ (7.3)
These $`\mathrm{\Sigma }`$’s are somewhat complicated functions such that the propagator reduces to the usual propagator away from $`\eta `$.k = 0.The above complicated structure is not important in itself; but only its behavior as $`ϵ`$ –\> 0.In Ref. 2,3 we have established a procedure for extracting the effective term in this limit. For $`\eta _0`$0,it can be effectively replaced by:
An integral over the contour in the complex k$`{}_{0}{}^{}plane`$ that passes below<sup>4</sup><sup>4</sup>4 An equivalent treatment where the contour goes *above* the singularity can also be given. the singular point $`\eta `$.k=0, over a semi-circle of radius \>\> $`\sqrt{ϵ}`$ ; so that we may replace here the propagator by its naive expression;
And a term<sup>5</sup><sup>5</sup>5 There was an error in the results of Ref.2,3;please see errata. ,<sup>6</sup><sup>6</sup>6 We have set $`\eta _0=1`$ in the following and are assuming k$`{}_{}{}^{2}0;.`$and $`\eta `$<sup>2</sup>$``$ 0.We may take light-cone gauge limit at the end however. :
D$`{}_{\mu \nu }{}^{extra}(k)=\delta (`$k$`{}_{0}{}^{}\frac{1}{2}\sqrt{iϵ\eta ^2}𝐤.\eta )`$ {k<sub>μ</sub>k<sub>ν</sub>\[i$`\sqrt{\frac{i\eta ^2}{ϵ}}`$ +$`\frac{\eta ^2}{(\eta ^2+iϵ)}`$\] $``$$`\frac{2\pi \eta ^2}{[𝐤^2(\eta .𝐤)^2](\eta ^2+iϵ)}`$
(+\[k$`{}_{\mu }{}^{}\eta _{\nu }^{}+k_\nu \eta _\mu ]`$$`\frac{i\pi \eta ^2}{[𝐤^2(\eta .𝐤)^2](\eta ^2+iϵ)}`$} (7.4)
We note in passing that the derivation we followed requires that we keep $`\eta ^2`$$``$0 in the intermediate stages of our calculations.We can however obtain light-cone gauge results as a limiting case in the final form. We find that in this limit \[$`\eta ^2`$$``$0, $`ϵ`$ fixed\], the extra terms D$`{}_{\mu \nu }{}^{extra}(k)`$ vanish.
## 8 APPLICATION TO THE WILSON LOOP AND THE THERMAL LOOP
The Wilson Loop is : W\[L\] = \< T P exp i $``$$`__C`$A<sub>μ</sub>dx<sup>μ</sup> \> (8.1)
The preservation of the Wilson Loop \[i.e. its having the same value for axial and the Lorentz Gauges\] has been taken as an important test of a prescription.For example, the principal value prescription has been found wanting this test at the O\[g<sup>4</sup>\] level. For the L-M prescription, the result has been verified to O\[g<sup>4</sup>\] .
Since, in our treatment, the expectation values of the gauge-invariant observables is always preserved under the FFBRS that takes one from the Axial to the Lorentz gauges, the formal proof that our procedure, of which the axial propagator expression is a special element, automatically preserves the Wilson Loop.
In reference 4, we do not rest at this formal argument:we verify it to O\[g$`{}_{}{}^{4}]`$.In doing so we find it useful to make a connection with the work of Cheng & Tsai. Using the argument in this work,we can in fact give a simple argument that verifies the above result to O\[g$`{}_{}{}^{4}]`$ for axial gauges \[$`\eta ^2`$ \< 0\] for arbitrary loops\[ not necessarily planar\] and in other cases \[i.e. $`\eta ^2`$$``$0\] for a large class of loops.
The periodic Wilson Loop \[Thermal loop\] is another test of a prescription.It is a rectangular Wilson loop in Euclidean space in the \[x,$`\tau `$\] plane with side along $`\tau `$ axis from 0 to 1/$`\beta .`$The numerical result for the loop to O\[g<sup>2</sup>\] is easy to evaluate explicitly . This result for the loop to O\[g<sup>2</sup>\] reads<sup>7</sup><sup>7</sup>7 There were unfortunate trivial typographical errors in related to D<sub>00</sub> \[k<sub>0</sub>=0, k\] in the result for the thermal loop,see erratum. :
W<sub>R</sub> = 1+$`\frac{(N_c^21)g^2}{2N_cT}`$$``$$`\frac{d^3k}{(2\pi )^3}`$\[1-$`\mathrm{cos}`$(k.R)\]D<sub>00</sub> \[k<sub>0</sub>=0, k\] (8.2)
The above quantity, depends on D<sub>00</sub> \[k<sub>0</sub>=0, k\]. It is easy to verify from the expression (8.2), that this quantity has the same value as in Lorentz gauges.Thus, the thermal Wilson loop is preserved to O\[g<sup>2</sup>\].Thus, Our prescription satisfies this test also.
## 9 SUMMARY AND COMMENTS:
\[A\] We generated a field-dependent BRS \[FFBRS\] transformation that connects different classes of gauges \[the procedure is general enough\];this transformation preserves the vacuum expectation value of gauge-invariant observables;
\[B\] We applied it to the connection between the Lorentz and the axial-type gauges ;
\[C\]We gave a compact result that can relate the Green’s functions in the two classes of gauges;
\[D\]We applied the procedure to the question of the treatment of the axial gauge pole prescription;
\[E\]We showed that the Wilson Loop is preserved under our transformation and verified the result to O\[g<sup>4</sup>\].
Recently,we have generalized these results to an arbitrary pair of gauges for the FPEA and to the planar gauges.
## 10 FUTURE PLANS
Finally we summarize the future plans of development along these lines of work.
Use the understanding developed to address to the existing problems with the axial and Coulomb gauges .
Use the relation (5.9) to study Green’s functions (both primary and of simple gauge-invariant operators ) in two sets of gauges; and their renormalization properties and anomalous dimensions.
Develop formal techniques for establishing gauge-independence of arbitrary observables such as observable anomalous dimensions and cross-sections.
Work along these lines is in progress.
Acknowledgement:
The work in part was supported by grant for the DST project No. DST/PHT/19990170.
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# Light Sgoldstino: Precision Measurements versus Collider Searches
## 1 Introduction
Superpartners of goldstino — longitudinal component of gravitino — may be fairly light. In a variety of models with low energy supersymmetry they are lighter than a few GeV. Such pattern emerges in a number of non-minimal supergravity models and also in gauge mediation models if supersymmetry is broken via non-trivial superpotential (see, e.g., Ref. and references therein). To understand that superpartners of goldstino may be light, it suffices to recall that in globally supersymmetric theories with canonical Kähler potential and in the absence of anomalous abelian gauge factors, the sum of scalar squared masses is equal to the sum of fermion squared masses in each separate sector of the spectrum. Since goldstino is massless, its spinless superpartners (scalar and pseudoscalar particles, $`S`$ and $`P`$, hereafter, sgoldstinos) are massless too; they are associated with a non-compact flat direction of the scalar potential. Higher order terms from the Kähler potential contribute to sgoldstino masses. Provided these terms are sufficiently suppressed, sgoldstinos remain light. Of course, these arguments in no way guarantee that sgoldstinos are always light, but they do indicate that small sgoldstino masses are rather generic. The theoretical discussion of sgoldstino masses is contained, e.g., in Ref. ; here we merely assume that sgoldstinos are light and consider their phenomenology.
Sgoldstinos couple to MSSM fields in the same way as goldstino ; constraints on their couplings may be translated into the limits on the supersymmetry breaking parameter $`F`$.
There are several papers devoted to astrophysical , cosmological and collider constraints on models with light sgoldstinos. However, the role of light sgoldstinos in low-energy laboratory measurements has not been studied in detail. To the best of author’s knowledge, the only paper discussing this issue, Ref. , concentrated on sgoldstino contribution (as well as the contribution from light gravitino) into anomalous magnetic moment of muon. Here we consider a variety of low energy experiments sensitive to light sgoldstinos.
In this paper we identify those experiments which are most sensitive to different sgoldstino vertices for various sgoldstino masses. These experiments provide constraints on the corresponding coupling constants. These constants are in fact proportional to the ratios of soft terms (squark and gaugino masses, trilinear coupling constants) and $`F`$. The latter parameter is related to the gravitino mass $`m_{3/2}`$ in a simple way, $`F=\sqrt{3/(8\pi )}m_{3/2}M_{pl}`$; small $`F`$ corresponds to light gravitino ($`m_{3/2}<M_{SUSY}`$). Hence, the constraints derived in this paper are of importance for models with light gravitino, whereas sgoldstino effectively decouple from the visible sector in models with heavy gravitino.
In principle, there are both flavor-conserving and flavor-violating sgoldstino couplings to fermions. We present our results in the form of bounds on $`\sqrt{F}`$ setting soft flavor-conserving terms to be of the order of electroweak scale, as motivated by the supersymmetric solution to the gauge hierarchy problem. Flavor-violating couplings are governed by soft off-diagonal entries in squark (slepton) squared mass matrices. When evaluating bounds on $`\sqrt{F}`$ we set these off-diagonal entries equal to their current limits derived from the absence of FCNC and lepton flavor violation in MSSM . In this way we estimate the sensitivity of various experiments to the supersymmetry breaking scale.
We consider only low-energy processes with sgoldstinos on mass-shell. Processes with sgoldstino exchange deserve separate discussion, though we do not expect that the results obtained in this paper will be altered significantly. Also, behind the scope of this paper are loop processes with virtual sgoldstinos running in loops (for instance, $`K^0\overline{K}^0`$-mixing, $`\mu e\gamma `$, etc.). These processes were analyzed in Ref. in models with heavy sgoldstinos. Constraints on $`F`$ obtained in Ref. are significantly weaker than ones presented in our paper, so the loop processes are less sensitive to $`F`$ in models with heavy sgoldstinos. However, models with light sgoldstinos have not been analyzed in detail yet, though it was pointed out in Ref. that enhancement effects may appear if sgoldstinos are light. In view of the results obtained in this paper we also find it conceivable that light virtual sgoldstinos may give significant contributions into rare processes considered in Ref. .
Let us briefly review the current status of experimental limits on $`F`$. If one ignores sgoldstino, then in models with light gravitino the strongest direct current bound on $`F`$ is obtained from Tevatron, $`\sqrt{F}>217`$ GeV . In models with light sgoldstinos, collider experiments become more sensitive to the scale of supersymmetry breaking. Namely, LEP and Tevatron provide constraints at the level of 1 TeV on the supersymmetry breaking scale in models with $`m_{S(P)}`$ of order of 20 GeV . The most stringent cosmological constraint comes from Big Bang Nucleosynthesis : models with light gravitino, $`m_{3/2}<1`$ eV, that corresponds to $`\sqrt{F}<710^4`$ GeV, are disfavored if sgoldstinos decouple at temperature not less than $`𝒪(100)`$ MeV ($`m_{S(P)}1`$ MeV). It has been argued in Ref. that among the astrophysical constraints, the strongest one comes from SN1987A: the gravitino mass is excluded in the range $`10^{1.5}\mathrm{eV}<m_{3/2}<30`$ eV for 1 keV$`<m_{S(P)}<10`$ MeV and in a wider range $`310^6\mathrm{eV}<m_{3/2}<50`$ eV for $`m_{S(P)}<1`$ keV. These excluded intervals correspond to $`10^4\mathrm{GeV}<\sqrt{F}<410^5`$ GeV and $`120\mathrm{GeV}<\sqrt{F}<510^5`$ GeV, respectively.
In this paper we consider various constraints on couplings of light ($`m_{S(P)}5`$ GeV) (pseudo)scalars to SM fields coming mostly from astrophysics and direct precision measurements. So, we partially fill the gap between constraints coming from collider experiments and cosmology.
As there are flavor-conserving and flavor-violating interactions of sgoldstino fields, we have to consider both flavor-symmetric and flavor asymmetric processes. Let us outline our results referring to these two cases in turn.
We begin with constraints independent of assumptions concerning breaking of flavor symmetry. As expected, strongest bounds arise from astrophysics and cosmology, that is $`\sqrt{F}10^6`$ GeV, or $`m_{3/2}>600`$ eV, for models with $`m_{S(P)}<10`$ keV and MSSM soft flavor-conserving terms being of the order of electroweak scale. For the intermediate sgoldstino masses (up to a few MeV) constraints from the study of SN explosion and reactor experiments lead to $`\sqrt{F}300`$ TeV. We will find that for heavier sgoldstinos, low energy processes (such as rare decays of mesons) provide limits comparable to ones from colliders but valid for different sgoldstino masses.
As concerns flavor-asymmetric processes, we find that these are generally very sensitive to light sgoldstino. Namely, with flavor-changing off-diagonal entries in squark (slepton) squared mass matrix close to the current bounds, direct measurements of decays of mesons (leptons) provide very strong bounds, up to $`\sqrt{F}900(15000)`$ TeV (valid at $`m_S5(0.34)`$ GeV), which is much higher than bounds expected from future colliders. If off-diagonal entries are small, the limits on $`\sqrt{F}`$ become weaker: they scale as square root of the corresponding off-diagonal elements.
We will see that the rates of processes with one sgoldstino in final state are proportional to $`F^2`$, whereas the rates of processes with two sgoldstinos in final state are proportional to $`F^4`$. Hence, under similar assumptions about soft terms governing sgoldstino couplings, processes with one sgoldstino are more sensitive to the supersymmetry breaking scale. Nevertheless, the coupling constants entering one-sgoldstino and two-sgoldstino processes are generally determined by different parameters, so the study of two-sgoldstino processes is also important.
Further progress in the search for sgoldstino is expected in several directions. Among the laboratory experiments, the most sensitive to flavor-conserving sgoldstino coupling for sgoldstino lighter than a few MeV are experiments with laser photons propagating in magnetic fields and reactor experiments. For heavier sgoldstinos, measurements of $`\mathrm{{\rm Y}}`$ partial widths exhibit the best discovery potential. If flavor violation in MSSM is sufficiently strong (say, at the level of current limits), the most promising is the study of charged kaon decays.
This paper is organized as follows. In section 2 the effective lagrangian for sgoldstinos is presented and sgoldstino decay modes are described. In section 3 we derive various constraints on the parameter of supersymmetry breaking $`\sqrt{F}`$ by considering low energy processes. There we study separately processes with one and two sgoldstinos in final states (sections 3.1 and 3.2, respectively). First, we discuss astrophysical and cosmological limits on sgoldstino interactions (section 3.1.1). Then we present laboratory bounds coming from search for light (pseudo)scalars in electromagnetic and strong processes (section 3.1.2). In sections 3.1.3 and 3.1.4. we discuss rare decays with one sgoldstino in final state due to flavor-conserving and flavor-violating sgoldstino couplings to SM fermions, respectively. Sections 3.2.1 and 3.2.2 are devoted to rare meson decays with two sgoldstinos in final state. Our conclusions and comparison of the results with ones coming from collider experiments are presented in section 4.
## 2 Effective lagrangian
Let us introduce the effective lagrangian for light goldstino supermultiplet: scalar $`S`$, pseudoscalar $`P`$ and goldstino $`\stackrel{~}{G}`$. The free part reads
$$=\frac{1}{2}\left(_\mu S^\mu Sm_S^2S^2\right)+\frac{1}{2}\left(_\mu P^\mu Pm_P^2P^2\right)+\frac{i}{2}\overline{\stackrel{~}{G}}\gamma ^\mu _\mu \stackrel{~}{G}.$$
There exist two types of interactions in the low-energy effective theory involving sgoldstino fields: these are terms that couple one sgoldstino and two sgoldstinos , respectively, to SM gauge fields (photons, gluons) and matter fields (leptons $`f_L`$, up- and down-quarks $`f_U`$ and $`f_D`$). Terms involving one sgoldstino are
$`_{eff}={\displaystyle \frac{1}{2\sqrt{2}F}}\left(m_S^2S\overline{\stackrel{~}{G}}\stackrel{~}{G}+im_P^2P\overline{\stackrel{~}{G}}\gamma _5\stackrel{~}{G}\right){\displaystyle \frac{1}{2\sqrt{2}}}{\displaystyle \frac{M_{\gamma \gamma }}{F}}SF^{\mu \nu }F_{\mu \nu }+{\displaystyle \frac{1}{4\sqrt{2}}}{\displaystyle \frac{M_{\gamma \gamma }}{F}}Pϵ^{\mu \nu \rho \sigma }F_{\mu \nu }F_{\rho \sigma }{\displaystyle \frac{1}{2\sqrt{2}}}{\displaystyle \frac{M_3}{F}}SG^{\mu \nu \alpha }G_{\mu \nu }^\alpha `$ (1)
$`+{\displaystyle \frac{1}{4\sqrt{2}}}{\displaystyle \frac{M_3}{F}}Pϵ^{\mu \nu \rho \sigma }G_{\mu \nu }^\alpha G_{\rho \sigma }^\alpha {\displaystyle \frac{\stackrel{~}{m}_{D_{ij}}^{LR2}}{\sqrt{2}F}}S\overline{f}_{D_i}f_{D_j}i{\displaystyle \frac{\stackrel{~}{m}_{D_{ij}}^{LR2}}{\sqrt{2}F}}P\overline{f}_{D_i}\gamma _5f_{D_j}{\displaystyle \frac{\stackrel{~}{m}_{U_{ij}}^{LR2}}{\sqrt{2}F}}S\overline{f}_{U_i}f_{U_j}i{\displaystyle \frac{\stackrel{~}{m}_{U_{ij}}^{LR2}}{\sqrt{2}F}}P\overline{f}_{U_i}\gamma _5f_{U_j}`$
$`{\displaystyle \frac{\stackrel{~}{m}_{L_{ij}}^{LR2}}{\sqrt{2}F}}S\overline{f}_{L_i}f_{L_j}i{\displaystyle \frac{\stackrel{~}{m}_{L_{ij}}^{LR2}}{\sqrt{2}F}}P\overline{f}_{L_i}\gamma _5f_{L_j}.`$
The direct coupling of two sgoldstinos is described by
$`_{eff}={\displaystyle \frac{1}{4F^2}}(S_\mu PP_\mu S)((\stackrel{~}{m}_{L_{ij}}^{LL2}+\stackrel{~}{m}_{L_{ij}}^{RR2})\overline{f}_{L_i}\gamma ^\mu \gamma _5f_{L_j}+(\stackrel{~}{m}_{L_{ij}}^{LL2}\stackrel{~}{m}_{L_{ij}}^{RR2})\overline{f}_{L_i}\gamma ^\mu f_{L_j}`$ (2)
$`+(\stackrel{~}{m}_{D_{ij}}^{LL2}+\stackrel{~}{m}_{D_{ij}}^{RR2})\overline{f}_{D_i}\gamma ^\mu \gamma _5f_{D_j}+(\stackrel{~}{m}_{D_{ij}}^{LL2}\stackrel{~}{m}_{D_{ij}}^{RR2})\overline{f}_{D_i}\gamma ^\mu f_{D_j}`$
$`+(\stackrel{~}{m}_{U_{ij}}^{LL2}+\stackrel{~}{m}_{U_{ij}}^{RR2})\overline{f}_{U_i}\gamma ^\mu \gamma _5f_{U_j}+(\stackrel{~}{m}_{U_{ij}}^{LL2}\stackrel{~}{m}_{U_{ij}}^{RR2})\overline{f}_{U_i}\gamma ^\mu f_{U_j}).`$
Here $`M_{\gamma \gamma }=M_1\mathrm{cos}^2\theta _W+M_2\mathrm{sin}^2\theta _W`$ and $`M_i`$ are gaugino masses; for down-quarks $`i=d,s,b`$, whereas for up-quarks $`i=u,c,t`$; $`\stackrel{~}{m}_{ij}^{LR2}`$, $`\stackrel{~}{m}_{ij}^{LL2}`$ and $`\stackrel{~}{m}_{ij}^{RR2}`$ are LR-, LL-, and RR-soft mass terms in squark squared mass matrix and for convenience we take them real. In what follows we do not discuss neutrino, so the corresponding couplings are omitted. Note that in MSSM the flavor-conserving one-sgoldstino coupling constants satisfy $`\stackrel{~}{m}_{ii}^{LR2}=m_{f_i}A_{f_i}`$, where $`m_{f_i}`$ are fermion masses and $`A_{f_i}`$ are corresponding soft trilinear coupling constants. Off-diagonal soft terms $`\stackrel{~}{m}_{ij}^{LR2}`$, $`\stackrel{~}{m}_{ij}^{LL2}`$ and $`\stackrel{~}{m}_{ij}^{RR2}`$ are subject to constraints from the absence of FCNC and lepton flavor violation (see, e.g., Ref. ).
The first part of the effective lagrangian, Eq. (1), is suppressed by $`F^1`$, whereas the second one, Eq. (2), is proportional to $`F^2`$, so processes with two sgoldstinos are very rare. The most stringent bounds on $`F`$ come from processes with one sgoldstino in final state. Nevertheless, as we will see, the absence of processes with two sgoldstinos gives rise to constraints on supersymmetry breaking parameter $`F`$ comparable to bounds from high-energy experiments. The latter constraints are, strictly speaking, independent of the constraints coming from one-sgoldstino processes: one-sgoldstino and direct two-sgoldstino processes are governed by $`\stackrel{~}{m}^{LR2}`$ and $`\stackrel{~}{m}^{LL2}`$, $`\stackrel{~}{m}^{RR2}`$, respectively.
Let us discuss decay modes of light sgoldstino. First, sgoldstino decay into two photons is always open ,
$$\mathrm{\Gamma }(S(P)\gamma \gamma )=\frac{m_{S(P)}^3M_{\gamma \gamma }^2}{32\pi F^2}.$$
(3)
Second, in models where $`m_{3/2}<m_{S(P)}`$ sgoldstinos may decay into gravitino pairs; however, the corresponding rates are suppressed by squared ratio of sgoldstino mass $`m_{S(P)}`$ and $`M_{\gamma \gamma }`$ in comparison with the decay into two photons, hence this mode may be disregarded. Third, relatively heavy sgoldstinos ($`m_{S(P)}\mathrm{\Lambda }_{QCD}`$) decay into gluons (light mesons) with larger width than into photons because of color enhancement and because the corresponding coupling is proportional to gluino mass which is usually the largest among the gaugino masses, i.e. $`M_3>M_{\gamma \gamma }`$. When analyzing hadronic decay modes of light sgoldstinos ($`m_{S(P)}<afew`$ GeV), corresponding rates into quarks and gluons should be rewritten in terms of light mesons. This step will be presented below. Fourth, sgoldstino can decay also into light leptons if this process is allowed kinematically ($`m_{S(P)}>2m_l`$). Since the corresponding coupling constants are proportional to fermion masses these rates are suppressed by a factor $`m_l^2/m_{P(S)}^2`$ apart from the phase space volume ,
$$\mathrm{\Gamma }(Sl\overline{l})=\frac{m_S^3A_l^2}{16\pi F^2}\frac{m_l^2}{m_S^2}\left(1\frac{4m_l^2}{m_S^2}\right)^{3/2},\mathrm{\Gamma }(Pl\overline{l})=\frac{m_P^3A_l^2}{16\pi F^2}\frac{m_l^2}{m_P^2}\left(1\frac{4m_l^2}{m_P^2}\right)^{1/2}.$$
(4)
Consequently, depending on MSSM mass spectrum, sgoldstino masses and the value of the supersymmetry breaking parameter $`F`$, there are three possible situations in experiments where light (pseudo)scalar particle appears. This particle may live long enough to escape from a detector. For instance, in the theory with the superpartner scale of order 100 GeV and $`\sqrt{F}=1`$ TeV this behavior would be exhibited by (pseudo)scalar particle with mass less than 10 MeV, at which sgoldstino width is saturated by two-photon mode. Another case is when (pseudo)scalar particle decays within detector into two photons or leptons. Apart from these cases, there is also a possibility of the decay into two gluons (quarks). For relatively light sgoldstinos (but with masses exceeding 270 MeV), the dominant hadronic decay is into two pions, while for heavier sgoldstinos $`KK`$ and $`\eta \eta `$ channels become available. Furthermore, there would be effects emerging due to $`P\pi ^0(\eta ,K^0)`$ mixing.
Let us estimate branching ratios of hadronic and photonic decay channels neglecting threshold factor. In order to estimate sgoldstino coupling to hadrons we make use of chiral theory of light hadrons. There are two different sources of sgoldstino-meson couplings in the effective lagrangian (1): interaction terms with gluons and coupling to quarks. We evaluate contributions from these two sources into meson-sgoldstino interactions separately.
First, we have to relate gluonic operators entering Eq. (1) to meson fields. We make use of the correspondence
$$(\pi \pi )_{J=0}|\frac{\beta (\alpha _s)}{8\pi \alpha _s}G_{\mu \nu }^aG^{a\mu \nu }|0=\frac{1}{2}q^2\phi _\pi ^\alpha \phi _\pi ^\alpha $$
(5)
derived in Ref. . Here $`q^2`$ is momentum of pion pair created with zero total angular momentum, $`J=0`$; $`\beta (\alpha _s)`$ is the $`\beta `$-function of QCD, $`\phi _\pi ^\alpha `$ is the pion isotopic amplitude,
$$\phi _\pi ^\alpha \phi _\pi ^\alpha =2\phi _{\pi ^+}\phi _\pi ^{}+\phi _{\pi ^0}\phi _{\pi ^0},$$
and quarks and mesons are considered massless. At higher energies also $`KK`$ and $`\eta \eta `$ pairs may be created by gluonic operator.
There is one more relation ,
$$A|\frac{N_f\alpha _s}{4\pi }G_{\mu \nu }^a\stackrel{~}{G}^{a\mu \nu }|0=constϵf_Am_A^2\phi _A,$$
(6)
where $`\stackrel{~}{G}^{a\mu \nu }`$ is a tensor dual to gluonic one, $`A`$ is a neutral pseudoscalar meson ($`\pi ^0`$, $`\eta `$) and $`const`$ is a normalization factor; $`f_A=f_\pi =130`$ MeV and $`ϵ`$ is a parameter responsible for $`SU(N_f)`$ flavor symmetry breaking ($`ϵ=(m_um_d)/(m_u+m_d)`$ for $`\pi ^0`$, $`ϵ1`$ for $`\eta `$).
In fact, the lagrangian (1) describes sgoldstino interactions at the superpartner scale. Sgoldstino coupling constants at low energies may be obtained by making use of renormalization group evolution. Thus for the gluonic operator one has
$$G_{\mu \nu }^2(M_3)=G_{\mu \nu }^2(\mu )\frac{\beta (\alpha _s(\mu ))}{\alpha _s(\mu )}\frac{\alpha _s(M_3)}{\beta (\alpha _s(M_3))}.$$
Hence, we estimate the matrix element of the gluonic operator between the scalar and meson pair as
$$(AA)_{J=0}|\frac{M_3}{2\sqrt{2}F}G_{\mu \nu }^aG^{a\mu \nu }S|S=\frac{\alpha _s(M_3)}{\beta (\alpha _s(M_3))}q^2\sqrt{2}\pi \phi _A\phi _A\frac{M_3}{F}\phi _S$$
(7)
and in a similar way we estimate the matrix element of another gluonic operator between the pseudoscalar and meson
$$A|\frac{M_3}{2\sqrt{2}F}G_{\mu \nu }^a\stackrel{~}{G}^{a\mu \nu }P|P=\frac{ϵconst}{\alpha _s(M_3)}\frac{\sqrt{2}\pi }{N_f}f_Am_A^2\phi _A\frac{M_3}{F}\phi _P.$$
(8)
Note, that these matrix elements are highly suppressed by squared sgoldstino or meson masses.
Since direct sgoldstino coupling to quarks contributes also to meson production, we remind basic relations of chiral theory
$$0|J_\mu ^{\pi ^0}(0)|\pi ^0(\stackrel{}{q})=\frac{i}{\sqrt{2}}f_\pi q_\mu ,0|J_\mu ^{\pi ^+}(0)|\pi ^+(\stackrel{}{q})=if_\pi q_\mu .$$
(9)
where
$$J_\mu ^{\pi ^0}=\frac{1}{2}\left(\overline{u}\gamma _\mu \gamma _5u\overline{d}\gamma _\mu \gamma _5d\right),J_\mu ^{\pi ^+}=\overline{d}\gamma _\mu \gamma _5u.$$
(10)
If we parameterize sgoldstino couplings to the triplet of light quarks $`q`$ as
$$=\overline{q}\left(S\widehat{\mathrm{\Sigma }}_Si\gamma _5P\widehat{\mathrm{\Sigma }}_P\right)q$$
with $`\widehat{\mathrm{\Sigma }}_S`$ and $`\widehat{\mathrm{\Sigma }}_P`$ being $`3\times 3`$ matrices of the corresponding coupling constants (which are read off from Eq. (1)), then the standard procedure (see, e.g., Ref. ) gives the following low-energy effective lagrangian
$$_{meson}=B_0\mathrm{Tr}\left(f_{\pi ^0}\widehat{\mathrm{\Phi }}\widehat{\mathrm{\Sigma }}_PPS\widehat{\mathrm{\Sigma }}_S\widehat{\mathrm{\Phi }}^2\right)$$
(11)
to the leading order in mesonic fields included in matrix $`\widehat{\mathrm{\Phi }}`$. The constant $`B_0`$ is related to quark condensate as $`0|\overline{q}q|0=\frac{1}{2}B_0f_{\pi ^0}^2`$ and may be evaluated from the masses of kaon and quarks, $`B_0=M_{K^0}^2/(m_d+m_s)`$. We account only for one-sgoldstino terms since others are suppressed by sgoldstino masses and additional inverse power of $`F`$.
The lagrangian (11) consists of two parts. The first one,
$$_{meson1}=\frac{B_0f_{\pi ^0}}{\sqrt{2}F}\left(\frac{\pi ^0}{\sqrt{2}}\left(\stackrel{~}{m}_{U_{11}}^{LR2}\stackrel{~}{m}_{D_{11}}^{LR2}\right)+\frac{\eta }{\sqrt{6}}\left(\stackrel{~}{m}_{U_{11}}^{LR2}+\stackrel{~}{m}_{D_{11}}^{LR2}2\stackrel{~}{m}_{D_{22}}^{LR2}\right)+K^0\stackrel{~}{m}_{D_{21}}^{LR2}+\overline{K^0}\stackrel{~}{m}_{D_{12}}^{LR2}\right)P,$$
(12)
is pseudoscalar sgoldstino mixing with $`\pi ^0`$, $`\eta `$, $`K^0`$ and $`\overline{K^0}`$ mesons, while the second one,
$`_{meson2}={\displaystyle \frac{B_0}{\sqrt{2}F}}((\stackrel{~}{m}_{U_{11}}^{LR2}+\stackrel{~}{m}_{D_{11}}^{LR2})(\pi ^+\pi ^{}+K^+K^{}+K^0\overline{K^0}+{\displaystyle \frac{1}{2}}\pi ^0\pi ^0)`$ (13)
$`+{\displaystyle \frac{1}{6}}\eta ^2\left(\stackrel{~}{m}_{U_{11}}^{LR2}+\stackrel{~}{m}_{D_{11}}^{LR2}+4\stackrel{~}{m}_{D_{22}}^{LR2}\right)+{\displaystyle \frac{1}{\sqrt{3}}}\pi ^0\eta \left(\stackrel{~}{m}_{U_{11}}^{LR2}\stackrel{~}{m}_{D_{11}}^{LR2}\right){\displaystyle \frac{1}{\sqrt{2}}}\pi ^0\overline{K^0}\stackrel{~}{m}_{D_{12}}^{LR2}{\displaystyle \frac{1}{\sqrt{2}}}\pi ^0K^0\stackrel{~}{m}_{D_{21}}^{LR2}`$
$`{\displaystyle \frac{1}{\sqrt{6}}}\overline{K^0}\eta \stackrel{~}{m}_{D_{12}}^{LR2}{\displaystyle \frac{1}{\sqrt{6}}}K^0\eta \stackrel{~}{m}_{D_{21}}^{LR2}+K^{}\pi ^+\stackrel{~}{m}_{D_{12}}^{LR2}+K^+\pi ^{}\stackrel{~}{m}_{D_{21}}^{LR2})S`$
describes scalar sgoldstino decays into mesons. Note that sgoldstino couplings with two different mesons is suppressed by off-diagonal term in squark mass matrix. In what follows we will not consider processes where real sgoldstino decays into such modes.
Now let us estimate matrix elements between sgoldstino and meson (i.e., sgoldstino-meson mixing) as a sum of two quantities, Eq. (8) and Eq. (12), while the amplitude of the scalar sgoldstino decay into pairs of light mesons is evaluated as a sum of Eq. (7) and Eq. (13). Let us compare contributions of gluon and quark operators into sgoldstino couplings to mesons. As an example, for the ratio of the corresponding contributions into coupling of the scalar to neutral pions and into pion-pseudoscalar mixing we obtain
$$\frac{\pi ^0\pi ^0|S_{gluon}}{\pi ^0\pi ^0|S_{quark}}=\frac{\alpha _s(M_3)}{\beta (\alpha _s(M_3))}4\pi \frac{m_S}{B_0}\frac{M_3}{A_Q}\frac{m_S}{m_u+m_d},$$
$$\frac{\pi ^0|P_{gluon}}{\pi ^0|P_{quark}}=\frac{2\pi \sqrt{2}}{\alpha _s(M_3)}\frac{M_3}{3A_Q}.$$
These ratios are larger than 10 for $`M_3=A_Q`$. Hence, gluonic operators give rise to stronger coupling of light sgoldstinos to light mesons, as compared to sgoldstino-quark interactions.
Let us evaluate the rate of the scalar sgoldstino decay into light mesons, assuming that this decay is allowed kinematically. As an example, for the neutral pion mode we obtain
$$\mathrm{\Gamma }(S\pi ^0\pi ^0)=\frac{\alpha _s^2(M_3)}{\beta ^2(\alpha _s(M_3))}\frac{\pi m_S}{324}\frac{m_S^2M_3^2}{F^2}\left(1\frac{\beta (\alpha _s(M_3))}{\alpha _s(M_3)}\frac{9}{4\pi }\frac{B_0}{m_S}\frac{m_u+m_d}{m_S}\frac{A_Q}{M_3}\right)^2\sqrt{1\frac{4m_{\pi ^0}^2}{m_S^2}}.$$
Taking into account only the largest contribution from the gluon operator and neglecting the threshold factor we estimate the ratio of rates of sgoldstino decays into photons and mesons,
$$\frac{\mathrm{\Gamma }(S\gamma \gamma )}{\mathrm{\Gamma }(S\pi ^0\pi ^0)}=\frac{81}{8\pi ^2}\frac{\beta ^2(\alpha _s(M_3))}{\alpha _s^2(M_3)}\frac{M_{\gamma \gamma }^2}{M_3^2}.$$
We see that this ratio is smaller than 1 at $`M_{\gamma \gamma }=M_3`$. Since in most models gluino is several times heavier than photino, for sufficiently heavy sgoldstinos hadronic modes usually dominate over photonic one.
Let us estimate now the contribution of pion-sgoldstino mixing into pseudoscalar sgoldstino width. Recall that the pion width is almost saturated by the two-photon decay mode. Then
$$\mathrm{\Gamma }(P\pi ^0\gamma \gamma )=\frac{1}{\alpha _s^2(M_3)}\frac{\pi ^2f_{\pi ^0}^2}{4\left(m_P^2m_{\pi ^0}^2\right)^2}\frac{M_3^2m_{\pi ^0}^4}{F^2}\left(\frac{m_um_d}{m_u+m_d}\right)^2\mathrm{\Gamma }^{}(\pi ^0\gamma \gamma ),$$
(14)
where the two-photon width of virtual pion is taken at $`p_{\pi ^0}^2=m_P^2`$ and may be approximated as
$$\mathrm{\Gamma }^{}(\pi ^0\gamma \gamma )\mathrm{\Gamma }_{tot}(\pi ^0)\frac{m_P^3}{m_{\pi ^0}^3}.$$
With account of only leading contributions from gluonic operator we obtain
$$\frac{\mathrm{\Gamma }^{direct}(P\gamma \gamma )}{\mathrm{\Gamma }(P\pi ^0\gamma \gamma )}=\frac{\alpha _s^2(M_3)}{8\pi ^3}\tau _{\pi ^0}m_{\pi ^0}\frac{M_{\gamma \gamma }^2}{M_3^2}\frac{m_{\pi ^0}^2}{f_{\pi ^0}^2}\left(\frac{m_P^2}{m_{\pi ^0}^2}1\right)^2.$$
(15)
As discussed above, $`\stackrel{~}{m}_{ii}^{LR2}=m_iA_{Q_i}`$, so at $`M_{\gamma \gamma }=M_3`$ and light $`P`$ ($`m_Pm_{\pi ^0}`$) we obtain that the ratio (15) is numerically $`810^2`$. In the opposite case of heavy $`P`$ ($`m_Pm_{\pi ^0}`$) the ratio becomes even larger. Hence mixing with pions gives negligible contribution to sgoldstino decay into photons (unless $`M_{\gamma \gamma }M_3/30`$; we do not consider this case). The only exception is the degenerate case, when sgoldstino and pion masses are close and this branching becomes of order 1. (In the case of strong degeneracy there is also a correction to pion life-time which may give rise to a constraint on $`F`$). We do not consider this unrealistic situation. The interference with $`\eta `$-meson gives nothing new. Indeed, Eq. (15) scales as $`\tau _{meson}m_{meson}^3`$ which is invariant under the variation of meson mass, if the meson width is (almost) saturated by anomalous decay into two photons. Decay via neutral kaon is also negligible because of large kaon life-time.
To conclude this section we summarize the situation with sgoldstino branching ratios. Let us begin with scalar sgoldstino. In Figure 1
we present scalar sgoldstino branching ratios into photons, leptons and neutral pions evaluated for four different sets of supersymmetry breaking soft terms, $`A`$, $`M_{\gamma \gamma }`$, $`M_3`$. To determine photonic and leptonic widths we make use of Eqs. (3) and (4), while the width into two neutral pions is calculated according to Eq. (5) generalized to non-zero pion masses. Estimating hadronic sgoldstino partial width we account only for $`\pi ^+\pi ^{}`$ and $`\pi ^0\pi ^0`$ decay modes. Other hadronic modes may be considered in the same way. Ratios between different hadronic channels are determined by chiral theory.
Scalar sgoldstino lighter than 270 MeV almost always predominantly decays into two photons. At sgoldstino mass close to $`2m_e`$ or $`2m_\mu `$, rates of the decays into pairs of corresponding leptons become comparable to the two-photon rate and even exceed the latter in models with large trilinear soft terms. Far from the lepton mass, the corresponding lepton branching ratio decreases as $`m_l^2/m_{S(P)}^2`$. At sgoldstino masses exceeding 270 MeV hadronic modes emerge. Their rates are somewhat higher than the rate of the two-photon decay except for models with large $`M_{\gamma \gamma }`$, in which the photonic mode always dominates.
As regards pseudoscalar sgoldstino, it does not have the decay mode into two pseudoscalar mesons to the zero order in $`G_F`$. Hence at $`M_{\gamma \gamma }M_3`$ its hadronic decay modes are suppressed unless $`m_P`$ is quite large (well above 1 GeV). In what follows we consider photonic and leptonic decay channels of the pseudoscalar sgoldstino only.
## 3 Searches for light sgoldstino
In accordance with the discussion of sgoldstino effective lagrangian presented in the previous section, there are two types of processes we are interested in. In the processes of the first type only one sgoldstino emerges while in the processes of the second type a pair of sgoldstinos appears in the final state. These processes are governed by different coupling constants and will be considered in turn.
### 3.1 Processes with one sgoldstino in the final state
#### 3.1.1 Bounds from astrophysics and cosmology
In subsections 3.1.1 and 3.1.2 we consider mainly pseudoscalar sgoldstino, though almost all constraints are valid for the scalar sgoldstino as well; a few exceptions will be pointed out.
Light pseudoscalar particles appear in particle physics models in various contexts; a well known example is an axion. There are numerous cosmological, astrophysical and laboratory bounds on interactions of light pseudoscalars which apply to sgoldstino. For completeness we collect in sections 3.1.1 and 3.1.2 the most stringent of these bounds and translate the bounds on sgoldstino coupling constants into bounds on supersymmetry breaking parameters $`\sqrt{F}`$ and $`m_{3/2}`$. Let us write the interactions of sgoldstino with photons and fermions as follows
$$_{\gamma P}=\sqrt{2}\frac{M_{\gamma \gamma }}{F}P\stackrel{}{E}\stackrel{}{B}g_\gamma P\stackrel{}{E}\stackrel{}{B},_{fP}=\frac{m_fA_f}{\sqrt{2}F}P\overline{f}\gamma _5fg_{fP}P\overline{f}\gamma _5f.$$
(16)
Then the limits on $`g_\gamma `$ and $`g_{fP}`$ imply limits on $`\sqrt{F}`$.
In what follows we set $`M_{\gamma \gamma }=A_f=100`$ GeV in our quantitative estimates, since superpartner scale $`M_{SUSY}`$ is expected to be close to the electroweak scale as motivated by supersymmetric solution to the hierarchy problem in SM.
One of the sources of pseudoscalars are stars: light pseudoscalars are produced there by Primakoff process, that is $`\gamma P`$ conversion in external electromagnetic field. Another place of sgoldstino creation is galactic space where magnetic fields produce pseudoscalars from propagating photons.
In “helioscope” method, a dipole magnet directed towards the Sun is used. Inside the volume with strong magnetic field, solar pseudoscalars can transform into X-rays by inverse Primakoff process. An alternative method, “Bragg diffraction”, was applied in SOLAX experiment to detect solar pseudoscalars. The absence of anomalous X-ray fluxes from SN1987A related to possible pseudoscalar conversion into photons in galactic magnetic field gives the strongest constraint on $`g_\gamma `$. Since this limit is valid only for unrealistically light pseudoscalar ($`m_P<10^9`$ eV), we consider the helium-burning life-time of Horizontal Branch Stars (HBS) in globular clusters as the most sensitive probe of $`F`$ at very small $`m_P`$.
There are two more constraints on $`g_\gamma `$ coming from cosmology and astrophysics. Light sgoldstinos are thermally produced in the early Universe via Compton process $`e\gamma eP(S)`$. Photons from sgoldstino decays contribute to the photon extragalactic background, if sgoldstinos outlive matter-radiation decoupling. If, on the other hand, sgoldstinos decay before matter-radiation decoupling, produced photons may heat electrons leading to distortion of CMBR spectrum, which is experimentally studied well enough to exclude wide range of $`F`$ at corresponding sgoldstino masses. The experiments on photon background and cosmic microwave background radiation, being combined, exclude a strip in $`(m_P,g_\gamma )`$ plane (see Ref. ). In Table 1 we present the corresponding limits for two typical values of $`m_P`$.
All these constraints <sup>1</sup><sup>1</sup>1See also Ref. for constraints on $`g_\gamma `$ coming from Deuterium fission by scalars decaying into two photons. on $`g_\gamma `$ are collected in Table 1. The limits on $`\sqrt{F}`$ are obtained at $`M_{\gamma \gamma }=100`$ GeV and scale as square root of $`M_{\gamma \gamma }`$. For completeness, we included in Table 1 also the limits obtained in Ref. by considering SN1987A.
Let us proceed with sgoldstino coupling to electrons. Restrictive limits come from delay of helium ignition in low-mass red giants. There are also two limits on coupling to electrons from bremsstrahlung process $`e^{}+(A,Z)(A,Z)+e^{}+P`$ and Compton process $`\gamma +e^{}e^{}+P`$ in stars: these processes lead to energy loss of stars and are constrained by helium-burning life-time of Horizontal Branch Stars. Note that the life-time of HBS gives stronger constraints on electron coupling to scalar than to pseudoscalar.
Let us turn to (pseudo)scalar coupling to nucleons. In order to relate the corresponding constant $`g_N`$ to $`F`$ we make use of the analogy to axion. Then effective lagrangian reads
$$_{eff}=i\overline{\psi }\gamma _5(g_N^{(0)}+g_N^{(3)}\tau _3)\psi P,$$
(17)
where $`\psi `$ denotes the nucleon dublet and
$$g_N^{(0)}\frac{A_Qm_N}{\sqrt{2}F},g_N^{(3)}\frac{m_um_d}{m_u+m_d}g_N^{(0)}.$$
(18)
The energy loss of Horizontal Branch Stars via Compton process $`\gamma +^4\mathrm{He}^4\mathrm{He}+S`$ gives rise to a bound on $`F`$. Also, nucleon-sgoldstino coupling leads to shortening of SN1987A neutrino burst.
Astrophysical constraints on sgoldstino-fermion interactions are presented in Table 2. Bounds on $`\sqrt{F}`$ are obtained at $`A_e=A_Q=100`$ GeV and scale as $`\sqrt{A_f}`$, $`f=Q,e`$. Note that the region $`\sqrt{F}1.510^4`$ GeV allowed by SN explosion is not ruled out by astrophysical arguments or direct measurements if sgoldstino is relatively heavy (10 keV$`m_{S(P)}`$10 MeV) and its interactions conserve flavor (see below).
For constraints coming from Big Bang Nucleosynthesis see Ref. .
#### 3.1.2 Laboratory bounds on very light sgoldstinos
Let us now consider direct laboratory limits on couplings of very light sgoldstinos.
The first set of bounds on $`F`$ comes from the study of laser beam propagation through transverse magnetic field. The production of real sgoldstinos would induce the rotation of the beam polarization, while the emission and absorption of virtual sgoldstinos would contribute to the ellipticity of the laser beam. Such effects have not been observed and their absence implies a constraint on pseudoscalar-photon coupling. There is also a constraint on the interaction of a pseudoscalar particle with photons coming from experiments on photon regeneration. In these experiments, light pseudoscalars produced via Primakoff effect penetrate through optic shield and then transform back into photons (“invisible light shining through walls”). Similar scheme is applied in NOMAD experiment. The results are presented in Table 3 at $`M_{\gamma \gamma }=100`$ GeV; limits on $`\sqrt{F}`$ scale as $`\sqrt{M_{\gamma \gamma }}`$ whereas bounds on $`m_{3/2}`$ scale as $`M_{\gamma \gamma }`$.
Another set of constraints is obtained from reactor experiments, where nuclear de-excitation is studied. Let us again make use of Eq. (17) and Eq. (18). Then we obtain for the isoscalar transition from excited nucleon state with the change of spin by $`J`$ and isospin by $`T`$ and with emission of photons and pseudoscalars with momenta $`k_\gamma `$ and $`k_P`$, respectively, the following ratio of rates
$$\frac{\omega _P^{J=1,T=0}}{\omega _\gamma ^{M1,J=0}}6\left(\frac{k_P}{k_\gamma }\right)^3\frac{g^{(0)2}g_{\pi NN}^2f_\pi ^2}{4\pi \alpha M_N^2},$$
where $`M1`$ refers to the type of electromagnetic transition and effective pion-nucleon coupling constant is $`g_{\pi NN}^2/4\pi =14.6`$. Products of pseudoscalar decay (two photons or $`e^+e^{}`$) are observed in detectors. In this way two constraints on the coupling of a pseudoscalar to nucleon have been obtained: $`\omega _P/\omega _\gamma \times \mathrm{Br}(Pe^+e^{})<10^{16}`$ and $`\omega _P/\omega _\gamma \times \mathrm{Br}(P\gamma \gamma )<1.510^{10}`$ (we set the pseudoscalar momentum equal to photon frequency, $`k_P=k_\gamma `$). Corresponding bounds on $`\sqrt{F}`$ are presented in Table 4 at $`A_Q=100`$ GeV and scale as $`\sqrt{A_Q}`$. The first constraint is valid for $`m_P<1.5`$ MeV and becomes weaker for heavier sgoldstinos, while the second limit is relevant only for light sgoldstino, $`m_P<1`$ MeV. The larger the branching ratio the stronger the corresponding bounds on $`\sqrt{F}`$: these bounds scale as quartic root of branching ratios. Although sgoldstino branching into $`e^+e^{}`$ is usually very small (see Figure 1), current experimental bounds on $`F`$ from sgoldstino decaying into $`e^+e^{}`$ are stronger than limits from decay into two photons. Note that reactor experiments give fairly strong bounds on $`F`$ but they should be considered as order-of magnitude estimates, as obtaining exact numbers requires accurate calculations involving nuclear matrix elements.
#### 3.1.3 Flavor conserving rare decays
Numerous bounds arise from precise measurements of partial widths of mesons and leptons (see Tables 6, 5, 7 and 8), if corresponding processes are allowed kinematically.
We begin with constraints independent of flavor violating terms in squark (slepton) mass matrix. One obtains a set of limits on supersymmetry breaking scale by considering Wilczek mechanism — decay of neutral vector meson $`V_{Q\overline{Q}}`$ ($`1^{}`$ state) into photon and (pseudo)scalar $`S(P)`$. There are two types of contributions into this process (see Fig. 2).
The first one is emission of real photons and (pseudo)scalars by quarks, while the second is decay of virtual photons, emitted by quarks, into photons and (pseudo)scalars. The first process is governed by fermion-sgoldstino coupling, while the second one emerges due to interaction with a pair of photons. The relevant candidates on the role of $`V_{Q\overline{Q}}`$ are $`J/\psi `$, $`\mathrm{{\rm Y}}`$ and $`\rho `$-, $`\omega `$-, $`\varphi `$-mesons.
Let us first consider heavy mesons, which may be described as quasistationary systems. With account of effective lagrangian (1) we obtain
$$\frac{\mathrm{\Gamma }(V_{Q\overline{Q}}S(P)\gamma )}{\mathrm{\Gamma }(V_{Q\overline{Q}}\gamma e^+e^{})}=\frac{M_V^2(A_QM_{\gamma \gamma })^2}{16\pi \alpha F^2},$$
(19)
where $`(+)`$ refers to decay into $`S(P)`$. We should compare the rate $`\mathrm{\Gamma }(V_{Q\overline{Q}}S(P)\gamma )`$ with current data on the rates $`\mathrm{\Gamma }(V_{Q\overline{Q}}\gamma +missingenergy)`$, $`\mathrm{\Gamma }(V_{Q\overline{Q}}3\gamma )`$ or $`\mathrm{\Gamma }(V_{Q\overline{Q}}\gamma +pair(s)ofleptons(lightmesons))`$ depending on $`m_{S(P)}`$ and superpartner mass spectrum (see discussion of sgoldstino decay modes in section 2). For illustration we set $`M_{\gamma \gamma }=A_Q=100`$ GeV in our quantitative estimates, so vector mesons would decay only into scalar sgoldstino. Eq. (19) shows that the corresponding constraints on $`\sqrt{F}`$ scale as a square root of the absolute value of the difference (sum) of $`A_Q`$ and $`M_{\gamma \gamma }`$, if one considers decay into scalar (pseudoscalar).
It turns out that constraints on $`F`$ from $`\mathrm{{\rm Y}}`$ decay into photons (leptons or light mesons), summarized in Table 5, are of the same order as limits from processes with single photon and missing energy (see Table 6) if corresponding branching ratios for sgoldstino decay are roughly of order one. The first type of constraints (Table 5) is relevant for sgoldstino decaying within detector (which is the case for $`m_{S(P)}10`$ MeV if $`M_{\gamma \gamma }=A=100`$ GeV); these constraints scale as quartic root of the corresponding sgoldstino branching ratios. The second type of bounds (Table 6) applies to lighter sgoldstino flying away from detector. We present in Table 5 only strongest constraints on $`\sqrt{F}`$. Besides these, there is a number of other $`\mathrm{{\rm Y}}`$ decay modes providing somewhat weaker constraints on $`F`$: $`\gamma \pi ^+\pi ^{}K^+K^{}`$, $`\gamma 2\pi ^+2\pi ^{}`$, $`\gamma 3\pi ^+3\pi ^{}`$, $`\gamma 2\pi ^+2\pi ^{}K^+K^{}`$.
One can show that limits on $`\sqrt{F}`$ from decays of light vector mesons ($`\rho `$, $`\omega `$, $`\varphi `$) are weaker at least by an order of magnitude.
Decays of $`\mathrm{{\rm Y}}`$ seem to have the best sensitivity to flavor-conserving sgoldstino couplings if $`M_\mathrm{{\rm Y}}m_{S(P)}afew`$ MeV. Since in the most part of the parameter space sgoldstino decays predominantly into two photons or two mesons, the most promising $`\mathrm{{\rm Y}}`$ decays are into three photons and into a photon and a pair of mesons. In models with large trilinear soft terms, leptonic widths of sgoldstinos become larger, and these modes become also interesting.
#### 3.1.4 Flavor violating rare decays
There is another type of processes to be considered. These are decays of charged particles: leptons or pseudoscalar mesons. The rates of these processes are more model dependent because they are governed by flavor violating soft terms.
While the bounds on $`\sqrt{F}`$ coming from decays of leptons are the same irrespectively of whether scalar or pseudoscalar sgoldstino is created in the final state, in the flavor-violating hadronic processes the creation of scalar sgoldstino is more important than the creation of pseudoscalar sgoldstino (if they have similar masses). When we discuss hadronic processes in what follows, we consider the emission of scalar sgoldstino only. The simplest example is kaon decay $`K^+\pi ^+S`$. In chiral theory kaon conversion into pion is described by matrix element
$$\pi ^+|\overline{s}\gamma _\mu d|K^+=\left(f_+(k_K+k_\pi )_\mu +f_{}(k_Kk_\pi )_\mu \right),$$
where $`f_+=1`$ and $`f_{}=0`$ in the case of exact $`SU(3)`$ flavor invariance. Then
$$\pi ^+|\overline{s}d|K^+=f_+\frac{m_K^2m_\pi ^2}{m_d+m_s}+f_{}\frac{m_S^2}{m_d+m_s}$$
and in what follows we neglect $`f_{}`$ contribution.
In principle, there are two mechanisms of the decay of charged particles with sgoldstinos in the final state. The first one is due to flavor-conserving sgoldstino interactions (with fermions and intermediate W-boson). The second one is due to flavor-changing terms in the low-energy effective interactions of light sgoldstinos (see Eqs.(1), (2); for instance, decay $`K^+\pi ^+S`$ is due to $`_{eff}=\frac{\stackrel{~}{m}_{D_{12}}^{LR2}}{\sqrt{2}F}S\overline{s}d`$). Hence, the second contribution emerges because of flavor violating interactions with fermions originating from off-diagonal insertions in squark(slepton) mass matrix.
As regards the first mechanism, it gives rise to constraints on $`\sqrt{F}`$ at the level of 100-250 GeV. We do not present these constraints explicitly, as they are at the same level or weaker than those summarized in Tables 5, 6.
The second mechanism is more interesting. The corresponding constraints are presented in Tables 7, 8,
where for definiteness we take flavor violating off-diagonal insertions in squark(slepton) mass matrix to be equal to their current experimental limits at $`\stackrel{~}{m}_{squark}=M_3=500`$ GeV, $`\stackrel{~}{m}_{slepton}=100`$ GeV. The limits on $`\sqrt{F}`$ scale as inverted quartic root of bounds on meson branchings and as square root of the off-diagonal elements $`m_{ij}^{LR2}`$ in squark squared mass matrix; they depend crucially on the strength of flavor violation in MSSM. Since hadronic and photonic modes usually dominate, limits on $`\sqrt{F}`$ coming from meson decays with a pair of leptons in the final state (say, $`K^+\pi ^+S(Se^+e^{})`$) are weaker, but not more than by one or two orders of magnitude, as compared to photonic and mesonic modes. Note, that similar constraints from three-body decays of neutral mesons (like $`B^0K^0S(S\mu ^+\mu ^{})`$) depend on the same coupling constants and are generally weaker than limits from rare decays of charged mesons.
From bounds presented in this section we conclude that sgoldstino interactions may give large contributions into flavor changing rare decays, including those forbidden in SM. In particular, in the case $`F=1`$ TeV<sup>2</sup>, the constraints from processes with final light sgoldstino significantly strengthen the bounds on off-diagonal elements in squark and slepton mass matrices in comparison with models where sgoldstinos decouple at low energies.
Our analysis suggests that contributions of intermediate (virtual) sgoldstinos into FCNC and lepton flavor violating processes may be also significant. For instance, pseudoscalar mesons may decay through light sgoldstino exchange. Also, there are potentially important contributions to loop processes like $`K^0\overline{K}^0`$, $`B^0\overline{B}^0`$ mixings, etc. These issues will be considered elsewhere.
### 3.2 Processes with two sgoldstinos
Processes with two final sgoldstinos appear due to the presence of two-sgoldstino interactions in low-energy effective lagrangian, Eq. (2), and due to the double contribution of one-sgoldstino interaction, Eq. (1). Of course, the corresponding amplitudes are highly suppressed (by additional $`F^1`$). Nevertheless, some of these processes are sensitive enough to place constraints on $`\sqrt{F}`$ at the level of 1 TeV.
Recall that two-sgoldstino coupling constants (2) differ from one-sgoldstino constants (1). Indeed, they are proportional to $`LL`$ and $`RR`$ insertions in scalar squared mass matrix, while one-sgoldstino coupling constants are proportional to $`LR`$ insertions. Note in this regard, that the current limits on flavor changing squark masses $`m_{ij}^{LL2}`$ and $`m_{ij}^{RR2}`$ are weaker than limits on $`m_{ij}^{LR2}`$. Hence, it makes sense to consider processes where two-sgoldstino couplings could be observed.
Complete analysis of low-energy processes with two sgoldstinos may be carried out along the same lines as for processes with one sgoldstino. Instead of going through the limits systematically, we discuss here only some examples in order to get the feeling of sensitivity to $`\sqrt{F}`$.
#### 3.2.1 Light neutral mesons
We begin with pion decay into two light sgoldstinos (see Fig. 3a). The relevant part of the effective lagrangian reads
$$=\frac{1}{4F^2}\left(S_\mu PP_\mu S\right)\left(\left(\stackrel{~}{m}_{U_{11}}^{LL2}+\stackrel{~}{m}_{U_{11}}^{RR2}\right)\overline{u}\gamma ^\mu \gamma ^5u+\left(\stackrel{~}{m}_{D_{11}}^{LL2}+\stackrel{~}{m}_{D_{11}}^{RR2}\right)\overline{d}\gamma ^\mu \gamma ^5d\right).$$
Then by making use of Eqs. (9), (10) we obtain
$$\mathrm{\Gamma }(\pi ^0SP)=\frac{f_\pi ^2}{m_\pi }\frac{[\stackrel{~}{m}_{U_{11}}^{LL2}+\stackrel{~}{m}_{U_{11}}^{RR2}\stackrel{~}{m}_{D_{11}}^{LL2}\stackrel{~}{m}_{D_{11}}^{RR2}]^2}{128\pi F^2}\frac{\left(m_S^2m_P^2\right)^2}{F^2}\sqrt{\left(1+\frac{m_P^2m_S^2}{m_\pi ^2}\right)^24\frac{m_P^2}{m_\pi ^2}}.$$
(20)
This rate is proportional to $`(m_S^2m_P^2)`$, so, as expected, it vanishes in the massless limit, $`m_S,m_P0`$.
In order to examine the sensitivity of this process, let us neglect the phase volume dependence and take $`|m_S^2m_P^2|m_\pi ^2/4`$. If we set the value in the square bracket equal to $`2\stackrel{~}{m}_Q^2`$ and choose $`\stackrel{~}{m}_Q=500`$ GeV, we obtain the limits presented in Table 9. A few remarks are in order. First, these bounds may be irrelevant in some theories because $`\sqrt{F}`$ should not be significantly smaller than any of the soft terms. Second, the constraint from pion disappearance (i.e., from Br($`\pi ^0SP`$)) is valid only if sgoldstinos fly away from detector. For $`m_{S(P)}m_\pi /2`$ this is the case if $`M_{\gamma \gamma },A_e<10`$ GeV, which is not forbidden by current experiments. Third, these limits are obtained at tuned sgoldstino masses and, in general, they are weaker (see Eq. (20)).
These results do not depend on flavor-violating couplings and are of the same order of magnitude as the limits presented in Tables 6, 5. However, the limits presented in Table 9 scale as inverted octopic root of the corresponding pion partial width (see Eq. (20)).
To illustrate that two-sgoldstino processes may impose more stringent constraints than one-sgoldstino processes with the same content of final SM particles, let us estimate the one-sgoldstino contribution to pion decay into four photons. Namely let us consider emission of sgoldstinos from the photon legs of pion (see Figure 3b). If sgoldstino decays within detector into photons, this would correspond to four-photon decay of pion. Pion-photon anomalous amplitude reads
$$A(\pi \gamma \gamma )=\frac{\alpha }{\pi f_\pi }ϵ_{\mu \nu \rho \sigma }ϵ_1^\mu ϵ_2^\nu q_3^\rho q_2^\sigma ,$$
where $`q_1`$, $`q_2`$ are the photon momenta. Then the corresponding squared matrix element of $`\pi ^0\gamma \gamma S(P)`$ is
$$|M|^2=4\frac{\alpha ^2}{\pi ^2}\frac{M_{\gamma \gamma }^2f_\pi ^2}{F^2}\left((q_1p)^2+(q_1q_3)^2\right),$$
where $`p`$ and $`q_1`$, $`q_3`$ are momenta of sgoldstino and outgoing photons, respectively. Neglecting sgoldstino mass we estimate the decay width as
$$\mathrm{\Gamma }(\pi ^0\gamma \gamma S(P))=\frac{1}{32}\frac{\alpha ^2m_\pi ^3}{f_\pi ^2\pi ^4}\frac{M_{\gamma \gamma }^2m_\pi ^2}{F^2}.$$
(21)
One can check that Eq. (21) gives weaker bound on $`\sqrt{F}`$ than the limit presented in the second row of Table 9 if $`M_{\gamma \gamma }=100`$ GeV and $`\stackrel{~}{m}_Q=500`$ GeV.
Let us now evaluate the bounds from decays of neutral kaons due to two-sgoldstino flavor violating couplings. The effective lagrangian reads
$`={\displaystyle \frac{1}{4F^2}}(S_\mu PP_\mu S)((\stackrel{~}{m}_{D_{21}}^{LL2}+\stackrel{~}{m}_{D_{21}}^{RR2})\overline{s}\gamma ^\mu \gamma ^5d+(\stackrel{~}{m}_{D_{21}}^{LL2}\stackrel{~}{m}_{D_{21}}^{RR2})\overline{s}\gamma ^\mu d`$
$`+(\stackrel{~}{m}_{D_{12}}^{LL2}+\stackrel{~}{m}_{D_{12}}^{RR2})\overline{d}\gamma ^\mu \gamma ^5s+(\stackrel{~}{m}_{D_{12}}^{LL2}\stackrel{~}{m}_{D_{12}}^{RR2})\overline{d}\gamma ^\mu s).`$
One can show, that only the measurements of branching ratios of $`K_L^0`$ impose interesting constraints on $`F`$, whereas current limits on rare $`K_S^0`$ decays provide weak constraints on $`\sqrt{F}`$. We obtain by making use of chiral theory
$$\mathrm{\Gamma }(K_L^0SP)=\frac{f_K^2}{m_K}\frac{[\stackrel{~}{m}_{D_{21}}^{LL2}+\stackrel{~}{m}_{D_{12}}^{LL2}+\stackrel{~}{m}_{D_{21}}^{RR2}+\stackrel{~}{m}_{D_{12}}^{RR2}]^2}{512\pi F^2}\frac{\left(m_S^2m_P^2\right)^2}{F^2}\sqrt{\left(1+\frac{m_P^2m_S^2}{m_K^2}\right)^24\frac{m_P^2}{m_K^2}}.$$
Note that in the limit of CP conservation, $`LL`$ and $`RR`$ squark mass matrices are real and symmetric, and the sum in the bracket equal to $`2\left(\stackrel{~}{m}_{D_{21}}^{LL2}+\stackrel{~}{m}_{D_{21}}^{RR2}\right)`$.
In analogy to the discussion of pion decays, let us neglect the phase volume dependence and set $`|m_S^2m_P^2|m_K^2/4`$, $`f_K=160`$ MeV. If we set the sum in the bracket equal to $`4\mathrm{R}\mathrm{e}\stackrel{~}{m}_{D_{21}}^{LL2}`$ and impose on $`\mathrm{Re}\stackrel{~}{m}_{D_{21}}^{LL2}`$ current constraints from the absence of FCNC at squark mass $`\stackrel{~}{m}_Q=500`$ GeV, we obtain the limits presented in Table 10.
Note that limits from kaon decays into a leptonic pair and a pair of mesons(photons) are usually more significant than limits from decays into four leptons, because of small sgoldstino decay branching ratio into leptons (see Figure 1). These bounds on $`\sqrt{F}`$ are obtained at tuned sgoldstino masses, $`|m_S^2m_P^2|m_K^2/4`$, and generally the bounds are somewhat weaker.
We are not aware of limits on decays $`K_L^04\gamma `$ and $`K_L^0\pi ^+\pi ^{}\gamma \gamma `$. If it would be possible to measure (or limit) their branching ratios at the level of $`10^7`$, the sensitivity of experiments to two-sgoldstino couplings would increase, because the photonic decay usually dominates over leptonic decay of sgoldstinos.
#### 3.2.2 Decays of heavy mesons
In analogy with light mesons we consider now heavy neutral mesons living sufficiently long, $`D^0`$, $`B^0`$ and $`B_s^0`$. We make use of the approach similar to the chiral theory in order to describe their interaction with sgoldstinos; in the following we set $`f_{B_s^0}=f_{B^0}=f_{D^0}`$=200 MeV. The limits obtained with the same assumptions as above about sgoldstino masses and values of flavor violating couplings are listed in Table 11. All remarks concerning these assumptions given in previous subsection may be repeated here.
It turns out that constraints on two-sgoldstino coupling from current bounds on $`B_s`$ decay modes are rather weak.
To conclude this subsection, we note that “two-sgoldstino” bounds scale as inverted octopic root of meson branching ratios, so eight order improvements (!) in measurements of corresponding partial widths are required to strengthen the bounds on $`\sqrt{F}`$ by an order of magnitude.
## 4 Discussions and Conclusions
We have considered models with light ($`m_{S(P)}<afew`$ GeV) superpartners of goldstino. The constraints on their effective couplings to SM particles have been presented. By making use of these constraints we have estimated the sensitivity of low energy experiments to the scale of supersymmetry breaking and gravitino mass.
Let us compare our results with constraints coming from high energy processes . If we ignore sgoldstino, then in models with light gravitino current direct bound on supersymmetry breaking scale is obtained from Tevatron, $`\sqrt{F}<`$217 GeV . The upgraded Tevatron may be able to cover the range of $`\sqrt{F}`$ up to $`\sqrt{F}`$290 GeV , while LHC will be capable to reach $`\sqrt{F}<`$1.6 TeV . In models with light sgoldstinos, collider experiments become more sensitive to the scale of sypersymmetry breaking. Most powerful among the operating machines, LEP and Tevatron, give a constraint of order 1 TeV on supersymmetry breaking scale in models with light sgoldstinos. Indeed, it was found in Ref. that with the LEP luminosity of 100 pb<sup>-1</sup>, at $`\sqrt{F}=1÷1.5`$ TeV one $`e^+e^{}SZ(\gamma )`$ event would occur, and ten $`e^+e^{}e^+e^{}S`$ events would appear at $`\sqrt{F}=1.5`$ TeV. At Tevatron, about 10 events in $`ppS\gamma (Z)`$ channel, and $`10^5`$ events in $`ppS`$ channel would be produced at $`\sqrt{F}=1`$ TeV . This gives rise to a possibility to detect sgoldstino, if it decays inside detector into photons and $`\sqrt{F}`$ is not larger than $`1.5÷2`$ TeV. However, these numbers have been obtained in a model with heavier superpartner scale, and, hence, with larger sgoldstino couplings, than we assumed in our paper. For that set of parameters bounds on $`\sqrt{F}`$ derived in our paper from processes originating due to flavor-conserved sgoldstino couplings should be stronger at least by a factor of 1.5 .
One important remark concerns the sensitivity of collider experiments to light particles. The currently available analyses carried out by experimental collaborations are relevant only for fairly heavy sgoldstino ($`M20`$ GeV). In this paper we have discussed lighter particles; in this sense our results may be considered as complementary to those derived up to now from LEP and Tevatron experiments.
From constraints presented in this paper we conclude that the sgoldstino signal is not likely to be observed at LEP and Tevatron if sgoldstinos are lighter than a few GeV and flavor-violating processes in MSSM are not extremely suppressed. One observes that precision measurements at low energies are promising for confirming directly such a model. The astrophysical bounds are usually stronger than laboratory ones, though they become invalid for $`m_S`$ and $`m_P`$ larger than a few MeV.
Among the laboratory processes, the most sensitive to very light sgoldstinos are propagation of laser beam in magnetic field and reactor experiments. For heavier sgoldstinos the most sensitive processes are $`\mathrm{{\rm Y}}`$ decays for flavor-conserving sgoldstino couplings and charged kaon decays for flavor-violating sgoldstino couplings.
Acknowledgments
The author is indebted to F. Bezrukov, P. Onyisi, A. Ovchinnikov, A. Penin, A. Pivovarov and V. Rubakov for useful discussions. This work was supported in part by RFBR grant 99-01-18410, CRDF grant 6603, Swiss Science Foundation, grant 7SUPJ062239, Russian Academy of Science, JRP grant # 37 and by ISSEP fellowship.
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# The Channel Capacity of a Fiber Optics Communication System: perturbation theory
## I Introduction
The performance of any communication system is ultimately limited by the signal to noise ratio of the received signal and available bandwidth. This limitation can be stated more formally by using the concept of channel capacity introduced withing the framework of information theory. The channel capacity is defined as the maximum possible bit rate for error-free transmission in the presence of noise. For a linear communication channel with additive white Gaussian noise, and a total signal power constraint at the input, the capacity is given by the celebrated Shannon formula
$`C`$ $`=`$ $`W\mathrm{log}\left(1+{\displaystyle \frac{P_0}{P_N}}\right)`$ (1)
where $`W`$ is the channel bandwidth, $`P_0`$ is the average signal power, and $`P_N`$ is the average noise power.
Current optical fiber systems operate substantially below the fundamental limitation, imposed by the Eq. (1). However, a considerable improvement in the coding schemes for lightwave communications, expected in the near future, may result in the development of systems, whose efficiency may approach this fundamental limit.
However, the representation of the channel capacity in the standard form (1) is unsuitable for applications to the actual fiber optics systems. It was obtained based on the assumption of linearity of the communication channel, while the modern fiber optics systems operate in a substantially nonlinear regime. Since the optical transmission lines must satisfy very strict requirements for bit-error-rate ($`10^{12}`$ to $`10^{15}`$), the pulse amplitude should be large enough so that is can be effectively detectable. The increase of the number of wavelength-division multiplexing (WDM) channels in the modern fiber optics communication systems also leads to a substantial increase of the electric field intensity in the fiber. As a consequence, the Kerr nonlinearity of the fiber refractive index $`n=n_0+\gamma I`$ (where $`I`$ is the pulse intensity) becomes substantial and should be taken into account.
In the present paper we consider corrections to the channel capacity of the optical fiber communication system, originating from the nonlinearity of the fiber. The technique that we use involves a perturbative computation of the relevant mutual information and subsequent optimization. To our knowledge, this method appears to be substantially new.
## II Fiber Optics Communication System as an Information Channel
We consider a typical fiber optics communication system, which consists of a sequence of $`N`$ fibers each followed by an amplifier (see Fig. 1). The amplifiers have to be introduced in order to compensate for the power loss in the fiber. An inevitable consequence of such design, however, is the generation of the noise in the system, coming from the spontaneous emission in the optical amplifiers. For simplicity, we will assume that all the fibers and the amplifiers of the link are identical.
The information is encoded in the electric field at the “imput” of the system, typically using the light pulses sent at different frequencies. The available bandwidth of the amplifiers as well as the increase of the fiber absorption away from the “transparency window” near the wavelength $`\lambda =1.55\mu m`$, limits the bandwidth of the fiber optic communication system.
The maximum amount of the information, that can be transmitted through the communication system per unit time, is called the channel capacity $`C`$. According to the Shannon’s basic result, this quantity is given by the maximum value of the mutual information per second over all possible input distributions:
$`C`$ $`=`$ $`\mathrm{max}_{p_x}\left\{H\left[y\right]H\left[y|x\right]_{p_x}\right\}`$ (2)
The mutual information
$`R`$ $`=`$ $`H\left[y\left(\omega \right)\right]H\left[y\left(\omega \right)|x\left(\omega \right)\right]_{p_x}`$ (3)
is a functional of the “input distribution” $`p_x\left[x\left(\omega \right)\right]`$, which represents the encoding of the information using the electric field components at different frequences
$`E_{\mathrm{in}}\left(t\right)`$ $`=`$ $`{\displaystyle _W}𝑑\omega x\left(\omega \right)\mathrm{exp}\left(i\omega t\right)`$ (4)
The entropy $`H\left[y\left(\omega \right)\right]`$ is the measure of the information received at the output of the communication channel. However, if the channel is noisy, for any output signal there is some uncertainty of what was originally sent. The conditional entropy $`H\left[y\left(\omega \right)|x\left(\omega \right)\right]`$ at the output for a given $`x\left(\omega \right)`$ represents this uncertainty.
The entropies $`H\left[y\left(\omega \right)\right]`$ and $`H\left[y\left(\omega \right)|x\left(\omega \right)\right]`$ are defined in terms of the corresponding distributions $`p\left(y\right)`$ and $`p\left(y|x\right)`$ via the standard relation
$`H`$ $``$ $`{\displaystyle 𝒟y\left(\omega \right)p\left[y\left(\omega \right)\right]\mathrm{log}\left(p\left[y\left(\omega \right)\right]\right)}`$ (5)
where $`pp_y(y)`$ for the entropy $`H\left[y\left(\omega \right)\right]`$, and $`pp\left(y|x\right)`$ for the entropy $`H\left[y\left(\omega \right)|x\left(\omega \right)\right]`$, and the functional integral is defined in the standard way
$`{\displaystyle 𝒟\xi \left(\omega \right)}\underset{M\mathrm{}}{lim}c_M\left[\mathrm{\Pi }_{m=1}^M{\displaystyle 𝑑\xi \left(\omega _m\right)}\right]`$ (6)
where is a normalization constant.
For any communication link, the signal power is limited by the system hardware. Therefore, the maximum of the mutual information in (1) should be found under the constraint of the fixed total power $`P_0`$ at the input:
$`P_0`$ $`=`$ $`{\displaystyle 𝒟x\left(\omega \right)\left|x\left(\omega \right)\right|^2p_x\left[x\left(\omega \right)\right]}`$ (7)
If the propagation in the communication channel is described by a linear equation, then the input-output relation for the system is given by
$`y\left(\omega \right)=K\left(\omega \right)x\left(\omega \right)+n\left(\omega \right)`$ (8)
where $`n\left(\omega \right)`$ is the noise in the channel. In this approximation, the problem of finding the maximum of the mutual information can be solved exactly, with the corresponding input distribution $`p_x`$ being Gaussian. If the amplifiers compensate exactly for the power losses in the fibers, the Channel Capacity is given by the Shannon formula (1).
As follows from (1), the better bit rates can be obtained for the higher signal-to-noise ratio $`P_0/P_N`$. With this in mind, the optics fiber communication systems are designed to operate with the pulses of high power. As a result, the optics fiber links operate in the regime, in which due to the Kerr nonlinearity the refraction index of the fiber strongly depends on the local electric field intensity. Therefore, a modern fiber optics communication system is, in fact, an essentially nonlinear communication channel, and cannot be adequately described within the framework of the Shannon’s linear theory.
## III The Model
The first step in the calculation of the channel capacity is to find the “input-output” relation for the communication channel. The time evolution of the electric field in the fiber $`E(z,t)`$, where $`z`$ is the distance along the fiber, can be accuarately described in the “envelope approximation”, when
$`E(z,t)`$ $`=`$ $`A(z,t)\mathrm{exp}\left(i\left(\beta _0z\omega _0t\right)\right)+\mathrm{c}.\mathrm{c}.`$ (9)
where the function $`A`$ represents the slowly (compared to the light frequency) varying amplitude of the electric field
in the fiber. The evolution of $`A(z,t)`$ is described by the equation
$`{\displaystyle \frac{A}{z}}+\beta _1{\displaystyle \frac{A}{t}}+{\displaystyle \frac{i}{2}}\beta _2{\displaystyle \frac{^2A}{t^2}}+{\displaystyle \frac{\alpha }{2}}A=i\gamma \left|A\right|^2A`$ (10)
Here the coefficients $`\beta `$ describe the frequency dependence of the wavenumber
$`\beta \left(\omega \right)`$ $`=`$ $`\beta _0+\beta _1\omega +{\displaystyle \frac{\beta _2}{2}}\omega ^2+O\left[\omega ^3\right]`$ (11)
where $`\omega `$ is measured from the center of the band $`\omega _0`$.
The equation (10) neglects the effects such as the stimulated Raman scattering and the stimulated Brillouin scattering, compared to the Kerr nonlinearity of the refraction index of the fiber, represented by the term $`\gamma \left|A\right|^2A`$.
The optical amplifiers incorporated into the communication system (see Fig. 1) compensate for the power losses in the fiber, but due to spontaneous emission each of them will inevitably introduce noise $`n(t)=\mathrm{exp}\left(i\omega _0t\right)𝑑\omega n_\omega \mathrm{exp}\left(i\omega t\right)`$ into the channel. Generally, even in a single optical amplifier, the noise distribution at any given frequency $`\omega _0+\omega `$ within the channel bandwidth $`n\left(\omega \right)`$ is close to a Gaussian:
$`p_n\left[n\left(\omega \right)\right]`$ $``$ $`\mathrm{exp}\left[{\displaystyle \frac{\left|n\left(\omega \right)\right|^2}{P_N^\omega }}\right]`$ (12)
This is even more so in a system with many independent amplifiers, due to the Central Limit Theorem. For simplicity, the noise spectrum $`P_N^\omega `$ can be assumed to be flat.
If the envelope function just before the amplifier is $`A\left(t\right)A_\omega ^0\mathrm{exp}\left(i\omega t\right)`$, then immediately after the amplifier
$`A_\omega `$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{\alpha }{2}}d\right)A_\omega ^0+n_\omega `$ (13)
where $`d`$ is the span of a single fiber.
The equations (10), (13) define the evolution of the electric field envelope over one “fiber-aplifier” link of the communication system. The total “input-output” relation will then involve solving the corresponding equations for all $`N`$ iterations of the single fiber-amplifier unit.
## IV The Perturbative Framework
If one is able to calculate the “output” signal $`y(\omega )`$ in terms of the “input” $`x(\omega )`$ and the noise contributions of each of the amplifiers $`n_\omega ^{\left\{\alpha \right\}},\alpha =1,\mathrm{},N`$,
$`y(\omega )=\mathrm{\Phi }(x(\omega );n_{\omega }^{}{}_{}{}^{\left\{1\right\}},\mathrm{},n_{\omega }^{}{}_{}{}^{\left\{N\right\}})`$ (14)
then the the conditional distribution $`p\left(y|x\right)`$ can be simply calculated as follows:
$`p\left(y|x\right)`$ $`=`$ $`\left\{\mathrm{\Pi }_{\alpha =1}^{N1}{\displaystyle 𝒟n_{\omega }^{}{}_{}{}^{\left\{\alpha \right\}}p_n\left[n_{\omega }^{}{}_{}{}^{\left\{\alpha \right\}}\right]}\right\}`$ (15)
$`\times `$ $`p_n\left[y(\omega )\mathrm{\Phi }(x(\omega );n_{\omega }^{}{}_{}{}^{\left\{1\right\}},\mathrm{},n_{\omega }^{}{}_{}{}^{\left\{N\right\}})\right]`$
where $`p_n`$ is the distribution function of the noise, produced by a single amplifier. The output distribution $`p_y\left(y\right)`$ can then be directly related to the input distribution $`p_x\left(x\right)`$ via the standard relation
$`p_y\left(y\right)={\displaystyle 𝒟x(\omega )p\left[y(\omega )|x(\omega )\right]p_x\left[x(\omega )\right]}`$ (16)
Using Eqns. (15),(16), one is able to express the mutual information in terms of a single distribution $`p_x`$. The calculation of the channel capacity then reduces to a standard problem of finding the maximum of a (nontrivial) functional.
The equation (10) is, in fact, the well studied nonlinear Shroedinger equation, with the time and distance variables interchanged. Only some partial solutions of this equation are known, corresponding to solitons. However, in order to calculate the channel capacity, one needs to find the general input-output relation for the communication system. This implies solving a set of $`N`$ essentially nonlinear equations (10) for arbitrary initial conditions. Even knowing some partial solutions, doing such calculation exactly for an essentially nonlinear system is not possible in a closed form.
In order to make progress, we note the presence of a natural perturbation parameter in the problem, namely $`\gamma `$.In fact, the fiber equation (10) is already an approximation, derived in the limit, when the change in the effective refraction index due to pulse propagation, described by the nonlinear term $`i\gamma \left|A\right|^2A`$, is small compared to the “unperturbed” value of the index of refraction $`n_0`$. We have developed a perturbative technique, when the solution of the nonlinear evolution equation, is represented as a power series in $`\gamma `$. Solving (10) separately for each power of $`\gamma `$, and using (13), for the input-output relation of a single fiber-amplifier unit $`\mathrm{\Phi }_\omega ^{(n)}`$, defined as
$`A_\omega ^{(n)}`$ $`=`$ $`\mathrm{\Phi }_\omega ^{(n)}\left(A_\omega ^{(n1)}\right),`$ (17)
we obtain:
$`\mathrm{\Phi }_\omega ^{(n)}\left(A_{\omega }^{}{}_{}{}^{(n1)}\right)`$ $`=`$ $`\left[A_\omega ^{\left(n1\right)}+{\displaystyle \underset{\mathrm{}=1}{\overset{\mathrm{}}{}}}\gamma ^{\mathrm{}}_\omega ^{(\mathrm{})}\left(A_\omega ^{\left(n1\right)}\right)\right]`$ (18)
$`\times `$ $`\mathrm{exp}\left(i\kappa _\omega d\right)+n_\omega `$
where $`d`$ is the length of a single fiber,
$`\kappa _\omega `$ $`=`$ $`\beta _1\omega {\displaystyle \frac{1}{2}}\beta _2\omega ^2`$ (19)
The procedure for the calculation of the functions $`_\omega ^{(\mathrm{})}`$, described in detail Appendix A, can be carried to an arbitrary order $`\mathrm{}`$.
The further calculation then involves the following steps:
* Iterating Eq. (18) $`N`$ times, to obtain the “input-output” relation for the whole communication system $`\mathrm{\Phi }_\omega [x\left(\omega \right);n_\omega ^{(1)},\mathrm{},n_\omega ^{(N)}]`$
* substituting the result into Eqns. (15), (16) to obtain the conditional distribution $`p(x|y)`$ and the output distribution $`p_y(y)`$ in terms of the input distribution $`p_x(x)`$ as expansions in powers of $`\gamma `$
* calculating the entropies $`H\left[y\left(\omega \right)\right]`$ and $`H\left[x\left(\omega \right)|y\left(\omega \right)\right]`$, and the mutual information $`R`$
Following these steps, the calculation of the channel capacity becomes a straightforward procedure. In Appendix B we describe it in detail, using a simple nonlinear channel $`y\left(\omega \right)=x\left(\omega \right)\mathrm{exp}\left(\varphi \left[x\left(\omega \right)\right]\right)+n\left(\omega \right)`$ as an example.
## V The Fiber Link Channel Capacity
After a tedious, but straigthforward calculation, we obtain:
$`H\left[y\left(\omega \right)\right]`$ $`=`$ $`H_0\left[y\left(\omega \right)\right]\mathrm{\Delta }C_1\mathrm{\Delta }H_y+𝒪\left(\gamma ^4\right)`$ (20)
and
$`H\left[y\left(\omega \right)|x\left(\omega \right)\right]`$ $`=`$ $`H_0\left[y\left(\omega \right)|x\left(\omega \right)\right]+\mathrm{\Delta }C_2+𝒪\left(\gamma ^4\right)`$ (21)
where $`H_0\left[y\right]`$ and $`H_0\left[x|y\right]`$ are given by the standard expressions for a linear channel. In the limit of large signal-to-noise ratio $`P_0P_N`$ we obtain:
$`\mathrm{\Delta }C_1`$ $`=`$ $`N^2W\left({\displaystyle \frac{\gamma P_0}{\alpha }}\right)^2Q_1(\alpha d,{\displaystyle \frac{\beta _2^2W^4}{\alpha ^2}})`$ (22)
$`\mathrm{\Delta }C_2`$ $`=`$ $`{\displaystyle \frac{4}{3}}\left(N^21\right)W\left({\displaystyle \frac{\gamma P_0}{\alpha }}\right)^2Q_2(\alpha d,{\displaystyle \frac{\beta _2^2W^4}{\alpha ^2}})`$ (23)
and the functions $`Q_1`$ and $`Q_2`$ are defined as follows:
$`Q_1(u,z)`$ $`=`$ $`{\displaystyle _{1/2}^{1/2}}𝑑x_1{\displaystyle _{1/2}^{1/2}}𝑑x_2{\displaystyle _{1/2}^{\overline{x}}}𝑑x`$ (24)
$`\times `$ $`f(u,z;x_1,x_2,x)`$
$`Q_2(u,z)`$ $`=`$ $`{\displaystyle _{1/2}^{1/2}}𝑑x_1{\displaystyle _{1/2}^{1/2}}𝑑x_2{\displaystyle _{1/2}^{1/2}}𝑑x`$ (25)
$`\times `$ $`f(u,z;x_1,x_2,x)`$
where $`\overline{x}\mathrm{max}[1/2,1/2+x_1+x_2]`$, and
$`f(u,z;x_1,x_2,x){\displaystyle \frac{\left|1\mathrm{exp}\left(uiv^2\right)\right|^2}{1+v^2}}`$ (26)
where
$`vz\left(xx_1\right)\left(xx_2\right)`$ (27)
The correction
$`\mathrm{\Delta }H_y`$ $`=`$ $`\gamma ^2W{\displaystyle 𝒟y\left(\omega \right)p_y^{(0)}\left[y\left(\omega \right)\right]\left(p_y^{(1)}\left[y\left(\omega \right)\right]\right)^2}`$ (28)
is caused by the deviations of the otput distribution
$`p_y`$ $`=`$ $`p_y^{(0)}+{\displaystyle \underset{\mathrm{}}{}}\gamma ^{\mathrm{}}p_y^{(\mathrm{})}`$ (29)
from the Gaussian form
$`p_y^{(0)}\left[y\left(\omega \right)\right]\mathrm{exp}\left({\displaystyle \frac{\left|y\left(\omega \right)\right|^2}{P_\omega +P_\omega ^N}}\right)`$ (30)
where $`P_\omega `$ (such that $`P_0=𝑑\omega P_\omega `$) is the input power at frequency $`\omega `$:
$`p_x^{(0)}\left[x\left(\omega \right)\right]\mathrm{exp}\left({\displaystyle \frac{\left|x\left(\omega \right)\right|^2}{P_\omega ^0}}\right)`$ (31)
Note, that the correction $`\mathrm{\Delta }H_y0`$ and equals to zero only when the distribution $`p_y^{(1)}=0`$. Therefore, as follows from Eqns. (20),(28), in the second order in nonlinearity $`\gamma `$ the mutual information $`R`$ has the maximum, when $`p_y^{(1)}=0`$, or, equivalently, when the output distribution is Gaussian up to the first order in nonlinearity.
For a general nonlinear channel, that would correspond to the input distribution, being non-Gassian already in the first order in $`\gamma `$. The corresponding correction can be obtained from Eq. (16), taken only up to the first orded in nonlinearity:
$`{\displaystyle \frac{}{\gamma }}\left[{\displaystyle 𝒟x(\omega )p\left[y(\omega )|x(\omega )\right]p_x\left[x(\omega )\right]}\right]|_{\gamma =0}`$ $`=`$ $`0`$ (32)
Generally, such an integral would yield $`p_x\left(x\right)=p_x^{(0)}\left(x\right)\left(1+\gamma p_x^{(0)}\left(x\right)\right)`$, where $`p_x^{(0)}`$ is Gaussian, and $`p_x^{(1)}0`$. However, it is straightforward to show, that for the fiber optics channel described by Eq. (10), a Gaussian input distribution leads to non-Gaussian corrections in the output distribution starting only from the second order. Therefore, the requrement $`p_y^{(1)}=0`$ is satisfied, when the input distribution is such that $`p_x^{(1)}=0`$.
For the channel capacity, defined as the maximum value of the mutual information, we obtain:
$`C`$ $`=`$ $`W\mathrm{log}\left(1+{\displaystyle \frac{P_0}{P_N}}\right)\mathrm{\Delta }C_1\mathrm{\Delta }C_2+𝒪\left(\gamma ^4\right)`$ (33)
The equation (33) yields the result for the fiber optics channel capacity in the second order in the nonlinearity $`\gamma `$. In the next section we will discuss the physical origins of the corrections $`\mathrm{\Delta }C_1`$ and $`\mathrm{\Delta }C_2`$.
## VI The Discussion and the Conclusions
In the spirit of the Shannon formula, the decrease of the capacity of a communication channel with a fixed bandwidth can be attributed to (i) the effective suppression of the signal power, and (ii) the enhancement of the noise. The corrections to the channel capacity, derived in the previous section, can be interpreted as resulting precisely from these two effects.
The four-wave scattering, induced by the fiber nonlinearity, inevitably leads to the processes, which generate photons with the frequencies outside the channel bandwidth. Such photons, are not recorded by the “receiver”, and are lost for the purpose of the information transmission. This corresponds to an effective bandwidth power dissipation, and should therefore lead to a decrease of the channel capacity. Since for small nonlinearity this power loss $`\mathrm{\Delta }P\gamma ^2`$, the dimension analysis implies
$`\mathrm{\Delta }P`$ $`=`$ $`{\displaystyle \frac{\gamma ^2P_0^2}{\alpha ^2}}F\left({\displaystyle \frac{\beta W^2}{\alpha }}\right)`$ (34)
Such scattering processes are suppressed, when the scattering leads to a subsantial change of the total momentum $`\delta \kappa [\omega _1,\omega _2\omega _3,\omega ]`$, so that the corresponding scattering rate
$`S[\omega _1,\omega _2\omega _3,\omega ]`$ $``$ $`{\displaystyle \frac{\delta \left(\omega _1+\omega _2\omega _3\omega \right)}{1+\left(\delta \kappa /\kappa _0\right)^2}}`$ (35)
In the spirit of the uncertainty realtion, $`\kappa _01/L_{\mathrm{eff}}`$, where $`L_{\mathrm{eff}}`$ corresponds to the length of the concentration of the power of the signal in the fiber. For a small absorption coefficient $`\alpha 1/d`$ the distance $`L_{\mathrm{eff}}`$ is of the order of the fiber length $`d`$, while in the opposite limit $`\alpha 1`$ the effective length $`L_{\mathrm{eff}}1/\alpha `$.
Using Eqn. (19), and the energy conservation $`\omega _3=\omega _1+\omega _2\omega `$, the momentum change $`\delta \kappa [\omega _1,\omega _2\omega _3,\omega ]`$ can be expressed as
$`\delta \kappa `$ $`=`$ $`\beta _2\left(\omega \omega _1\right)\left(\omega \omega _1\right)`$ (36)
Substituting (36) into (35), for the channel capacity loss due to the bandwidth power “leakage”, in the limit $`P_0P_N`$, and $`\alpha d1`$, we obtain
$`\mathrm{\Delta }C_P`$ $``$ $`W{\displaystyle \frac{\mathrm{\Delta }P}{P}}W{\displaystyle \frac{\gamma ^2P_0^2}{\alpha ^2}}{\displaystyle _W}𝑑\omega _1{\displaystyle _W}𝑑\omega _2{\displaystyle _W}𝑑\omega _3`$ (37)
$`\times `$ $`{\displaystyle _{\omega W}}d\omega S[\omega _1,\omega _2\omega _3,\omega ]`$
$`=`$ $`W{\displaystyle \frac{\gamma ^2P_0^2}{\alpha ^2}}{\displaystyle _W}𝑑\omega _1{\displaystyle _W}𝑑\omega _2{\displaystyle _{\omega W}}𝑑\omega `$
$`\times `$ $`{\displaystyle \frac{1}{1+\left(\beta _2/\alpha \right)^2\left(\omega \omega _1\right)^2\left(\omega \omega _2\right)^2}}`$
which in the appropriate limit is consistent with $`\mathrm{\Delta }C_1`$.
In Fig. 2 we plot the dependence of $`\mathrm{\Delta }C_1`$ on the dimensionless parameter $`\beta _2W^2/\alpha `$. Since momentum change $`\delta \kappa `$ is proportional to $`\beta _2`$, the increase of the dispersion leads to a strong suppression of the power leakage from the bandwidth window, and of the corresponding correction to the channel capacity.
In a communication system with many “fiber-amplifier” units, the fiber nonlinearity leads not only to the mixing of the signals at different frequencies, but also to the mixing of the signal with the noise. Qualitatively, this would correspond to an effective enhancement of the noise power in the system, and therefore to a loss of the channel capacity. This effect is not present, when the system has only one “fiber-amplifier” link, which explains the appearance of the $`(N1)`$ factor in $`\mathrm{\Delta }C_2`$ and $`\mathrm{\Delta }C_3`$.
The effective noise enhancement is caused by the scattering processes, which involve a “signal photon” and a photon, produced due to spontaneous emission in one of the amplifiers. The total power of this extra noise can be expressed as
$`{\displaystyle \frac{\mathrm{\Delta }P_N}{P_0}}`$ $``$ $`\gamma ^2{\displaystyle \frac{P_0P_N}{\alpha ^2}}{\displaystyle _W}𝑑\omega _1{\displaystyle _W}𝑑\omega _2{\displaystyle _W}𝑑\omega _3{\displaystyle _W}𝑑\omega `$ (38)
$`\times `$ $`S[\omega _1,\omega _2\omega _3,\omega ]`$
The corresponding correction to the capacity
$`\mathrm{\Delta }C_N`$ $``$ $`W{\displaystyle \frac{\mathrm{\Delta }P_N}{P_N}}W{\displaystyle \frac{\gamma ^2P_0^2}{\alpha ^2}}{\displaystyle _W}𝑑\omega _1{\displaystyle _W}𝑑\omega _2{\displaystyle _W}𝑑\omega `$ (39)
$`\times `$ $`{\displaystyle \frac{1}{1+\left(\beta _2/\alpha \right)^2\left(\omega \omega _1\right)^2\left(\omega \omega _2\right)^2}}`$
where we assumed $`\alpha d1`$. In this limit (39) is up to a constant factor identical to $`\mathrm{\Delta }C_2`$.
The dependence of $`\mathrm{\Delta }C_2`$ on $`\beta _2W^2/\alpha `$ is also shown in Fig. 2. Note, that $`\mathrm{\Delta }C_2`$ also decreases with the increase of the dispersion, but more slowly than $`\mathrm{\Delta }C_1`$. Since the scattering processes, which contribute to $`\mathrm{\Delta }C_1`$, need to “move” one of the frequencies out of the bandwidth window, they generally involve a substantial change of the total momentum, and are therefore more strongly affected by the dispersion.
The two physical effects, described above, determine the fundamental limit to the bit rate for a fiber optics communication system. As follows from our analysis (see Fig. 2), the relative contributions of $`\mathrm{\Delta }C_1`$ and $`\mathrm{\Delta }C_2`$, often referred to as the “four-wave mixing”, can be suppressed by choosing a fiber with a large dispersion, or when using a larger bandwidth.
In our analysis, we treated the whole available bandwidth as a single channel. As a result, the cross-phase modulation, which severely limits the performance of advanced wavelength-division multiplexing systems (WDM), does not affect the channel capacity. The reason for this seemingly contradictory behaviour, is that in a WDM system, the “receiver”, tuned to a particular WDM channel, has no information on the signals at the other channels. Therefore, even in the absense of the “geniune” noise, the nonlinear interaction between different channels, leading to a change in the signal in any given channel, will be an effective noise source, thus limiting the communication rate. This limit however is not fundamental, and can be overcome by using the whole bandwidth all together.
In conclusion, we developed a perturbative method for the calculation of the channel capacity for fiber optics communication systems. We obtained analytical expressions for the corrections to the Shannon formula due to fiber nonlinearity. We have shown that, compared to the Shannon limit, the actual channel capacity is substantially suppressed by the photon scattering processes, caused by the fiber nonlinearity.
## Appendix A Perturbative Solution of the Propagation Equation
In this Appendix we describe the perturbative solutuion of the nonlinear equation (10) with the boundary condition
$`A(0,t)`$ $`=`$ $`{\displaystyle 𝑑\omega x\left(\omega \right)\mathrm{exp}\left(i\omega t\right)}`$ (40)
We represent $`A(z,t)`$ as a power series
$`A(z,t)`$ $`=`$ $`{\displaystyle 𝑑\omega \mathrm{exp}\left(i\omega t\left[\frac{\alpha }{2}+i\kappa _\omega \right]z\right)\underset{n=0}{\overset{\mathrm{}}{}}\gamma ^{\mathrm{}}}`$ (41)
$`\times `$ $`_{\mathrm{}}(z,\omega )`$
where $`\kappa _\omega `$ is defined is (19). Substituting (41) into Eqns. (10),(40), we obtain:
(i) for $`\mathrm{}=0`$
$`{\displaystyle \frac{_0(z,\omega )}{z}}`$ $`=`$ $`0`$ (42)
$`_0(0,\omega )`$ $`=`$ $`x\left(\omega \right)`$ (43)
(ii) for $`\mathrm{}0`$
$`{\displaystyle \frac{_{\mathrm{}}(z,\omega )}{z}}`$ $`=`$ $`i{\displaystyle \underset{\mathrm{}_1,\mathrm{}_2,\mathrm{}_3=1}{\overset{\mathrm{}1}{}}}\delta _{\mathrm{}1,\mathrm{}_1+\mathrm{}_2+\mathrm{}_3}{\displaystyle 𝑑\omega _1𝑑\omega _2}`$ (44)
$`\times `$ $`_\mathrm{}_1(z,\omega _1)_\mathrm{}_2(z,\omega _2)_\mathrm{}_3^{}(z,\omega _1+\omega _1\omega )`$
$`\times `$ $`\mathrm{exp}\left[i\left(\kappa _{\omega _1}+\kappa _{\omega _1}\kappa _{\omega _1+\omega _2\omega }\kappa _\omega \right)\right]`$
$`_{\mathrm{}}(0,\omega )`$ $`=`$ $`0`$ (45)
where $`\delta `$ in the Kronekker’s delta-function.
For any $`\mathrm{}`$, these equations reduce to linear first order differential equation, and can be solved straightforwardly. For example, the solutions for the first three terms in the anzats (41) are given by:
$`_0(z,\omega )`$ $`=`$ $`x\left(\omega \right)`$ (46)
$`_1(z,\omega )`$ $`=`$ $`{\displaystyle 𝑑\omega _1𝑑\omega _2F_{\omega _1\omega _2}^\omega x\left(\omega _1\right)x\left(\omega _2\right)}`$ (47)
$`\times `$ $`x^{}\left(\omega _1+\omega _2\omega \right)`$
$`_2(z,\omega )`$ $`=`$ $`{\displaystyle }d\omega _1{\displaystyle }d\omega _2{\displaystyle }d\omega _3{\displaystyle }d\overline{\omega }[G_{\omega _1\omega _2\omega _3\overline{m}}^\omega `$ (48)
$`\times `$ $`x^{}\left(\omega _1\right)x^{}\left(\omega _2\right)x\left(\omega _1+\omega _2\overline{\omega }\right)x\left(\omega _3\right)`$
$`\times `$ $`x\left(\omega +\overline{\omega }\omega _3\right)+H_{\omega _1\omega _2\omega _3\overline{\omega }}^\omega x\left(\omega _1\right)x\left(\omega _2\right)x\left(\omega _3\right)`$
$`\times `$ $`x^{}(\omega _1+\omega _2\overline{\omega })x^{}(\overline{\omega }+\omega _3\omega )]`$
Here the functions $`F`$, $`G`$ and $`H`$ are given by
$`F_{\omega _1\omega _2}^\omega `$ $`=`$ $`i{\displaystyle \frac{1\mathrm{exp}\left(\alpha di\varphi _{\omega _1\omega _2}^\omega d\right)}{\alpha +i\varphi _{\omega _1\omega _2}^\omega }}`$
$`G_{\omega _1\omega _2\omega _3\overline{\omega }}^m`$ $`=`$ $`\left(F_{\omega _1\omega _2}^{\overline{\omega }}\right)^{}\left(F_{\omega \overline{\omega }}^{\omega _3}\right)^{}`$
$`\times `$ $`\left(1\mathrm{exp}\left(\alpha d+i\varphi _{\omega \overline{\omega }}^{\omega _3}d\right)\right)`$
$``$ $`\left({\displaystyle \frac{F_{\omega _1\omega _2}^{\overline{\omega }}}{\frac{1}{F_{\omega _1\omega _2}^{\overline{\omega }}}+\frac{1}{F_{\omega \overline{\omega }}^{\omega _3}}}}\right)^{}`$
$`\times `$ $`\left(1\mathrm{exp}\left(2\alpha d+\left(\varphi _{\omega _1\omega _2}^{\overline{\omega }}+\varphi _{\omega \overline{\omega }}^{\omega _3}\right)d\right)\right)`$
$`H_{\omega _1\omega _2\omega _3\overline{\omega }}^\omega `$ $`=`$ $`2\left(F_{\omega _1\omega _2}^{\overline{\omega }}\right)\left(F_{\overline{\omega }\omega _3}^m\right)`$
$`\times `$ $`\left(1\mathrm{exp}\left(\alpha di\varphi _{\overline{\omega }\omega _3}^\omega d\right)\right)`$
$``$ $`2\left({\displaystyle \frac{F_{\omega _1\omega _2}^{\overline{\omega }}}{\frac{1}{F_{\omega _1\omega _2}^{\overline{\omega }}}+\frac{1}{F_{\overline{\omega }\omega _3}^\omega }}}\right)^{}`$
$`\times `$ $`\left(1\mathrm{exp}\left(2\alpha di\left(\varphi _{\overline{\omega _1}\omega _2}^{\overline{\omega }}+\varphi _{\overline{\omega }\omega _3}^\omega \right)d\right)\right)`$
where
$`\varphi _{\omega _1\omega _2}^\omega `$ $`=`$ $`\left(\kappa _{\omega _1}+\kappa _{\omega _2}\kappa _{\omega _1+\omega _2\omega }\kappa _\omega \right)`$
## Appendix B The Perturbative Calculation of the Channel Capacity
In this Appendix, we describe the perturbative calculation of the capacity of the simple nonlinear channel
$`y\left(\omega \right)=x\left(\omega \right)\mathrm{exp}\left(i\gamma \varphi \left[x\left(\omega \right)\right]\right)+n\left(\omega \right)`$ (49)
Here $`\varphi \left[x\right]`$ is an arbitrary real function, and the noise $`n\left(\omega \right)`$ is a Gaussian random variable:
$`p_n\left[n\left(\omega \right)\right]`$ $``$ $`\mathrm{exp}\left[{\displaystyle \frac{\left|n\left(\omega \right)\right|^2}{P_N}}\right]`$ (50)
The fact, that the noise in this model is additive, implies that the conditional distribution $`p\left(y|x\right)`$ is fixed, and defined by the noise distribution $`p_n`$:
$`p\left(y|x\right)`$ $`=`$ $`p_n\left[yx\mathrm{exp}\left(i\gamma \varphi \left(x\right)\right)\right]`$ (51)
It is therefore straightforward to show, that the entropy $`H\left[y|x\right]`$ does not depend on $`\gamma `$:
$`H\left[y|x\right]`$ $`=`$ $`{\displaystyle 𝑑yp_n\left(yxe^{i\gamma \varphi (x)}\right)}`$ (52)
$`\times \mathrm{log}\left[p_n\left(yxe^{i\gamma \varphi (x)}\right)\right]`$
$`=`$ $`{\displaystyle 𝑑zp_n\left(z\right)\mathrm{log}\left[p_n\left(z\right)\right]}`$
$`=`$ $`\mathrm{log}\left[2\pi eP_N\right]`$
In order to calculate the entropy $`H\left[y\right]`$, we represent the output distribution as a power series in $`\gamma `$:
$`p_y\left(y\right)`$ $`=`$ $`p_y^0\left(y\right)\left[1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}p_y^{(n)}\left(y\right)\right]`$ (53)
where $`p_y^0(y)`$ is the “unperturbed”, Gaussian distribution
$`p_y^0\left(y\right)`$ $`=`$ $`{\displaystyle \frac{1}{\pi \left(P_0+P_N\right)}}\mathrm{exp}\left({\displaystyle \frac{\left|y\right|^2}{P_0+P_N}}\right)`$ (54)
corresponding to the linear channel $`y=x+n`$. Substituting (53) into the definition of the entropy $`H_y`$, Eq. (5), we obtain:
$`H_y`$ $`=`$ $`\mathrm{log}\left[2\pi e\left(P_0+P_N\right)\right]{\displaystyle \frac{1}{2}}{\displaystyle p_y^{(1)}\left(y\right)^2p_y^0\left(y\right)}`$ (55)
$`+`$ $`{\displaystyle \frac{1}{P_0+P_N}}\left[{\displaystyle }dy\right|y|^2p_y\left(y\right)`$
$``$ $`{\displaystyle }dy\left|y|^2p_y^0\left(y\right)\right]`$
The second term in Eq. (55), $`(1/2)p_y^{(1)}\left(y\right)^2p_y^0\left(y\right)`$, represents the difference of the output distribution from Gaussian, and corresponds to the contribution $`\mathrm{\Delta }H_y`$ in Eq. (20). Note, that in the second order in nonlinearity the deviations of the output distribution from Gaussian lead to a decrease of capacity.
The third term, $`(1/(P_0+P_N))\left[𝑑y\left|y\right|^2p_y\left(y\right)𝑑y\left|y\right|^2p_y^0\left(y\right)\right]`$, is proportional to the change of the output power, $`𝑑y\left|y\right|^2p_y\left(y\right)`$, due to nonlinearity, and corresponds to $`\mathrm{\Delta }C_1`$ in Eq. (20). Generally, the nonlinearity leads to energy exchange between different degrees of freedom in the channel (e.g. between different frequencies), and to the power leakage out of the bandwidth window. However, for the specific (nad non-generic) example, chosen in the present Appendix, this exchange is absent, since the output power
$`\left|y\right|^2`$ $`=`$ $`\left|x\right|^2+\left|n\right|^2=P_0+P_N`$ (56)
does not depend on the nonlinearity.
Substituting (56) in Eq. (55), and using Eq. (52), for the mutual information $`R`$ we obtain:
$`R`$ $`=`$ $`\mathrm{log}\left[1+{\displaystyle \frac{P_0}{P_N}}\right]{\displaystyle \frac{1}{2}}{\displaystyle p_y^{(1)}\left(y\right)^2p_y^0\left(y\right)}`$ (57)
As immediately follows from Eq. (57), the channel capacity, equal to the maximum of the mutual information, is given by the Shannon formula (1), and is achieved when
$`p_y^{(1)}\left(y\right)`$ $`=`$ $`0`$ (58)
The next step is to calculate the input distribution
$`p_x\left(x\right)`$ $`=`$ $`p_x^0\left(x\right)\left[1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}p_x^{(n)}\left(x\right)\right]`$ (59)
corresponding to (58). The general relation between the input and the output distributions is defined by the conditional distribution $`p(y|x)`$:
$`p_y\left(y\right)={\displaystyle 𝑑xp\left(y|x\right)p_x\left(x\right)}`$ (60)
and, considered as an equation for $`p(x)`$, is a Fredholm intergal equation of the first kind. Note however, that since Eq. (58) represents not the whole output distribution, but only it’s first order term $`p_y^{(1)}(y)`$, we can expand Eq. (60) and keep only the terms up to the first order in $`\gamma `$. We obtain:
$`{\displaystyle 𝑑xp\left(y|x\right)}|_{\gamma =0}p_x^0\left(x\right)p_x^{(1)}\left(x\right)`$ (61)
$`+`$ $`{\displaystyle 𝑑x\frac{}{\gamma }p\left(y|x\right)}|_{\gamma =0}p_x^0\left(x\right)=0,`$
Substituting here the conditional distribution from Eq. (51), we obtain:
$`{\displaystyle 𝑑xp_n\left(yx\right)p_x^0\left(x\right)p_1\left(x\right)}={\displaystyle \frac{i}{P_N}}`$
$`\times {\displaystyle }dxp_n(yx)p_x^0\left(x\right)\varphi \left(x\right)(x^{}yy^{}x),`$ (62)
Using the identity
$`yp_n\left(yx\right)`$ $`=`$ $`\left(x+P_N{\displaystyle \frac{}{x^{}}}\right)p_n\left(yx\right),`$ (63)
and integrating by parts, we can represent the right hand side of (62) as follows:
$`i{\displaystyle 𝑑xp_n\left(yx\right)p_x^0\left(x\right)\frac{i}{P_N}\varphi \left(x\right)\left(x^{}yy^{}x\right)}=`$
$`i{\displaystyle 𝑑xp_n\left(yx\right)p_x^0\left(x\right)\left(x^{}\frac{\varphi \left(x\right)}{x^{}}x\frac{\varphi \left(x\right)}{x}\right)}`$ (64)
Therefore, as follows from Eqns. (62) and (64), the input distribution
$`p_x^{(1)}\left(x\right)`$ $`=`$ $`i\left(x^{}{\displaystyle \frac{\varphi \left(x\right)}{x^{}}}x{\displaystyle \frac{\varphi \left(x\right)}{x}}\right)`$ (65)
This procedure can be followed up for all orders in $`\gamma `$. By a direct calculation, it is straightforward to show, that the channel capacity is represented by the Shannon result (1), which is achieved when for any $`n>1`$
$`\{\begin{array}{cc}p_x^{(n)}\left(x\right)=0\hfill & \\ p_y^{(n)}\left(y\right)=0\hfill & \end{array}`$ (68)
Subsitituting (68) and (65) into Eq. (59), for the input distibution we finally obtain:
$`p_x\left(x\right)`$ $`=`$ $`{\displaystyle \frac{1}{\pi P_0}}\left[1+i\gamma \left(x^{}{\displaystyle \frac{\varphi \left(x\right)}{x^{}}}x{\displaystyle \frac{\varphi \left(x\right)}{x}}\right)\right]`$ (69)
$`\times `$ $`\mathrm{exp}\left[{\displaystyle \frac{\left|x\right|^2}{P_0}}\right]`$
with the corresponding channel capacity
$`C`$ $`=`$ $`W\mathrm{log}\left[1+{\displaystyle \frac{P_0}{P_N}}\right]`$ (70)
This result has a simple physical meaning. When the input distribution is organized in such a way, that the quantity $`z=x\mathrm{exp}\left(i\gamma \varphi \left(x\right)\right)`$ has the Gaussian distribution, then, considering $`z`$ as input, the communication channel becomes linear: $`y=z+n`$, and the channel capacity is therefore given by the Shannon formula (1),(70). The corresponding input distribution is then defined by the Jacobian of the transformation from $`xx_R+ix_I`$ to $`zz_R+iz_I`$, $`(z_R,z_I)/(x_R,x_I)`$ (note, that $`x_R`$, $`x_I`$, $`z_R`$, $`z_I`$ are defined as real variables):
$`p_x\left(x\right)`$ $`=`$ $`{\displaystyle \frac{1}{\pi P_0}}{\displaystyle \frac{(z_R,z_I)}{(x_R,x_I)}}\mathrm{exp}\left[{\displaystyle \frac{\left|x_R^2+x_I^2\right|^2}{P_0}}\right]`$ (71)
which reduces to the distribution (69), since
$`{\displaystyle \frac{(z_R,z_I)}{(x_R,x_I)}}`$ $`=`$ $`1+\gamma x_R{\displaystyle \frac{\varphi (x_R,x_I)}{x_I}}\gamma x_I{\displaystyle \frac{\varphi (x_R,x_I)}{x_R}}`$ (72)
$``$ $`1+i\gamma x^{}{\displaystyle \frac{\varphi (x,x^{})}{x^{}}}i\gamma x{\displaystyle \frac{\varphi (x,x^{})}{x}}`$
This result should be contrasted to the so called “Gaussian estimate” of the channel capacity. In the latter appoach, the information channel is described by the joint Gaussian distribution
$`𝒫(x_\omega ,y_\omega )\mathrm{exp}\left(\left[x_\omega ^{}y_\omega ^{}\right]𝒜\left[\begin{array}{c}x_\omega \\ y_\omega \end{array}\right]\right)`$ (75)
where
$`𝒜`$ $`=`$ $`\left[\begin{array}{cc}x_\omega ^{}x_\omega & x_\omega ^{}y_\omega \\ y_\omega ^{}x_\omega & y_\omega ^{}y_\omega \end{array}\right]^1`$ (78)
The channel capacity is then estimated as the mutual information, corresponding to the distribution (75):
$`C_G={\displaystyle 𝑑\omega \mathrm{log}\left[\frac{x_\omega ^{}x_\omega y_\omega ^{}y_\omega }{x_\omega ^{}x_\omega y_\omega ^{}y_\omega x_\omega ^{}y_\omega y_\omega ^{}x_\omega }\right]}`$ (79)
Under the constraint of the fixed input power $`𝑑\omega \left|x_\omega \right|^2`$, the estimate (79) was shown to give the low bound to the channel capacity.
For the model channel considered in the present Appendix, the “Gaussian estimate” yields an expression, different from the Shannon result. For example, when $`\varphi \left(x\right)=\left|x\right|^2`$, we obtain
$`C_G`$ $`=`$ $`W\mathrm{log}\left[1{\displaystyle \frac{P_0}{\left(P_0+P_N\right)\left(1+\gamma ^2P_0^2\right)^2}}\right]`$ (80)
$`=`$ $`W\mathrm{log}\left[1+{\displaystyle \frac{P_0}{P_N}}\right]2W\gamma ^2P_0^2{\displaystyle \frac{P_0}{P_N}}+𝒪\left(\gamma ^4\right)`$
which, as expected, is smaller that the actual channel capacity (70). Note, that the difference between the exact channel capacity and the Gaussian estimate
$`\delta C`$ $``$ $`CC_G=W\mathrm{log}\left[1+{\displaystyle \frac{P_0}{P_N}}\left(1{\displaystyle \frac{1}{\left(1+\gamma ^2P_0^2\right)^2}}\right)\right]`$ (81)
$`=`$ $`2W\gamma ^2P_0^2{\displaystyle \frac{P_0}{P_N}}+𝒪\left(\gamma ^4\right)`$
is not merely a constant scale factor, but a nontrivial function of the signal to noise ratio, and the nonlinerity.
Even when the input distribution is Gaussian, like e.g. when the phase $`\varphi `$ depends on $`x`$ via the “power” $`\left|x\right|^2`$, the Gaussian Estimate does not yield the exact result. The reason for this behaviour is that the joint Gaussian distribution does not correctly reproduce the conditional distribution $`p\left(y|x\right)`$.
For an essentially noinlinear system (e.g. a fiber optics communication channel), there is generally very little apriori knowledge about the parametric dependence of the Channel Capacity on the signal to noise ratio and other system parameters. In this case, the Gaussian Estimate for the channel capacity can be (should be?) viewed as a very unreliable method, as there is no way to separate it’s artefacts from the actual behaviour of the channel capacity.
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# 1 Introduction
## 1 Introduction
The purpose of this paper is first to show that the type IIB domain wall solution of Bremer et al. provides a supersymmetric realization of the Randall-Sundrum brane-world and secondly to give its ten-dimensional interpretation in terms of IIB 3-branes.
The idea that our universe may be a 3-brane in a higher dimensional spacetime has a history going back nearly two decades . More recently, another viewpoint on this basic idea has grown out of the Hořava-Witten model for M–theory/heterotic string duality, based upon an $`S^1/Z_2`$ orbifold in $`D=11`$ spacetime. This orbifold construction was later realized in a $`D=5`$ compactification by a concrete solution to semiclassical M–theory, i.e. $`D=11`$ supergravity . A key point in this construction was the introduction of flux for the M–theory 4-form field strength $`G_{[4]}`$ wound around the compact dimensions, which were taken to be a Calabi-Yau 3-fold. The resulting reduced theory is a specific version of matter-coupled $`D=5`$, $`N=2`$ supergravity. This dimensionally-reduced theory has a scalar potential arising from the $`G_{[4]}`$ flux, which rules out flat space or indeed any maximally symmetric space as a solution to the equations of motion. But this $`D=5`$ reduced theory readily admits domain wall, i.e. 3-brane, solutions. A natural configuration is a pair of two 3-branes in a $`Z_2`$ symmetric configuration; projecting the fields of this theory into the subspace of $`Z_2`$ invariant configurations then reproduces the Hořava-Witten orbifold. As in the original $`D=11`$/$`D=10`$ discussion, the massless brane-wave excitations of this scenario are not easily deduced by direct analysis of the solution, but one may obtain information about the zero-modes by anomaly inflow arguments. These may either be carried out in $`D=10`$, leading to the original Hořava-Witten prediction of a $`D=10`$, $`N=1`$ super Yang-Mills $`E_8`$ gauge multiplet residing on each of the two fixed planes of the orbifold, with the resulting structure subsequently reduced to $`D=5`$, or one may carry out the anomaly analysis directly in $`D=5`$, yielding more general possibilities for gauge structure .
Another theory in which similar constructions can be made is $`D=10`$ type IIB supergravity. This has a self-dual 5-form field strength that supports the D3-brane, which is the basis for much recent discussion of the Maldacena conjecture, linking string theory in the near-horizon region of the D3-brane to a Yang-Mills theory quantized on the boundary of the associated asymptotic anti de Sitter space, which is the near-horizon limiting spacetime. In the pure supergravity context, relations between $`p`$-branes in higher dimensions and domain walls arising after dimensional reduction on spheres was developed in , including the case of the D3-brane of type IIB theory.
Meanwhile, another development was brewing. Randall and Sundrum proposed a simple model of physics on 3-branes embedded in $`D=5`$ anti de Sitter space, first in a model with two 3-branes, one of positive and one of negative tension. This model was criticized for the apparent danger of non-physical modes from the negative tension brane, and also because the modulus related to the distance between the two branes gave another parameter needing fixing in any phenomenological analysis. Subsequently, a revision of this scenario was put forward , in which there was only one 3-brane, of positive tension, essentially obtained from the first scenario by sending the negative tension brane to the Cauchy horizon of anti de Sitter space. The striking result found in this second scenario is that, although the fifth dimension of space-time is now infinite, the effective gravity theory on the single remaining 3-brane nonetheless has $`D=4`$ and not $`D=5`$ leading behavior. The gravitational potential for static sources starts out with a Newtonian $`1/r`$, corrected by terms of order $`\mathrm{\Lambda }^1/r^3`$, where $`\mathrm{\Lambda }`$ is the $`D=5`$ cosmological constant. This “binding of gravity” to the 3-brane happens when a $`D=5`$ spacetime has a warped product structure, with the warp factor, i.e. the factor multiplying the $`D=4`$ submetric, decreasing as one recedes on either side from the single Randall-Sundrum 3-brane. This corresponds in general terms to the 3-brane acting as a positive-tension source on the right-hand side of the Einstein equations. It was not clear, however, whether this scenario could arise from an explicit solution to a supergravity theory.
Links between the Randall-Sundrum model and supergravity were made in Refs. . In , the $`D=5`$ 3-brane solutions to the type II theory presented in were used to make an analogy to the Randall-Sundrum model. The explicit relation between this construction and the specific Randall-Sundrum model was not fully pinned down, however. This perspective was further elaborated in . Despite the existence of these works, there still seems to be some confusion in the literature as to whether the Randall-Sundrum model can in fact be obtained from an explicit supergravity solution<sup>1</sup><sup>1</sup>1The equivalence of the graviton propagator calculated from closed loops of the $`N=4`$ SCFT in the Maldacena picture and that calculated from tree graphs in the Randall-Sundrum picture was already strongly indicative of a supersymmetric Randall-Sundrum brane-world .. Moreover, there are powerful general arguments as to why smooth supersymmetric solutions obtained from $`D=5`$ gauged supergravity coupled to various combinations of $`D=5`$ matter cannot reproduce a Randall-Sundrum scenario with binding of gravity to the 3-brane. A key word here is ‘smooth.’ Although one might well like to replace the Randall-Sundrum scenario, with its delta-function source, by a smooth solution, experience with domain walls in supergravity (i.e. codimension-one brane solutions) shows them always to be based upon a linear harmonic function in the $`d=1`$ codimension. In order for such a solution to have a localized energy concentration, i.e. a ‘brane,’ some kind of ‘kink’ must be introduced into the linear harmonic function so as to give a location to the domain wall. Thus, the search for a smooth codimension one solution looks rather unlikely to be successful.<sup>2</sup><sup>2</sup>2Some rigorous results along these lines have recently been spelled in Ref. . See also Ref. . Moreover, the remainder of the argument of Refs. concerns the general behavior of renormalization group flows between critical points of coupled supergravity-matter potentials. This gives the impression that even if one were to relax the requirement of smoothness, there would be no solution leading to the binding of gravity to the 3-brane.
In this paper, we shall first explicitly obtain the original (kinked) Randall-Sundrum geometry from type IIB supergravity. This follows from the work of Refs. . This construction makes essential use of the ‘breathing mode’ of the $`S^5`$ dimensional reduction of type IIB supergravity of Ref. . We shall show why this massive mode escapes the constraints on supersymmetric flows by reason of its transforming in a representation with AdS lowest energy $`E_0=8>4`$, thus falling outside the scope of the analysis of Ref. . The breathing-mode solutions, although Kaluza-Klein consistent in a purely bosonic context containing just the breathing mode and gravity, do not really correspond to a pure $`D=5`$ supergravity theory. The construction retains an essential memory of its $`D=10`$ type IIB origin. This is particularly so when one considers the superpartners of the breathing mode, which include massive spin two modes that cannot be retained in a consistent truncation to a finite number of $`D=5`$ fields.
Another memory of $`D=10`$ supergravity in the supersymmetric realization of the Randall-Sundrum geometry resides in the $`Z_2`$ symmetry of this geometry. This geometry is very similar to the $`Z_2`$ symmetric configuration of two M–theory 3-branes in $`D=5`$ that explicitly realizes the Hořava-Witten construction as an M-theory brane solution . In the M–theory solution, the $`Z_2`$ symmetry is central to the appearance of the orbifold, and it also plays a critical rôle in the preservation of unbroken $`D=4`$ supersymmetry on the brane world-volumes . The same is true in the double 3-brane type IIB solution that we present as the supergravity realization of the Randall-Sundrum geometry: continuity of the unbroken supersymmetry Killing spinor depends on the way the $`Z_2`$ symmetry is implemented. In particular, in the M–theory case as well as in the type IIB construction , the constant parameter determining the flux of the relevant underlying form field is $`Z_2`$ odd, and so flips sign upon crossing either of the 3-branes; this flip is crucial for the continuity of the unbroken supersymmetry parameter. Accordingly, in the type IIB case as in the M-theory case, the $`D=5`$ theory is really obtained from a dimensional reduction on a pair of Kaluza-Klein ansätze, one on each side of the $`Z_2`$ symmetric spacetime. Although this construction requires the presence of brane sources for the form-field flux, it is natural in the context of the higher-dimensional M– or type IIB theory. This split ansatz, however, means that it is much less natural to view the geometry as arising in a single $`D=5`$ theory.
Having shown how to obtain the Randall-Sundrum model from type IIB supergravity, we next set out to study the brane-wave oscillations of the solution. This analysis is quite natural in the type IIB analogue of the M–theory $`Z_2`$ symmetric double 3-brane construction . Although, as in , this configuration involves both a positive and a negative tension brane, thus leading to concerns about negative energies, we show that there is a ‘mode-locking’ phenomenon that reduces the zero-modes to just one (positive energy) $`D=4`$, $`N=4`$ Maxwell multiplet in the case of one singly charged brane. This happens because the $`Z_2`$-odd modes turn out to be non-zero modes constrained to be related to Kaluza-Klein massive modes by the Bianchi identities for the type IIB 5-form field strength $`H_{[5]}`$ and for the gravitational curvature. Thus, one does not have to make an explicit projection by hand into a $`Z_2`$-invariant subspace of fields: this projection happens spontaneously, by normal Kaluza-Klein dynamical mechanisms freezing out massive Kaluza-Klein modes. The type IIB models considered here have the great advantage that one can carry out more of the Kaluza-Klein analysis explicitly than in the analogous discussion of M–theory reduced on Calabi-Yau spaces . But it is to be expected that an analogous mode-locking mechanism will operate there as well. And in that case, the mode-locking can be expected to lead to a spontaneous appearance of $`D=4`$ chirality, thus generalizing the appearance of chirality by explicit $`Z_2`$ projection.
## 2 Supersymmetric domain walls and renormalization group flows
While there are many ways of representing a metric on anti-de Sitter space, perhaps the most natural form of the metric from a domain wall point of view is given in terms of Poincaré coordinates
$$ds^2=e^{2gy}\eta _{\mu \nu }dx^\mu dx^\nu +dy^2.$$
(2.1)
Written in this manner, the Minkowski signature boundary of AdS is reached when $`y\mathrm{}`$, while the point $`y\mathrm{}`$ is instead a null surface, the AdS Killing horizon. In the AdS/CFT correspondence, this metric is viewed as the near-horizon geometry of $`N`$ coincident D3-branes, which is described by $`𝒩=4`$ super Yang-Mills living on the boundary. Furthermore, the distance to the boundary is regarded as an energy; from the bulk point of view $`y\mathrm{}`$ is a flow to the UV, while $`y\mathrm{}`$ is a flow to the IR.
The Randall-Sundrum brane-world is obtained by taking two Poincaré patches of AdS, both given by (2.1), and joining them at the brane location $`y=0`$. The resulting Randall-Sundrum metric has the form
$$ds^2=e^{2g|y|}\eta _{\mu \nu }dx^\mu dx^\nu +dy^2,$$
(2.2)
and its geometry gives rise to a localized graviton on the ‘Planck’ brane. Presented as ‘an alternative to compactification’, much has been made of the fact that this scenario binds gravity even though the $`y`$ direction has an infinite extent. Nevertheless, it is apparent from the form of (2.2) that the Planck brane only lives in a tiny portion of AdS, and that movement away from the brane flows towards the Killing horizon and not towards the Minkowski boundary of AdS. Had one instead chosen to join together the $`y<0`$ regions of (2.1), the resulting geometry would preserve the vast majority of the original space, including all of the portion of AdS near the boundary. This then would yield a divergent ‘localization’ volume and give rise to a brane of opposite character to the Randall-Sundrum brane, namely one that does not bind gravity.
In fact, the above observation motivated the authors of Ref. to view the Randall-Sundrum geometry as a warped compactification of F-theory on a Calabi-Yau four-fold. In this picture, the warped geometry arises from the presence of D3-branes situated on the elliptically fibered Calabi-Yau manifold. The five dimensional Randall-Sundrum universe then corresponds to the noncompact four-dimensional spacetime with the addition of a single $`y`$ coordinate which provides a preferred slicing of the internal space along flows between separated stacks of D3-branes. One thus sees that the Randall-Sundrum brane itself is not identified with any one of the D3-branes, but is instead viewed as an effective geometry that arises in interpolating between the near-horizon locations of the D3-branes. In terms of the parametrization in (2.2), the D3-branes are located at the horizons, $`y=\pm \mathrm{}`$, and the apparent infinite extent of the $`y`$ coordinate is simply a result of the warping of the compact space by the D3-branes themselves. The localization of gravity is then explained by the compactness of the underlying F-theory construction. Heterotic and M-theory realizations based on warped Calabi-Yau compactifications have been examined in Ref. .
Returning to a five-dimensional picture, there have been many attempts to explain the Randall-Sundrum scenario from a supersymmetric domain-wall point of view. The advantage of this approach is that one can generally ignore the added complications of the compactification of the underlying IIB theory, and instead focus only on brane constructions in the resulting $`D=5`$ gauged supergravity theory. However, as we emphasize below, it is important to realize that there is no reason (other than simplicity) to expect that the relevant degrees of freedom lie only in the massless supergravity sector. In fact, as emphasized in , massless gauged supergravity precludes the localization of gravity on a brane. Thus massive fields are a necessity.
For the Randall-Sundrum picture to be realistic, where the Planck brane is a dynamical object, it would have to be supported by bulk scalar fields. Thus, in the language of bulk renormalization group flow, we seek a brane solution with stable flows to AdS critical points in the IR on both sides of it. This approach has been studied extensively in both the AdS/CFT and brane-world pictures, with considerable overlap. Nevertheless, the distinction between flows of massless and massive scalars has not always been made clear, so we wish to do so below.
Since we demand that the flow away from the brane is towards an AdS background, the scalars must reach some fixed values corresponding to a critical point in the potential. Then, independent of any specific model, at that point, we may expand the scalars about their fixed values. However before doing so, it is worth realizing that representations in AdS differ from those in a flat background.
Recall that, for AdS<sub>5</sub>, general representations of $`SU(2,2)`$ may be labeled by $`D(E_0,j_1,j_2)`$ where $`E_0`$ is the lowest energy (which may be given in terms of the natural mass scale of the AdS background). For scalars, $`D(E_0,0,0)`$, unitarity requires $`E_01`$ with $`E_0=1`$ corresponding to the singleton representation. General unitarity bounds for $`SU(2,2)`$ as well as for the $`SU(2,2|N/2)`$ superalgebras have been obtained in (see also ). For a scalar field in AdS<sub>5</sub>, the mass is given in terms of $`E_0`$ by $`m^2=E_0(E_04)`$, so that ‘massless’ scalars in fact correspond to $`E_0=4`$. Of course, negative mass squared is not to be feared in an AdS background, provided the Breitenlohner-Freedman bound is satisfied. For this case it corresponds to $`m^24`$, which is saturated for $`E_0=2`$.
To be specific, we now consider the case of a brane supported by a single scalar coupled to gravity, where the Lagrangian takes the form
$$e^1=R\frac{1}{2}\varphi ^2V(\varphi ).$$
(2.3)
While one may generalize by including more scalars, this single scalar example is sufficient to bring out our conclusion. The resulting equations of motion have the form
$`R_{MN}`$ $`=`$ $`\frac{1}{2}_M\varphi _N\varphi +\frac{1}{3}g_{MN}V(\varphi ),`$
$`^2\varphi `$ $`=`$ $`_\varphi V(\varphi ).`$ (2.4)
Note that we do not insist that (2.3) necessarily originates from a supersymmetric theory. However in many cases we are of course interested in supersymmetry. This suggests the identification of a putative ‘superpotential’ $`W(\varphi )`$ with
$$V=(_\varphi W)^2\frac{2}{3}W^2,$$
(2.5)
and the putative ‘transformations’
$`\delta \psi _\mu `$ $`=`$ $`[_\mu \frac{1}{6\sqrt{2}}W\gamma _\mu ]ϵ,`$
$`\delta \lambda `$ $`=`$ $`\frac{1}{2}[\gamma \varphi +\sqrt{2}_\varphi W]ϵ.`$ (2.6)
Identification of the above transformations with those of an actual supergravity theory requires some care.<sup>3</sup><sup>3</sup>3In this paper, we have not performed a full Kaluza-Klein reduction of the type IIB supersymmetry transformations, so all discussions of supersymmetry here are somewhat schematic. They will be sufficient, however, to determine the necessary $`Z_2`$ behavior for having an unbroken $`D=5`$ supersymmetry. In particular, from an $`N=2`$ point of view,<sup>4</sup><sup>4</sup>4We take $`N=2`$ to label minimal supersymmetry in $`D=5`$, corresponding to $`D=4`$, $`N=2`$ supersymmetry. the field $`\varphi `$ may reside in either a vector, tensor or hypermatter multiplet, with possibly different forms of coupling to the fermions. In all cases, the fields $`(g_{\mu \nu },\psi _\mu )`$ and $`(\varphi ,\lambda )`$ would be part of a (not necessarily consistent) truncation of the actual supergravity theory.
As emphasized previously in discussions of holographic renormalization group flows, the equations of motion following from a domain-wall ansatz take on a simple form. Starting with the metric
$$ds^2=e^{2A(y)}\eta _{\mu \nu }dx^\mu dx^\nu +e^{2B(y)}dy^2$$
(2.7)
one obtains the following equations
$`A^2=\frac{1}{24}\varphi ^2\frac{1}{12}e^{2B}V,`$
$`A^{\prime \prime }A^{}B^{}=\frac{1}{6}\varphi ^2,`$
$`\varphi ^{\prime \prime }+(4A^{}B^{})\varphi ^{}=e^{2B}_\varphi V,`$ (2.8)
where primes denote $`y`$ derivatives. The first two equations were obtained by combining components of the Einstein equation. Note that the three equations are not all independent, and we find it convenient to focus only on the first two.
In codimension one, the second metric factor $`e^{2B}`$ is redundant, and may be removed by defining a new coordinate $`\stackrel{~}{y}=e^B𝑑y`$ (keeping in mind that explicit domain wall solutions often have a simpler form when presented in terms of a metric with the $`e^{2B}`$ factor). We proceed by setting $`B=0`$, so the equations resulting from (2) take the form
$$A^{\prime \prime }=\frac{1}{6}\varphi ^2,A^2=\frac{1}{24}\varphi ^2\frac{1}{12}V,\varphi ^{\prime \prime }+4A^{}\varphi ^{}=_\varphi V.$$
(2.9)
As emphasized in , the first of these equations indicates that $`A^{\prime \prime }0`$, with saturation of the inequality corresponding to sitting in the pure AdS vacuum. For the present case, this has the consequence that the function $`A(y)`$ must be concave-down, which is in fact exactly what is needed to support a ‘kink-down’ (i.e. positive tension) Randall-Sundrum brane of the form (2.2) with a continuous metric function.
To study the behavior of the flow to the IR fixed point, we may expand about the fixed value, $`\varphi _{}`$, of the scalar. To quadratic order, the potential then has the form
$$V=12g^2+\frac{1}{2}m^2(\varphi \varphi _{})^2+\mathrm{},$$
(2.10)
where the constant factor is chosen to give the conventional normalization of the AdS curvature,
$$R_{MNPQ}=g^2(g_{MP}g_{NQ}g_{MQ}g_{NP}).$$
(2.11)
While in some cases $`g`$ may coincide with the coupling constant of gauged supergravity, we only take it to parameterize the AdS background at the specific fixed point we are interested in.
We now insert (2.10) into the second equation of (2.9) to find that $`A(y)\pm gy`$, at least up to linear order in $`\varphi `$. Thus we recover the expected linear behavior giving rise to an AdS background. Continuing with the $`\varphi `$ equation of motion, and again working to linear order in $`\varphi `$ (which amounts to making the substitution $`A^{}\pm g`$), we find
$$\varphi ^{\prime \prime }\pm 4g\varphi ^{}m^2\varphi 0,$$
(2.12)
which has in general two solutions,
$$\varphi \varphi _{}+ce^{E_0A(y)},\text{and}\varphi \varphi _{}+ce^{(4E_0)A(y)},$$
(2.13)
where $`E_0=2+\sqrt{(m/g)^2+4}2`$ is given exactly by the mass/$`E_0`$ relation for a scalar field in AdS space. Additionally, for either flow, the metric function behaves like
$$A\pm gy\frac{1}{24}(\varphi \varphi _{})^2.$$
(2.14)
Finally, this allows us to examine the IR flow, corresponding the to behavior in the direction $`A\mathrm{}`$. We see that IR stability is ensured for $`E_0>4`$ by taking the second solution of (2.13), while the flow is always unstable for $`2E_0<4`$, and the massless case, $`E_0=4`$, is marginal.
As a result, the above analysis indicates that $`E_0>4`$ is a necessary condition for IR stability, and hence for the construction of a Randall-Sundrum brane. Note, furthermore, that this result was derived without having to appeal to supersymmetry. Thus it holds in general for both BPS and non-BPS flows. However, as we see below, BPS flows impose a further condition on the relative signs of the terms in the superpotential. This powerful and completely general result was in fact present, although hidden in the discussion of Ref. . However, in , only scalars residing in massless vector multiplets of $`N=2`$ gauged supergravity (i.e. the $`𝒟(2,0,0,0)`$ representation, where the last value denotes the $`U(1)_r`$ charge) were considered. In particular, the authors of relied on the relation $`(_i_jW)_{cr}=\frac{1}{3}g_{ij}W_{cr}`$ (in our normalization) arising from very special geometry. Such scalars always have $`E_0=2`$, yielding the negative reported result. Curiously, while it may not have been appreciated that scalars in the decomposition of an $`N=8`$ gauged supergravity multiplet reside in $`N=2`$ tensor and hypermatter multiplets as well as vector multiplets, such $`N=8`$ scalars all have $`E_0=2`$, 3 or 4 so that they also do not lead to IR stable branes.
Turning now to the case of a supersymmetric flow, it is straightforward to see from (2) that the Killing spinor conditions yield the first order equations
$`A^{}`$ $`=`$ $`\pm {\displaystyle \frac{1}{3\sqrt{2}}}e^BW`$
$`\varphi ^{}`$ $`=`$ $`\sqrt{2}e^B_\varphi W,`$ (2.15)
for a domain wall preserving exactly half of the supersymmetries. This result may in fact also be derived without using the transformations (2), but instead by a traditional BPS argument for finding static minimum energy configurations . Combining both equations gives rise to a holographic renormalization group flow
$$\frac{d\varphi }{dA}=6\frac{_\varphi W}{W},$$
(2.16)
consistent with the second order equations (2). In this case, we expand the superpotential as
$$W=\pm 3\sqrt{2}g(1+\frac{1}{12}\lambda (\varphi \varphi _{})^2+\mathrm{})$$
(2.17)
corresponding to the potential (2.10), provided $`\lambda `$ is identified with either $`E_0`$ or $`4E_0`$. Note that this introduces a two-fold ambiguity. However this is in fact somewhat artificial, since knowledge of the actual supersymmetric theory would fully determine the superpotential (but see e.g. Ref. for a discussion on the relation between $`V`$ and $`W`$ without supersymmetry). In contrast to (2.13), the supersymmetric flow condition, (2.16), gives only a single approach to the fixed point
$$\varphi \varphi _{}+ce^{\lambda A(y)}.$$
(2.18)
As a result, for a BPS flow, not only do we require $`E_0>4`$, but also we learn from the above analysis that $`\lambda =4E_0`$ must be negative in the superpotential, (2.17). The requirement of $`\lambda <0`$ was previously noted in .
This connection between $`E_0`$ and the behavior of a scalar field in AdS was initially made in investigations of the Maldacena conjecture , where $`E_0`$ was related to the conformal dimension of appropriate operators on the CFT side of the AdS/CFT conjecture. In this case, (2.12) taken with exact equality is simply the massive scalar equation in the reference AdS background (2.1). This in itself highlights the similarity between the brane-world scenario and the AdS/CFT conjecture. In some sense, the Randall-Sundrum brane, being inserted at some fixed location in AdS, cuts off the flow to the UV and hence may be described by a Maldacena CFT cut off at some energy scale related to the location of the brane.
## 3 Breathing mode domain walls and the brane-world
Based on the preceding analysis, it is clear that consideration of the massless sector of ($`N=2`$, 4 or 8) gauged supergravities alone does not lead to realistic brane-world configurations. However, for a five dimensional model originating from IIB theory, many other degrees of freedom may come into play. While round $`S^5`$ compactifications of IIB supergravity yield $`N=8`$ gauged supergravity at the massless level, this is also accompanied by a Kaluza-Klein tower of massive states. In general, consistent truncations of sphere reductions are a delicate matter . However it is consistent to include the breathing mode $`\phi `$ in the truncation: although it lives in a massive supermultiplet, it is nevertheless a gauge singlet.<sup>5</sup><sup>5</sup>5Note that we use $`\phi `$ to denote the breathing mode rather than $`\varphi `$, in order to emphasize that it is distinct from the $`D=10`$ dilaton of the type IIB theory. Domain walls supported by the breathing mode have been investigated in Refs. , and have recently been suggested as possible realizations of the brane-world scenario.
To make connection with the Randall-Sundrum model, we examine type IIB string theory compactified on $`S^5`$. This sphere reduction, with the inclusion of a single squashing mode along with the breathing mode, was investigated in . Focusing only on the scalar modes, the resulting five dimensional Lagrangian is
$$e^1_5=R\frac{1}{2}\stackrel{~}{\phi }^2\frac{1}{2}\stackrel{~}{f}^2V(\stackrel{~}{\phi },\stackrel{~}{f}).$$
(3.1)
The scalar potential has the form
$$V(\stackrel{~}{\phi },\stackrel{~}{f})=8m^2e^{\frac{10}{\sqrt{15}}\stackrel{~}{\phi }}+e^{\frac{4}{\sqrt{15}}\stackrel{~}{\phi }}(\mu ^2e^{\frac{6}{\sqrt{10}}\stackrel{~}{f}}R_4e^{\frac{1}{\sqrt{10}}\stackrel{~}{f}}),$$
(3.2)
where the constants $`(m,\mu ,R_4)`$ are parameters of the compactification .
While this potential may now be expanded in the form of Eqn. (2.10), it is perhaps more enlightening to first express it in the form of a ‘superpotential’ according to (2.5). We find
$$W=2\sqrt{2}me^{\frac{5}{\sqrt{15}}\stackrel{~}{\phi }}e^{\frac{2}{\sqrt{15}}\stackrel{~}{\phi }}(\sqrt{2}\mu e^{\frac{3}{\sqrt{10}}\stackrel{~}{f}}+\frac{R_4}{2\sqrt{2}\mu }e^{\frac{2}{\sqrt{10}}\stackrel{~}{f}}).$$
(3.3)
Note that there is a slight sign ambiguity in inverting (2.5); here we have chosen the signs so that $`W`$ has a critical point at
$$e^{\frac{3}{\sqrt{15}}\stackrel{~}{\phi }_{}}=\frac{\mu }{2m}\left(\frac{R_4}{6\mu ^2}\right)^{\frac{3}{5}},e^{\frac{5}{\sqrt{10}}\stackrel{~}{f}_{}}=\frac{R_4}{6\mu ^2},$$
(3.4)
corresponding to that of $`V`$ as well. Expansion of $`W`$ then gives
$$W=3\sqrt{2}m\left(\frac{\mu }{2m}\right)^{\frac{5}{3}}\left(\frac{R_4}{6\mu ^2}\right)[1\frac{1}{3}(\stackrel{~}{\phi }\stackrel{~}{\phi }_{})^2+\frac{1}{2}(\stackrel{~}{f}\stackrel{~}{f}_{})^2+\mathrm{}].$$
(3.5)
Comparison with (2.17) then demonstrates explicitly that the breathing mode $`\stackrel{~}{\phi }`$ has $`E_0=8`$ while the squashing mode $`\stackrel{~}{f}`$ has $`E_0=6`$. Curiously, the two modes enter with opposite signs in $`W`$. While this $`N=8`$ symmetric critical point is indeed a minimum of the potential, it is only a saddle point of $`W`$.
The consequences for the resulting supersymmetric flow were investigated in the previous section. For supersymmetric flows, this critical point is IR stable for the breathing mode, while it is unstable for the squashing mode. This indicates explicitly that simply having a domain wall supported by a scalar with $`E_0>4`$ may be insufficient to ensure the stability of a supersymmetric Randall-Sundrum configuration. Nevertheless, we have now seen why use of the massive breathing mode of sphere reductions has been successful in constructing brane-world domain walls , avoiding the limitations on supersymmetric flows presented in Ref. .
To proceed, we now truncate out the squashing mode by setting $`\stackrel{~}{f}=0`$ and $`R_4=6\mu ^2=\frac{6}{5}R_5`$. After dropping tildes, the resulting potential for the breathing mode is simply
$$V(\phi )=8m^2e^{\frac{10}{\sqrt{15}}\phi }R_5e^{\frac{4}{\sqrt{15}}\phi },$$
(3.6)
and has an AdS minimum at
$$e^{\frac{6}{\sqrt{15}}\phi _{}}=\frac{R_5}{20m^2}.$$
(3.7)
Here $`R_5`$ is the curvature scalar of the round $`S^5`$, arising from the type IIB Kaluza-Klein ansatz
$`ds_{10}^2`$ $`=`$ $`e^{2\alpha \phi }ds_5^2+e^{2\beta \phi }ds^2(S^5)`$
$`H_{[5]}`$ $`=`$ $`4me^{8\alpha \phi }ϵ_{[5]}+4mϵ_{[5]}(S^5),`$ (3.8)
where
$$\alpha =\frac{1}{4}\sqrt{\frac{5}{3}},\beta =\frac{3}{5}\alpha .$$
(3.9)
This also indicates that $`m`$ is essentially the 5-form flux of the Freund-Rubin compactification. Thus the two parameters $`m`$ and $`R_5`$ of the five dimensional potential, (3.6), have their origin in the Kaluza-Klein compactification from ten dimensions. Note that the Kaluza-Klein ansatz (3) for the $`H_{[5]}`$ field strength implies that the Freund-Rubin parameter $`m`$ must be odd under transformations $`yy`$. In order for this to be realized as a symmetry of the type IIB theory, this ‘lower’ $`D=5`$ transformation must be accompanied by an orientation-reversing transformation of $`S^5`$, so that the self-dual structure of $`H_{[5]}`$ is preserved, but with $`mm`$. By the ‘skew-whiffing theorem’ , both orientations have the same (maximal) supersymmetry in the case of $`S^5`$. For any other compactifying 5-manifold the supersymmetries would not match.
For a complete truncation of the sphere compactification down to $`D=5`$, in which all Kaluza-Klein modes except for the breathing mode are discarded, the two parameters $`m`$ and $`R_5`$ satisfy trivial Bianchi identities, and hence must be constant. In this case only a single combination of the two is actually physical. The constant parameter $`R_5`$ may then be viewed as a necessary dimensionful parameter for measuring coordinate distances on the five sphere (much as one would have to introduce a length scale $`L`$ for toroidal compactification, where periodic coordinates are identified as $`y=y+2\pi L`$). The actual invariant (physical) size of the five sphere is then set by the expectation of the breathing mode $`\phi `$. To see formally how $`R_5`$ may be scaled away, consider a shift of $`\phi `$ along with a scaling of $`m`$
$$\phi \phi +\frac{\sqrt{15}}{4}\mathrm{log}\lambda ,mm\lambda ^{\frac{5}{4}}.$$
(3.10)
This transformation has the effect of multiplying $`R_5`$ by $`\lambda `$ in the potential (3.6), so that an appropriate choice of $`\lambda `$ may be used to scale $`R_5`$ to any desired value. A particularly natural choice would be to set $`R_5=20m^2`$, so that the AdS critical point is reached at $`\phi _{}=0`$. From a ten dimensional point of view, the transformation (3.10) results in
$`ds_{10}^2`$ $`=`$ $`\lambda ^{\frac{5}{8}}[e^{2\alpha \phi }ds_5^2+e^{2\beta \phi }\lambda ^1ds^2(S^5)]`$
$`H_{[5]}`$ $`=`$ $`\lambda ^{\frac{5}{2}}[4me^{8\alpha \phi }ϵ_{[5]}+4m\lambda ^{\frac{5}{2}}ϵ_{[5]}(S^5)],`$ (3.11)
which is thus a rescaling of $`S^5`$ combined with a $`D=5`$ ‘trombone’ symmetry.
However, as we will discuss in the following Section, if one no longer truncates out the additional Kaluza-Klein modes, then both $`m`$ and $`R_5`$ no longer need to be taken as constant. In this case, attempts to scale away $`R_5(x)`$ would result in a dynamical scaling by $`\lambda (x)`$. In this sense one simply trades one parameter for another, and cannot fully eliminate $`R_5`$. With this in mind, we maintain both parameters $`m`$ and $`R_5`$ in the solution below.
Breathing mode domain wall solutions follow by making the standard ansatz (2.7) and by solving the resulting equations (2). As mentioned above, keeping two independent factors in the ansatz, $`A(y)`$ and $`B(y)`$, is redundant. For $`B=0`$, the solution was presented in , while it was originally presented in with a different choice of coordinates. The advantage of the original choice is its highlighting of a linear harmonic function as a natural feature of codimension one $`p`$-brane solutions. This solution has the basic form
$`e^{\frac{7}{\sqrt{15}}\phi }=H,e^{4A}=e^B=\stackrel{~}{b}_1H^{\frac{2}{7}}+\stackrel{~}{b}_2H^{\frac{5}{7}},`$
$`H=e^{\frac{7}{\sqrt{15}}\phi _0}+ky,`$ (3.12)
where
$$\stackrel{~}{b}_1=\eta _1\frac{28m}{3|k|},\stackrel{~}{b}_2=\eta _2\frac{14\sqrt{5R_5}}{15|k|}.$$
(3.13)
Here $`\eta _{1,2}=\pm 1`$ are in general independent choices of signs for the solution. For our purposes they are fixed by requiring an appropriate AdS limit for $`\phi \phi _{}`$. This gives $`\eta _2=\eta _1`$ and $`\eta _1`$ chosen so that $`e^{4A}>0`$ in order for the metric to be real at a given initial value of $`y`$.
The linear harmonic function $`H`$ is restricted to be nonnegative. Examination of the solution indicates that the AdS horizon is located at $`H=H_{}e^{7\phi _{}/\sqrt{15}}`$, where $`e^{4A}`$ vanishes. For initial $`H>H_{}`$ the five dimensional space asymptotically flattens out as $`H\mathrm{}`$, with a corresponding limit for the scalar field $`\phi \mathrm{}`$, yielding an asymptotically vanishing scalar potential. This case is the second branch of Ref. , where it was referred to as a hybrid Type II and dilatonic domain wall. On the other hand, for initial $`H<H_{}`$, the solution soon runs into a singularity at $`H=0`$. Note, however, that if one starts with a solution with $`H>H_{}`$ initially and signs $`\eta _{1,2}`$ chosen so as to make $`e^{4A}>0`$ initially, but then follows the evolution of $`H`$ within the spacetime through the $`H=H_{}`$ horizon, the metric in the region with $`H<H_{}`$ becomes complex, so one should really treat the region below the horizon using different, appropriately chosen coordinates. Both the $`H>H_{}`$ and $`H<H_{}`$ cases have a natural interpretation in the lifting of (3) to ten dimensions. In the IIB theory, (3) lifts directly to the geometry of $`N`$ coincident D3-branes with total charge $`\stackrel{~}{k}=m(20/R_5)^{5/2}`$ . The two regions $`H\begin{array}{c}>\hfill \\ <\hfill \end{array}H_{}`$ then correspond to the regions either ‘outside’ or ‘inside’ the D3-brane horizon. This furthermore demonstrates that the first, $`H>H_{}`$, case is nothing but the conventional near horizon limit occurring prominently in the Maldacena conjecture. The second, $`H<H_{}`$, case is unphysical as it stands, however, as it sees a different region of the D3-brane geometry containing a singularity.
Neither case by itself provides a suitable framework for a Randall-Sundrum configuration. While in one direction one may reach an AdS horizon, in the other direction one will either run into a singularity or on out into unbounded flat space. One obvious possibility for obtaining an asymptotically AdS space on both sides of a domain wall is to reflect the solution at $`y=0`$, imposing thus a $`yy`$ $`Z_2`$ symmetry. The resulting two-sided domain wall, supported by an absolute value kink in the linear harmonic function
$$H=e^{\frac{7}{\sqrt{15}}\phi _0}+k|y|,$$
(3.14)
was in fact how the solution was originally presented in . The presence of such a kink is rather natural for a codimension one object. Supergravity $`p`$-brane solutions are generally supported by $`\delta `$-function sources at the locations of the branes themselves, and this remains true for domain walls. Passing through a domain wall, one jumps through a sheet of charge, and this jump in charge manifests itself in a change in the slope of the linear harmonic function. A priori, the slope could take any values on the two sides of the domain wall, but clearly the $`Z_2`$ symmetric jump from $`k`$ to $`k`$ is a natural configuration. We shall see that this configuration is distinguished also by preserving unbroken supersymmetry.
For either the plain unkinked (3) or the kinked (3.14) solution, the slope $`k`$ may be scaled away by taking $`yy/|k|`$ and $`x^\mu x^\mu |k|^{1/4}`$. This explains why the apparent domain wall charge $`k`$ is not directly related to lifted quantities such as the D3-brane charge $`\stackrel{~}{k}`$. However, note that this scaling does not eliminate the sign of $`k`$, leaving thus a distinction between the slope up and slope down possibilities. For discussions of multiple domain wall configurations or brane fluctuations, it is more convenient to retain $`k`$.
If one chooses to restrict the coordinate $`y`$ in (3.14) to range only over the interval $`y_0yy_0`$, identifying the points $`y_0`$ and $`y_0`$, then one obtains a $`Z_2`$ symmetric solution that can serve as the background for a $`Z_2`$ orbifold construction. This orbifold construction is analogous to the treatment of M–theory 3-branes given in as a brane realization of the Hořava-Witten $`S^1/Z_2`$ orbifold, and has also been proposed in the Randall-Sundrum context in . The identification of $`y_0`$ and $`y_0`$ essentially reproduces the original Randall-Sundrum model with both an attractive and a repulsive brane (if one chooses $`k<0`$, then the attractive brane is the one located at $`y=0`$). From the five-dimensional point of view, the $`yy`$ $`Z_2`$ map is a parity flip. As we have mentioned above, however, this alone is not a good symmetry of the underlying type IIB theory. In order for this transformation to be compatible with the round-sphere compactification of the IIB theory, this $`Z_2`$ transformation must combine the flip in $`y`$ with an orientation-reversing transformation of the $`S^5`$. For example, an allowable transformation flips all six of the coordinates transverse to the underlying $`D=10`$ D3-brane. The net effect is to send $`mm`$ as well as $`yy`$.
This orientation reversal has important consequences for the supersymmetry transformations (2), since the superpotential $`W`$ also flips, $`WW`$, under these transformations. Actually, this is what one wants, because if the superpotential were to not to flip in this way, then all supersymmetries would be broken by the domain wall, and it would then no longer be BPS. To see this, consider for example the $`\delta \lambda `$ transformation for the solution (3) with the linear harmonic function (3.14). By truncating out the squashing mode from (3.3), one arrives at the breathing mode superpotential:
$$W=\sqrt{2}m\left[2e^{\frac{5}{\sqrt{15}}\phi }5\sqrt{\frac{R_5}{20m^2}}e^{\frac{2}{\sqrt{15}}\phi }\right].$$
(3.15)
Written as above, this clearly changes sign as $`mm`$. On the other hand, If one were to assume instead that $`W`$ remains invariant, one would find
$`_\phi W`$ $`=`$ $`\frac{2}{3}\sqrt{30}m\left[e^{\frac{5}{\sqrt{15}}\phi }\sqrt{{\displaystyle \frac{R_5}{20m^2}}}e^{\frac{2}{\sqrt{15}}\phi }\right]`$ (3.16)
$`=`$ $`{\displaystyle \frac{\sqrt{30}}{14}}|k|H^1(|\stackrel{~}{b}_1|H^{\frac{2}{7}}|\stackrel{~}{b}_2|H^{\frac{5}{7}})`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|\phi ^{}|e^B,`$
where the signs $`\eta _{1,2}`$ have been chosen to obtain the outside (i.e. $`H>H_{}`$) AdS solution. Inserting this into (2) we would find
$$\delta \lambda =\frac{1}{2}e^B(\gamma ^{\overline{y}}\phi ^{}|\phi ^{}|)ϵ,$$
(3.17)
where $`\overline{y}`$ denotes a local Lorentz index. Because of the absolute value in the linear harmonic function (3.14), $`\phi ^{}`$ changes sign on opposite sides of $`y=0`$. Therefore the assumption of an invariant $`W`$ would leave no possibility of obtaining a Killing spinor that is consistently defined on both sides of $`y=0`$. If one were to attempt to patch together separate Killing spinors on both sides of $`y=0`$, in the case of an invariant $`W`$, the $`y\begin{array}{c}>\hfill \\ <\hfill \end{array}0`$ projections on the supersymmetry parameter would be into mutually orthogonal components, $`(1+\gamma ^{\overline{y}})ϵ_+=0`$ versus $`(1\gamma ^{\overline{y}})ϵ_{}=0`$. However, since the superpotential does change sign under the $`Z_2`$, the absolute value in (3.17) is in fact not present, and we accordingly find global Killing spinors of the form $`ϵ=e^{A/2}(1+\gamma ^{\overline{y}})ϵ_0`$. Similar considerations apply at the location of the second kink in the $`Z_2`$ invariant background. If one expands the theory in modes about this $`Z_2`$ invariant background, keeping only the $`Z_2`$ invariant modes, the resulting theory is equivalent to one defined on an $`S^1/Z_2`$ orbifold.
As we have just demonstrated, the domain wall solution is always one half supersymmetric, with or without the absolute value kink. In particular, the $`Z_2`$ orbifolding has not destroyed any further supersymmetry beyond the original half-BPS solution. On the other hand, there is no restoration of supersymmetry either in the presence of of a kink. Consider taking a simultaneous limit $`k0`$ and $`\phi _0\phi _{}`$. Without the kink, this limit would yield pure AdS, i.e. the D3-brane near horizon limit in which full supersymmetry is restored. But with the kink, one obtains instead a $`Z_2`$ symmetric patching of AdS, with a Randall-Sundrum brane located, say, at $`y=0`$. The presence of the orbifold fixed point prevents the full supersymmetry from being restored. However, this is fully expected when a domain wall is present. Although the $`Z_2`$ symmetrization introduces an absolute value into functions, the Killing spinor equations are of first order, and so do not see any $`\delta `$-function singularities. As long as the conditions (2) are satisfied, the solution remains supersymmetric.
Of course the second order equations of motion will see the $`\delta `$-function brane source. For the solution (3.14), we find that the extra source terms at $`y=0`$ are
$`T_{MN}^{\mathrm{brane}}`$ $`=`$ $`{\displaystyle \frac{3k}{14}}(2\stackrel{~}{b}_1^2e^{\frac{3}{\sqrt{15}}\phi _0}+5\stackrel{~}{b}_2^2e^{\frac{3}{\sqrt{15}}\phi _0}7|\stackrel{~}{b}_1\stackrel{~}{b}_2|)\delta (y)\delta _M^\mu \delta _N^\nu g_{\mu \nu },`$
$`Q^{\mathrm{brane}}`$ $`=`$ $`4\sqrt{{\displaystyle \frac{5}{3}}}{\displaystyle \frac{3k}{14}}(\stackrel{~}{b}_1^2e^{\frac{3}{\sqrt{15}}\phi _0}+\stackrel{~}{b}_2^2e^{\frac{3}{\sqrt{15}}\phi _0}2|\stackrel{~}{b}_1\stackrel{~}{b}_2|)\delta (y),`$ (3.18)
assuming $`\stackrel{~}{b}_1\stackrel{~}{b}_2<0`$ as indicated above. These ‘brane sources’ enter in the equations of motion as
$`R_{MN}\frac{1}{2}g_{MN}R`$ $`=`$ $`T_{MN}^\phi +T_{MN}^{\mathrm{brane}},`$
$`^2\phi `$ $`=`$ $`_\phi V(\phi )+Q^{\mathrm{brane}},`$ (3.19)
where
$$T_{MN}=\frac{1}{2}(_M\phi _N\phi \frac{1}{2}g_{MN}\phi ^2)\frac{1}{2}g_{MN}V(\phi ).$$
(3.20)
Depending on the sign of $`k`$, the branes have either positive or negative energy density. However, in both cases the relation between charge and tension is the same, so the branes may be stacked up in BPS configurations.
We have thus seen that the kinks at the brane locations have different consequences for supersymmetry and for the equations of motion. Since the supersymmetry variations and Killing spinor conditions are of first order, the kinks give rise to possibly discontinuous quantities, but no $`\delta `$-function singularities. On the other hand, the equations of motion will be sensitive to the additional $`\delta `$-function sources. Although one may view the equations of motion as a composition of two supersymmetries, there is no contradiction in the presence and absence of the $`\delta `$-function terms since the Killing spinor equations only give rise to a subset of the full equations of motion. To see this consider again for simplicity the $`\delta \lambda `$ transformation, (2)
$`\delta \lambda `$ $`=`$ $`\frac{1}{2}[\gamma _\varphi +\sqrt{2}_\phi W]ϵ`$ (3.21)
$`=`$ $`\frac{1}{2}e^B[\phi ^{}\gamma ^{\overline{y}}+\sqrt{2}e^B_\phi W]ϵ.`$
Partial breaking of supersymmetry then demands the BPS condition<sup>6</sup><sup>6</sup>6One may of course choose the other sign if so desired. However this is a global choice, and must be consistent in all patches of space. $`\phi ^{}=\sqrt{2}e^B_\phi W`$, relating the scalar to its potential. Similarly, vanishing of the gravitino relates the metric to the scalar potential, $`A^{}=e^BW/3\sqrt{2}`$, as given in (2). Now consider deriving the second order $`A^{\prime \prime }`$ equation of motion by taking a derivative of $`e^BA^{}`$
$$e^B(A^{\prime \prime }A^{}B^{})=\frac{1}{3\sqrt{2}}W^{}$$
(3.22)
For a continuous $`W`$, one simply uses the chain rule, $`W^{}=_\phi W\varphi ^{}`$, and substitutes in the $`\phi ^{}`$ equation to arrive at the $`A^{\prime \prime }`$ equation of motion given in (2). However, the assumption of a continuous $`W`$ is actually too strong. For the $`Z_2`$ invariant case, where $`W`$ changes sign at the brane (say at $`y=0`$), one would also pick up a source term upon differentiating, resulting in
$$A^{\prime \prime }A^{}B^{}=\frac{1}{6}\phi ^2+\frac{\sqrt{2}}{3}e^BW\delta (y).$$
(3.23)
Thus, while supersymmetry implies most of the equations of motion, it does not in fact determine all of them. In fact, for higher codimension branes, there is even more slack between the BPS conditions and the equations of motion. The harmonic function condition, of primary importance in brane constructions, is generally a consequence of the equations of motion, and not supersymmetry .
## 4 D3-branes and the world in ten dimensions
Until now we have focused almost exclusively on the five-dimensional viewpoint of the Randall-Sundrum scenario. Since the breathing mode domain wall has its origins in the $`S^5`$ compactification of IIB theory, it has a natural interpretation in terms of IIB D3-branes . Following this connection from the brane-world geometry to breathing mode branes and then to D3-branes, one is led to a realization of the Randall-Sundrum scenario in terms of IIB theory in an appropriate D3-brane background.
While the lifting of the breathing mode brane to patches of the D3 geometry is straightforward, the resulting configuration has unusual features. Following , lifting of the solution given in (3) proceeds by identifying a ten-dimensional Schwarzschild coordinate
$$\rho =\sqrt{\frac{20}{R_5}}H^{\frac{3}{28}}$$
(4.1)
Using the charge relation $`\stackrel{~}{k}=m(20/R_5)^{5/2}`$ and the Kaluza-Klein ansatz (3), one finds the resulting ten-dimensional metric
$$ds_{10}^2=\stackrel{~}{b}_2^{\frac{1}{2}}\left(1\frac{\stackrel{~}{k}}{\rho ^4}\right)^{\frac{1}{2}}dx_\mu ^2+\left(1\frac{\stackrel{~}{k}}{\rho ^4}\right)^2d\rho ^2+\rho ^2d\mathrm{\Omega }_5^2,$$
(4.2)
which is that of $`N`$ D3-branes of total charge $`\stackrel{~}{k}`$ . A further change of coordinates, $`r^4=\rho ^4\stackrel{~}{k}`$, may be performed to transform this into standard isotropic form
$$ds_{10}^2=\sqrt{\stackrel{~}{b}_2}H_{D3}^{1/2}dx_\mu ^2+H_{D3}^{1/2}(dr^2+r^2d\mathrm{\Omega }_5^2),$$
(4.3)
with a harmonic function $`H_{D3}=1+\stackrel{~}{k}/r^4`$. Note that the constant $`\stackrel{~}{b}_2`$ may easily be scaled out of the longitudinal coordinates.
For the $`Z_2`$ symmetric configuration, obtained by kinking the linear harmonic function, (3.14), we see that $`H`$ is a double valued function of $`y`$. This has the consequence that the lifting relation (4.1) is similarly double valued; opposite sides of the breathing-mode brane lift to identical $`\rho `$ values. While the orbifold picture corresponds to a single slice of the D3-brane geometry, $`\rho [\rho _{},\rho _+]`$, the full circle compactification instead corresponds to two copies of the D3-brane geometry patched together at $`\rho _{}`$ and $`\rho _+`$. Note that the AdS horizon, located at $`H_{}`$, lifts to the D3-brane horizon, located at $`\rho _{}=\stackrel{~}{k}^{1/4}`$. Thus taking the Randall-Sundrum configuration (kink down with $`H>H_{}`$) and pushing the second brane off to the Cauchy horizon corresponds in ten dimensions to taking two copies of the near-horizon geometry of $`N`$ D3-branes, and gluing them together at a value $`\rho _0`$ of the Schwarzschild coordinate corresponding to the initial value $`H_0`$ of the linear harmonic function.
For this Randall-Sundrum configuration, it is instructive to ‘unfold’ the doubled metric, (4.3), by defining a new radial coordinate $`\xi [r_0,r_0]`$ such that $`r=r_0|\xi |`$. After scaling out $`\stackrel{~}{b}_2`$ from (4.3), the lifted Randall-Sundrum metric has the form
$$ds_{10}^2=\left(1+\frac{\stackrel{~}{k}}{(r_0|\xi |)^4}\right)^{\frac{1}{2}}dx_\mu ^2+\left(1+\frac{\stackrel{~}{k}}{(r_0|\xi |)^4}\right)^{\frac{1}{2}}(d\xi ^2+(r_0|\xi |)^2d\mathrm{\Omega }_5^2).$$
(4.4)
The positive tension brane is located at $`\xi =0`$, while the negative tension brane is pushed off to the AdS horizon at $`\xi =\pm r_0`$ (the two values are identified under the $`Z_2`$ orbifolding). As seen explicitly here, this act of patching together two stacks of D3-branes essentially compactifies the six-dimensional space transverse to the branes, and also introduces a curvature discontinuity at $`\xi =0`$, the location of the patching. Furthermore, this compactification introduces a charge conservation condition, implying that the net D3 charge must vanish. Thus the resulting kink at $`\xi =0`$ must include a stack of $`2N`$ negative tension D3-branes, with $`2N`$ units of charge soaking up the $`N+N`$ units of charge from the two stacks of positive tension D3-branes.
The question arises, however, whether placing this stack of $`2N`$ negative tension D3-branes at $`\xi =0`$ is sufficient for generating the kinked Randall-Sundrum geometry. Furthermore, the reduction of D3-brane tension from $`D=10`$ to $`D=5`$ yields the simple result $`T_{D=5}=T_{D=10}`$. In addition to giving rise to the tension discrepancy pointed out in , it also leaves unexplained how positive $`D=5`$ tension arises from negative $`D=10`$ tension. As it turns out, the resolution to both issues is the realization that the $`Z_2`$ orbifolding, or the doubling of spacetime, itself gives rise to a positive tension contribution at $`\xi =0`$, the location of the kink. Of course, it is easy to see that the net tension has to be positive, as that is what is required to ‘fold up’ or compactify the space transverse to the branes. The resulting picture is one of negative tension D3-branes trapped on a positive tension $`Z_2`$ orbifold plane giving rise to a composite description of the Randall-Sundrum configuration .
By starting with a brane-world scenario on a circle, one obviously obtains a compact Kaluza-Klein geometry, corresponding to expanding IIB theory about a $`^{1,3}\times S_1\times S^5`$. The $`S_1`$ coordinate $`y`$ lifts to the radial coordinate $`\rho `$, living in a restricted annular range between the two D3 source shells in a double D3-brane background. Of course there is no surprise in starting with a compact geometry and lifting it to another compact scenario. However, by taking the limit of placing the second brane at the Cauchy horizon of AdS, one effectively decompactifies the original Randall-Sundrum geometry of into the picture of . Nevertheless, from a ten-dimensional point of view, this corresponds to simply extending the range of $`\rho `$ a finite distance so as to reach the doubled D3-brane horizon: the internal space remains compact (at least if the inside-horizon brane cores are disregarded). By smoothing out the patching of the double D3-brane configuration, one presumably obtains a warped compactification with an internal six-manifold in the spirit of .
To complete this D3-brane picture of the brane-world, we present the limit in which the $`Z_2`$ symmetric supergravity solution literally reproduces the Randall-Sundrum configuration of a single positive-tension ‘kink-down’ brane between two patches of anti de Sitter space . Starting from the $`D=5`$ 3-brane metric (3) with $`\stackrel{~}{b}_2>0`$, $`\stackrel{~}{b}_1<0`$, $`k<0`$, we want to take a limit as $`k0_{}`$. However, the inverse power of $`k`$ in $`\stackrel{~}{b}_1`$ and $`\stackrel{~}{b}_2`$ (3.13) makes this appear singular. The cure for this is to take a coordinated limit as $`k0_{}`$ and $`\phi \phi _{}`$. We implement this explicitly by taking
$$e^{\frac{7}{\sqrt{15}}\phi _0}=\left(\frac{20m^2}{R_5}\right)^{\frac{7}{6}}+\beta |k|,\beta >0.$$
(4.5)
Note that for $`\beta >0`$, one has $`e^{\frac{7}{\sqrt{15}}\phi _0}>e^{\frac{7}{\sqrt{15}}\phi _{}}`$, i.e. $`H_0>H_{}`$. Accordingly, for finite $`k<0`$, the harmonic function $`H`$ decreases from its value $`H_0`$, reaching the Cauchy horizon value $`H_{}`$ at $`y=y_h`$. This is the natural point at which to make an identification $`y_hy_h`$, putting the second (negative tension) 3-brane at the horizon. For finite $`k`$, one thus has a “semi-interpolating soliton” in the sense that one of the asymptotic limits of the solution, but not both, corresponds to a vacuum solution of the theory, in this case the AdS space with asymptotic scalar $`\phi _{}`$. At the Randall-Sundrum brane, however, there is no horizon.
Taking the joint limit defined by (4.5) as $`k0_{}`$, the difference between the two harmonic functions in $`e^{2A}`$ partially cancels, giving an expression proportional to $`k`$, which cancels the $`k`$ in the denominators of $`\stackrel{~}{b}_1`$ and $`\stackrel{~}{b}_2`$. The resulting metric function is then given by
$`e^{4A}`$ $`=`$ $`4m\left({\displaystyle \frac{R_5}{20m^2}}\right)^{\frac{5}{6}}(\beta |y|)`$ (4.6)
$`=`$ $`{\displaystyle \frac{4}{L}}(\beta |y|),`$
where $`L=m^1(20m^2/R_5)^{5/6}`$ and the $`y`$ coordinate remains restricted to a compact range, $`|y|<\beta `$. This corresponds to the line element
$$ds^2=\frac{2}{\sqrt{L}}(\beta |y|)^{\frac{1}{2}}\eta _{\mu \nu }dx^\mu dx^\nu +\frac{L^2}{16}\frac{dy^2}{(\beta |y|)^2}.$$
(4.7)
The apparent infinite range of the fifth dimension<sup>7</sup><sup>7</sup>7As always in anti de Sitter space, proper distance to the Cauchy horizon is infinite, but the affine parameter along a null geodesic to the horizon is finite. is obtained by making a change of variables
$`\beta |y|`$ $`=`$ $`\beta e^{4|\stackrel{~}{y}|/L},`$
$`x^\mu `$ $`=`$ $`\left({\displaystyle \frac{L}{4\beta }}\right)^{\frac{1}{4}}\stackrel{~}{x}^\mu ,`$ (4.8)
resulting in the five-dimensional metric
$$ds^2=e^{2|\stackrel{~}{y}|/L}\eta _{\mu \nu }d\stackrel{~}{x}^\mu d\stackrel{~}{x}^\nu +d\stackrel{~}{y}^2,$$
(4.9)
which is literally the Randall-Sundrum solution . This sign of the kink ($`k<0`$) thus corresponds to a binding of gravity to the 3-brane at $`y=0`$, with a metric corresponding to segments of pure anti de Sitter space everywhere off this brane surface.
In taking the above Randall-Sundrum limit $`k0_{}`$, $`\phi _0\phi _{}`$, the ten-dimensional coordinate $`\rho `$ is restricted to a progressively limited range near $`\stackrel{~}{k}`$, or, equivalently, $`r`$ is progressively restricted to a range near $`r=0`$. Thus, from a $`D=10`$ perspective, the ‘infinite’ Randall-Sundrum scenario corresponds to shrinking the outer (RS) brane source tightly around the inner horizon brane. Clearly, what is infinite and what is infinitesimal in this subject is frame-dependent.
It is instructive to see in addition the scaling of the ‘brane sources’, (3), in the Randall-Sundrum limit. Taking $`k0_{}`$, we find
$`T_{MN}^{\mathrm{brane}}`$ $`=`$ $`{\displaystyle \frac{24}{L^2}}\delta (y/\beta )\delta _M^\mu \delta _N^\nu g_{\mu \nu }`$ (4.10)
$`=`$ $`2V_{}\delta (y/\beta )\delta _M^\mu \delta _N^\nu g_{\mu \nu },`$
while $`Q^{\mathrm{brane}}=0`$. This vanishing of the scalar charge is in fact forced on us since $`\phi `$ decouples from the solution in this limit. This brings up a key observation that it is not so much the breathing mode $`\phi `$ that supports the brane, but rather $`H_{[5]}`$ flux corresponding to D3 charge. In addition, it is also the behavior of $`H_{[5]}`$ flux that saves the BPS condition with $`Q^{\mathrm{brane}}=0`$; the variation $`\delta \lambda `$ becomes trivial (as it must for a decoupling scalar), while the gravitino transformation becomes that of pure AdS but with a sign flip $`W_{}W_{}`$ at $`y=0`$ (corresponding to a Freund-Rubin compactification with opposite $`S^5`$ orientations). This preservation of supersymmetry further supports the D3-brane origin of the Randall-Sundrum brane-world, via the double 3-brane configuration that we have presented.
The above successful reproduction of the Randall-Sundrum scenario with a ‘kink-down’ (i.e. positive tension) domain wall embedded into $`D=5`$ anti de Sitter space depends crucially upon use of the breathing mode $`\phi `$, which we have shown to transform in a necessary $`E_0>4`$ anti de Sitter representation. Noted as a possibility for a Randall-Sundrum scenario in , this mode escapes the analysis of because it belongs to a massive spin-two multiplet, and thus does not belong to an intrinsically $`D=5`$ supergravity theory. This is because the full multiplet of the breathing mode’s superpartners cannot be retained in a ‘consistent’ Kaluza-Klein reduction, since it involves a massive spin two mode, which never can be kept in a consistent reduction on spheres . With respect to the $`D=5`$, $`N=8`$ supersymmetry, the breathing mode belongs to a multiplet containing 20 copies of the following sets of fields: 1 spin 2, 4 spin 3/2, 26 spin 1, 20 spin 1/2, 15 spin 0. With respect to a $`D=5`$, $`N=2`$ decomposition, it belongs to a long massive vector supermultiplet<sup>8</sup><sup>8</sup>8The $`N=8N=2`$ decomposition of this bulk multiplet in terms of superconformal $`N=4`$, $`D=4`$ boundary supermultiplets was given in . which is another way of explaining why it escaped the analysis of . Since the breathing mode is an $`SO(6)`$ singlet, only the inclusion of the breathing mode’s non-singlet superpartners leads to difficulties with Kaluza-Klein consistency; truncation to the purely bosonic theory involving just $`D=5`$ gravity and the breathing mode is fully consistent.
## 5 Mode locking and spontaneous reduction to an orbifold
The $`Z_2`$ symmetric scenario presented above, with two branes of opposite tension and opposite magnetic charge, corresponding to (3.14), is clearly similar to the brane constructions of Hořava-Witten orbifolds in M–theory given in . The analogous type IIB situation has the great advantage that one can work out explicitly many features of the dynamics, whereas the analogous discussions in M–theory reduced on Calabi-Yau 3-folds must necessarily remain rather implicit. Here, we wish to explore further the properties of this $`Z_2`$ symmetric solution, and see to which extent it naturally corresponds to an orbifold compactification.
The orbifold compactification may be viewed as a compactification on a circle with an additional projection of all the fluctuations into $`Z_2`$ even states only. In a Kaluza-Klein spirit, however, one can investigate the possibility of removing the enforced $`Z_2`$ projection, in order to see what the theory does purely of its own accord when compactified about the double 3-brane background. Thus, we start without making any $`Z_2`$ projections, but still shall take the $`y`$ direction to be a circle. As explained above, from a ten-dimensional point of view, the D3-branes now have no noncompact transverse directions. Thus there is an added cohomology constraint, which demands that there cannot be any non-zero net magnetic charge in the compact transverse space. Unlike general warped compactifications, which allow for additional fields and non-trivial topology, we shall maintain our focus on the round $`S^5`$ and the breathing mode of the compactification. Then, the simplest allowed configuration on the circle is to have a simple pair of 3-branes with opposite magnetic charges. Placing the branes at opposite points on the circle gives rise to a $`Z_2`$ symmetric configuration. However, without imposing the $`Z_2`$ orbifold symmetry, it would appear that the branes are free to move independently. But we shall now demonstrate that this is not the case; instead, there is a mode-locking phenomenon that links the fluctuations of the two 3-branes into a $`Z_2`$-invariant combination.
Consider the $`y`$ coordinate to be periodic with length $`2\mathrm{}`$, making the identification at $`y=\rho _1\rho _2`$. For bosonic fields on this circle, one must impose continuity conditions at both the locations of the 3-branes. Demanding continuity of the scalar field $`\phi `$ and the metric component $`e^{2A}`$ at $`y=0`$ and also at $`y=\rho _1y=\rho _2`$, one has four continuity conditions to satisfy. In this discussion we shall take the overall periodicity length $`2\mathrm{}`$ to be fixed, so $`\rho _1+\rho _2=2\mathrm{}`$. ¿From continuity of the scalar field $`\phi `$, one simply obtains at $`y=0`$ that the value $`\phi _0`$ must be a common limit of $`\phi `$ as one approaches the $`y=0`$ RS brane either from the left or from the right. Continuity at $`y=\rho _1\rho _2`$ implies continuity of the harmonic function $`H`$, so one obtains $`|k_1|\rho _2=|k_2|\rho _2`$, or, using $`\rho _1+\rho _2=2\mathrm{}`$, that $`\left|\frac{k_1}{k_2}\right|=\frac{2\mathrm{}}{\rho _1}1`$. Imposing as well the periodicity conditions on the metric function $`e^{2A}`$ at $`y=0`$ and $`y=\rho _1\rho _2`$, one obtains the continuity conditions
$$\left|\frac{k_2}{k_1}\right|=\frac{|m_2|\sqrt{\frac{R_{5(2)}}{20}}e^{\frac{3}{\sqrt{15}}\phi _0}}{|m_1|\sqrt{\frac{R_{5(1)}}{20}}e^{\frac{3}{\sqrt{15}}\phi _0}}=\frac{|m_2|\sqrt{\frac{R_{5(2)}}{20}}(e^{\frac{3}{\sqrt{15}}\phi _0}+|k_2|\rho _2)}{|m_1|\sqrt{\frac{R_{5(1)}}{20}}(e^{\frac{3}{\sqrt{15}}\phi _0}+|k_1|\rho _1)}.$$
(5.1)
These conditions are solved by matching relations for $`m`$ and $`R_5`$ between the two regions:
$$m_2=(\frac{2\mathrm{}}{\rho _1}1)^1m_1,\sqrt{R_{5(2)}}=(\frac{2\mathrm{}}{\rho _1}1)^1\sqrt{R_{5(1)}}.$$
(5.2)
Accordingly, if one now makes a standard soliton-physics ansatz by letting the $`Z_2`$-odd modulus $`\rho _1`$ become dependent upon the $`D=4`$ coordinates $`x^\mu `$, then upon substitution back into the field equations, one obtains the effective equation for $`\rho _1(x^\mu )`$. Because the oscillations of this coordinate are linked by (5.2) to the Kaluza-Klein ansatz parameters $`m`$ and $`R_5`$, however, this specific modulus has special restrictions on its oscillations. Both $`m`$ and $`R_5`$ are curvature components, and are thus subject to Bianchi identities. To see this for $`m`$, consider the Kaluza-Klein ansatz (3), together with the Bianchi identity
$$dH_{[5]}+\frac{1}{2}ϵ_{ij}F_{[3]}^iF_{[3]}^j0.$$
(5.3)
Letting $`mm(x)`$ and substituting the original ansatz (3), one obtains directly a suppression of $`m`$ fluctuations, $`_\mu m(x)=0`$. For this reason, parameters entering into generalized Kaluza-Klein ansätze like (3) have been sometimes been called “non-zero modes” . In order to see the dynamics of such modes in more detail, one should restore the massive Kaluza-Klein modes that are normally set to zero in a compactification. In the case of $`m`$, this means replacing the ansatz (3) by
$$H_{[5]}=4m(x)e^{8\alpha \phi }ϵ_{[5]}+4m(x)ϵ_{[5]}(S^5)+h_{[5]},$$
(5.4)
where $`h_{[5]}`$ represents the fluctuating massive Kaluza-Klein modes. Re-performing the analysis of the Bianchi identity (5.3) for this generalized ansatz, one now shows that a non-vanishing $`_\mu m`$ must be proportional to $`ϵ^{z_1z_2z_3z_4z_5}_{[z_1}h_{|\mu |z_2z_3z_4z_5]}`$, where $`h_{\mu z_2z_3z_4z_5}`$ is a Kaluza-Klein massive mode, with mass determined as usual by the inverse radius of the $`S^5`$ internal sphere, i.e. corresponding to the length scale of the $`D=5`$ anti de Sitter space. Thus, $`m(x)`$, and hence $`\rho _1(x)`$ are in fact Kaluza-Klein massive modes, and become ‘frozen out’ at energies lower than the AdS scale. Similar considerations apply to the non-zero mode $`R_5`$, which is the Ricci scalar of the internal $`S^5`$ sphere, upon use of the gravitational curvature Bianchi identity. Specifically, in the simple case with Kaluza-Klein massive modes set to zero, if one sets to zero the $`D=5`$ Bianchi identity $`^M(R_{MN}\frac{1}{2}g_{MN}R)=0`$ and uses the dimensionally reduced field equations, one finds, for $`R_5R_5(x^\mu )`$, the constraint $`_\nu R_5\mathrm{exp}(\frac{1}{2}\sqrt{\frac{5}{3}}\phi )m_\nu m=0`$, thus locking out the low energy $`R_5(x^\mu )`$ fluctuations as well.
Given that the $`Z_2`$ odd modes are linked via Bianchi identities to massive Kaluza-Klein modes, one expects the theory to settle down into a low-energy effective theory that is $`Z_2`$ symmetric. Strictly speaking, all that has been demonstrated above so far is that the $`D=4`$ derivatives $`_\mu m`$, $`_\mu R_5`$ are locked out at low energies. In order to show that the theory settles down into a $`Z_2`$ symmetric lowest energy configuration, one would need either to analyze in detail the energy functional for the compactified theory, or to study in more detail the equations of motion of the massive modes. It is likely that the analysis of $`Z_2`$ odd modes can only be done fully consistently if one keeps the entire Kaluza-Klein towers of massive states.
However, one can get an idea of the situation that is obtained with non-$`Z_2`$-symmetric configurations if one considers in a little more detail the question of supersymmetry preservation in a patched background with the matching conditions (5.1, 5.2). Locally, in a patch, there is no difficulty in finding a Killing spinor. However, once one declares that the overall compact part of the spacetime is $`S^5\times S^1`$, one is required to impose continuity and periodicity conditions both for bosons and for fermions.
In the $`Z_2`$ symmetric configuration of the two 3-branes, we have already demonstrated while discussing the unbroken supersymmetry transformation of Section 3,
$$\delta \lambda =\frac{1}{2}e^B(\gamma ^{\overline{y}}\phi ^{}\phi ^{})ϵ,$$
(5.5)
that there is a consistently defined and continuous unbroken supersymmetry transformation with a $`Z_2`$ even global Killing spinor $`ϵ=e^{A/2}(1+\gamma ^{\overline{y}})ϵ_0`$. Now consider the form of the broken supersymmetry transformations in the double 3-brane background. As one can see from the supersymmetry algebra the anticommutator $`\{Q_{\mathrm{broken}},Q_{\mathrm{preserved}}\}`$ involves a translation in the fifth coordinate $`y`$, which is clearly $`Z_2`$ odd. Indeed, the broken supersymmetry parameters will have $`Z_2`$ odd projection conditions. This $`Z_2`$ odd character is canceled, however, in expressions for Goldstone spinor zero modes like (5.5), by the $`Z_2`$ odd character of $`\phi ^{}`$. Combining the Goldstino expression for $`y>0`$ with the $`Z_2`$ map for $`y<0`$ amounts to inserting an absolute value sign around $`\phi ^{}`$ in (5.5), taking a broken supersymmetry parameter for $`ϵ`$. Thus, overall, the Goldstino zero mode is $`Z_2`$ even, as it must be in a consistent truncation.<sup>9</sup><sup>9</sup>9The ‘kink’ in the Goldstino expression resulting from (5.5) with the replacement $`\phi ^{}|\phi ^{}|`$ corresponds to the sign flip of the superpotential $`W`$. That $`W`$ flips without necessarily passing through zero is what allows the Goldstino mode to be normalizable in the present case, thus circumventing the normalizability problems for Goldstinos described in Ref. . Overall, the zero modes of the double 3-brane geometry form a single $`D=4`$, $`N=4`$ super Maxwell multiplet.
Now consider what happens if one tries to expand around a non-$`Z_2`$-symmetric configuration of 3-branes. For the Killing spinor itself, one may observe that $`ϵ=e^{A/2}(1+\gamma ^{\overline{y}})ϵ_0`$ is in fact still continuous and well-behaved in the non-symmetric case, since the metric function $`e^{2A}`$ is by construction matched at the branes. However the situation is different for the candidate Goldstinos. For a non-$`Z_2`$-symmetric configuration the derivative $`\phi ^{}`$ differs by more than a sign as one crosses a 3-brane: in this case one has $`|k_2||k_1|`$, so there is a non-unimodular factor present as well. This prevents one from having continuity both of the unbroken supersymmetry parameter and of the Goldstinos. Thus, although things look locally like one has a BPS configuration with unbroken supersymmetry for a non-$`Z_2`$-symmetric configuration, analysis of the putative zero-mode supermultiplets finds them to be inconsistent with the available matching conditions. So, one is lead to conclude that only the $`Z_2`$ symmetric configuration has a proper unbroken supersymmetry and zero-mode multiplets transforming correctly with respect to it.
The configuration with globally unbroken supersymmetry should be the proper ‘vacuum’ in this double 3-brane sector of type IIB theory compactified on $`S^5`$. A fuller analysis of this spontaneous reduction to a $`Z_2`$ invariant effective theory on the basis of energy functionals and the equations of motion for the Kaluza-Klein massive modes would be desirable. But it is already clear that this double 3-brane model displays a remarkable spontaneous appearance of an orbifold structure. This happens not by insistent projection into a $`Z_2`$ invariant sector of the theory, but naturally by virtue of the Kaluza-Klein dynamics of the theory.
Our discussion has indicated that the original Randall-Sundrum model arises naturally when the fifth dimension $`y`$ direction is taken to be compact, and one may view the model as a system of two D3-branes transverse to the internal $`S_1\times S^5`$. From the $`D=5`$ point of view, there are two branes, one with positive and one with negative tension, constrained by Kaluza-Klein dynamics to live at diametrically opposed points on the circle. While the presence of a negative tension brane might appear troublesome, we have shown that it does not contribute the naïvely anticipated negative energy modes; these are non-zero modes and mix with higher Kaluza-Klein massive modes. The negative tension 3-brane has the effect of protecting the spacetime from curvature singularities in the geometry that might reside behind the Cauchy horizon. Of the a priori two independent types of motion of the 3-branes along the $`S^1`$ direction, only the $`Z_2`$ even modes, corresponding to an overall ‘rotation’ of both branes along the circle, localized in the $`D=4`$ coordinates $`x^\mu `$, correspond to genuine zero modes.
## 6 Conclusions
We have found that an appropriately constructed D3-brane configuration provides a supersymmetric and dynamically stable Randall-Sundrum scenario. This is achieved in a solution to the $`D=10`$ type IIB supergravity equations which can be given a $`D=5`$ interpretation, but is not fully a $`D=5`$ solution, for it employs an intrinsically massive Kaluza-Klein mode, the $`S^5`$ breathing mode. This mode has AdS energy $`E_0=8`$, satisfying the bound $`E_0>4`$ that is required for an asymptotic approach to AdS space from a downwards-facing warp-factor kink in a Randall-Sundrum scenario. There is also a $`Z_2`$ flip in the sign of the Freund-Rubin parameter $`m`$. This is natural enough in a $`D=10`$ context where $`m`$ is a field-strength value, but it is less natural from a $`D=5`$ viewpoint, where $`m`$ normally would appear as a parameter. We have found, moreover, that although one can decide to exclude the $`Z_2`$ odd modes when expanding the theory around the presented $`Z_2`$ invariant background, and thus reproduce an $`S^1/Z_2`$ orbifold reduction, it is not actually necessary to make this projection by hand. Bianchi identities for the curvature values entering in the solution relate the $`Z_2`$ odd modes to Kaluza-Klein massive states of the theory, and so they decouple naturally at low energy. Although charge conservation on the circle requires branes to come in oppositely charged pairs, we have seen that one can recover a single brane Randall-Sundrum model by pushing the second brane off to the Cauchy horizon (i.e. by taking $`\phi _1=\phi _{}`$ for the second brane). From the $`D=10`$ point of view, however, this corresponds to shrinking an outer RS shell of D3 brane tightly around an inner ‘horizon’ D3 brane of opposite charge and tension. Clearly, an important problem is whether this geometry can be realized in a string theory context.
## Acknowledgments
We would like to acknowledge stimulating discussions with Shanta de Alwis, Mirjam Cvetic, Marc Henneaux, Renata Kallosh, Neil Lambert, Finn Larsen, Ruben Minasian, Herve Partouche, Chris Pope, Lisa Randall, Hisham Sati and Herman Verlinde.
## Note added
As this paper was in the final stages of preparation, a very interesting paper appeared that sheds light on the relationship between constructions such as those of Refs. or the present paper and the supersymmetry scheme of Ref. , which was otherwise puzzling. Ref. discusses supersymmetry in orbifolds, and in particular the $`D=5`$ case of interest here. In order to obtain a continuous Killing spinor at orbifold singularities (necessary for the Killing equation to be realized everywhere, including at the singular points), Ref. introduces a 5-form ‘theory of nothing’ field strength, which has just a constant as a solution, but allows for this variable to be only piecewise constant. This allows for a $`Z_2`$ sign flip in the prepotential that is critical for having a preserved supersymmetry allowing coupling to supermatter. This sign flip was not made in the discussion of Ref. , leading to problems with matter coupling. This difficulty of , and the resolution of was also investigated in and independently worked out by . We anticipate that a fuller Kaluza-Klein treatment of the type IIB theory, including all fermions and making a careful reduction of the type IIB supersymmetry transformations, will show that the $`D=5`$ supersymmetry realization adopted in Ref. can also be viewed as the natural dimensional reduction of the type IIB theory using a $`Z_2`$ symmetrized ansatz of the type employed in Refs. and the present paper. In particular, we expect that the 5-form ‘theory of nothing’ field introduced in Refs. can be identified with the $`D=5`$ residue of the type IIB self-dual 5-form field strength.
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# Precise Physics of Simple Atoms
## Low-energy tests of QED
The precise physics of simple atom is the most interesting part of the so-called low-energy tests of Quantum Electrodynamics (QED). Low energy tests of QED offer a number of different options:
* A study with free particles provides the possibility of testing the QED Lagrangian for free particles. The most accurate data arise from anomalous magnetic moments of the electron (Kinoshita<sup>1</sup><sup>1</sup>1References marked with $``$ correspond to presentations at a satellite meeting to ICAP named Hydrogen atom, 2: Precision physics of simple atomic system and its Proceedings will be published by Springer in 2001.) and the muon Hughes .
* However, one knows that the bound problem makes all calculations more complicated. Bound state QED is not a well–established theory. It involves different effective approaches to solve a two–body problem. These approaches can be essentially checked with low $`Z`$ atomic systems like for e. g. hydrogen and deuterium (Hänsch), neutral helium (Drake) and helium ions, muonium (Jungmann), positronium (Conti), muonic hydrogen and helium etc. We consider here most of the low–$`Z`$ atoms.
* Study of high–$`Z`$ ions (Myers, Stölker) cannot further the test of the bound state QED because of large contributions due to the nuclear structure. Rather such an investigation is useful for trying different nuclear models. However, in some particular cases, atomic systems with a not too high $`Z`$ can give some important information on higher order terms of the QED $`Z\alpha `$ expansion.
* There are some other two–body atoms under investigation. They contain a hadron as an orbiting particle. Different antiprotonic (Yamazaki) and pionic (Nemenov) atoms provide a unique opportunity to study particle property with spectroscopic means with a high precision. In some sense it is not possible to have low precision: if a signal is detected the accuracy is granted.
The precision study of the simple atoms is not only limited by experiments with simple atoms. The theory is not able to predict anything to be comparable to the experimental data. What theory can do is to express a measurable quantity in terms of fundamental constants and particle (or nuclear) properties.
First of all we need to determine somehow the Rydberg constant ($`R_{\mathrm{}}`$), the fine structure constant ($`\alpha `$) and the electron mass in some appropriate units (e. g. in atomic units or in terms of the proton mass ($`m_e/m_p`$)). Uncertainties arising from these constants are sometimes compatible with other items of the uncertainty budget or they are even sometimes the most important source of inaccuracy. One should remember that the electron is the most fundamental particle for physics, chemistry, and metrology and the constants associated with its properties go through any atomic spectroscopic effects and any quantum electromagnetic effects. Due to that a number of different studies, which are very far from the spectroscopy of simple atoms (like e. g. Watt balance experiment (see Mohr for detail)), are really strongly connected with the precision physics of simple atoms.
However, a knowledge of the universal fundamental constants is not enough for precision theoretical predictions and we need to learn also some more specific constants like for e. g. the muon mass or the proton electric charge radius. The former is important for the muonium hyperfine structure, while the latter is for calculating the hydrogen Lamb shift.
## Spectra of simple atoms
Let us discuss the spectrum of simple atoms in more detail. The gross structure of atomic levels in a hydrogen–like atom comes from the Schödinger equation with the Coulomb potential and the result is well–known<sup>2</sup><sup>2</sup>2We use the relativistic units in which $`\mathrm{}=c=1`$. $`E_n=(Z\alpha )^2m_e`$. There are a number of different corrections: the relativistic ones (one can find them from the Dirac equation), the hyperfine structure (due to the nuclear magnetic moment) and the QED ones. A structure of levels with the same value of the principal quantum number $`n`$ is a signature of any atomic system. In Fig. 1 we present three different spectra of the structure at $`n=2`$. The first one is realized in “normal” (electronic) hydrogen–like atoms (hydrogen, deuterium, helium ions etc). The muonium spectrum is the same. The largest splitting, of order $`(Z\alpha )^4m_e`$, is the fine structure (i. e. the splitting between levels with a different value of the electron angular momentum $`j`$), the Lamb shift arising from the electron self–energy effects is of order $`\alpha (Z\alpha )^4m_e\mathrm{ln}\left(1/(Z\alpha )\right)`$ and it splits the levels with the same $`j`$ and different values of the electron orbital momentum $`l`$. Some nuclei are spinless (like e. g. in <sup>4</sup>He), while others have a non–zero spin (in hydrogen, deuterium, muonium, helium–3). In the latter case, the nuclear spin splits levels with the same electronic quantum number. The splitting are of order $`(Z\alpha )^4m_e^2/M`$ or $`\alpha (Z\alpha )^3m_e^2/m_p`$, where $`M`$ is the nuclear mass, and the structure depends on the value of the nuclear spin. The scheme in Fig. 1 is for nuclear spin 1/2 (hydrogen).
The structure of levels in positronium and muonic atoms is different because some other effects enter into consideration. For positronium, an important feature is a real (into two and three photons) and virtual (into one photon) annihilation. The former is responsible for the decay of the s-states, while the latter shifts triplet levels (and 2<sup>3</sup>s<sub>1</sub> in particular). The shift is of the order of $`\alpha ^4m_e`$. Contributions of the same order arise from relativistic effects and hyperfine interactions. As a result the positronium level structure at $`n=2`$ has no hierarchy (Fig. 1).
Another situation is that for the muonic atoms. A difference comes from a contribution due to the vacuum polarization effect (the Uehling potential). Effects of electronic vacuum polarization shift all levels to the order of $`\alpha (Z\alpha )^2m_\mu `$. This shift is a nonrelativistic one and it splits 2s and 2p levels. The fine and hyperfine structures are of the same form as for the normal atoms (i. e. $`(Z\alpha )^4m_\mu `$ and $`(Z\alpha )^4m_\mu ^2/M`$ respectively) and at low $`Z`$ the Lamb shift is a dominant correction to the energy levels.
## Hydrogen Lamb shift
A number of different splittings have been precisely studied for about a century. Bound state QED and maybe even QED itself was essentially established after a study of the Lamb shift and the fine and hyperfine structures in hydrogen, deuterium and helium ions. In the last decades, progress with such measurements was quite slow. The results of the last twenty years are presented in Fig 2 (Lamb shift) and 3 (fine structure recalculated in terms of the Lamb shift), while the older experiments are averaged (see Ref. CJP for references). To reach the Lamb shift from the fine structure (2p<sub>3/2</sub>–2s<sub>1/2</sub>) measurement we need to use a value of the 2p<sub>3/2</sub>–2p<sub>1/2</sub> splitting which was found theoretically. The most direct results of the Lamb shift need no QED theory. A result claimed to be the most accurate one (Sokolov) has an uncertainty of about 2 ppm. It is corrected because of a recalculation of the lifetime of the 2p<sub>1/2</sub> state 2p :
$$\tau ^1(2p_{1/2})=\frac{2^{10}\pi }{3^8}\alpha ^3R_{\mathrm{}}\frac{m_R}{m}\left\{1+\mathrm{ln}\left(\frac{9}{8}\right)\left(Z\alpha \right)^2+\frac{\alpha \left(Z\alpha \right)^2}{\pi }\left(8.045\mathrm{}\right)\mathrm{ln}\frac{1}{\left(Z\alpha \right)^2}\right\}.$$
There is some criticism by E. Hinds Hinds and it is not clear if this result is as accurate as claimed. We wish to note, however, that common opinion on the direct Lamb shift measurement contains two contradicting statements. Firstly, it is generally believed that a Lamb splitting of 2s<sub>1/2</sub> and 2p<sub>1/2</sub> (about 1 GHz), with a decay width of 2p being 0.1 GHz, cannot be measured better than 10 ppm. This means that the statistic error should be larger than 10 ppm. Secondly, it is believed that Sokolov’s experiment is incorrect only because of a possible systematic error claimed by Hinds. However, nobody insists that the statistical treatment of Sokolov’s data was incorrect and we can hope that traditional methods can go far beyond 10 ppm level. Measurement of the deuterium Lamb shift within the Sokolov scheme will provide a chance to test some systematics of his experiment.
Essential progress in study of the hydrogen Lamb shift comes recently from the optical two–photon Doppler–free experiments (see Hänsch and Schwob et al. fort detail). The Doppler–free measurement offers a determination of some transition frequency in the gross structure with a accuracy high enough to use the results to find the Lamb shift. However, two problems arise due to these experiments.
The transition energy between different levels of the gross structure is mainly determined by the Rydberg energy: $`R_{\mathrm{}}/n^2`$. To extract the Lamb shift we first need to find a value of the Rydberg constant. There are two ways to manage this. Following the first of them one has to measure two different frequencies within one experiment with the ratio of the frequencies being an integer number. Obtaining a beat frequency one can avoid the problem of determining the Rydberg constant. Three experiments have been performed in this way: the Garching experiment dealt with the 1s–2s transition and the 2s–4s (and 4d), at Yale 1s–2s frequency was compared with the one–photon 2s-4p transition and that was the only precision optical experiment with a one–photon transition. The recent Paris experiment worked with 1s–3s and 2s–6s (and 6d). The values derived from these experiments are collected in Fig. 4.
Another way to manage the problem with the Rydberg constant is to do two independent absolute frequency measurements (i. e. measurements in respect to the primary cesium standard) and to compare them afterwards, hence determining both the Rydberg constant and the Lamb shift. Such an approach, combining results from Garching (1s–2s) and from Paris (on 2s–8s, –8d, –12d), gave another optical value (Fig. 4). Some of the optical experiments were also performed for deuterium and that may improve the accuracy in the determination of the Rydberg constant and, thence, of the hydrogen Lamb shift.
However, the values in Fig. 4 derived from the optical measurements need further theoretical treatment. The experiments involved a number of levels (1s, 2s, 3s etc) and with optical experimental data there was also a problem of an increasing number of levels with an unknown Lamb shifts. The problem was solved with the help of a specific difference JETP94 :
$`\mathrm{\Delta }(n)`$ $`=`$ $`E_L(1s)n^3E_L(ns)`$ (1)
$`=`$ $`{\displaystyle \frac{\alpha (Z\alpha )^4}{\pi }}{\displaystyle \frac{m_R^3}{m^2}}\times \{{\displaystyle \frac{4}{3}}\mathrm{ln}{\displaystyle \frac{k_0(1s)}{k_0(ns)}}(1+Z{\displaystyle \frac{m}{M}})^2+C_{Rec}{\displaystyle \frac{Zm}{M}}`$
$`+`$ $`(Z\alpha )^2[A_{61}\mathrm{ln}{\displaystyle \frac{1}{(Z\alpha )^2}}+A_{60}^{VP}(n)+G_n^{SE}(Z\alpha )+{\displaystyle \frac{\alpha }{\pi }}\mathrm{ln}^2{\displaystyle \frac{1}{(Z\alpha )^2}}B_{62}(n)]\},`$
where the coefficients $`A_{61}`$, $`A_{60}^{VP}(n)`$, $`C_{Rec}`$ and $`B_{62}(n)`$ and a table for $`G_n^{SE}(Z\alpha )`$ can be found in Refs. ZP97 ; CJP . The uncertainty was also discussed there. The difference has a better established status than that for 1s (or 2s) Lamb shifts (see Table 1). The uncertainty budget was improved recently after calculations of one–loop corrections, exactly at $`Z=1`$, (Jentschura et al.) and leading three–loop contributions (Melnikov and van Ritbergen).
The theory of 2p<sub>3/2</sub>–2p<sub>1/2</sub> splitting is also well established. Perhaps, we have to clarify here the word “theoretical”. A value is a theoretical one if it is sensitive to theoretical problems (like the problem of the proton radius and of higher–order QED corrections for the Lamb shift). An insensitive, sterile value is not theoretical, it is rather a mathematical one, and that is the case for the difference $`\mathrm{\Delta }(2)`$ and the 2p<sub>j</sub> energy. Details of theoretical calculations can be found in the review Eides .
## Nuclear structure effects
Now we can compare theory and experiment for the 2s Lamb shift. We summarize them in Fig. 5, where we present average values for the Lamb shift, fine structure and optical beat frequency and comparison experiments. What is important is the influence of the nuclear charge distribution on the energy levels
$$\mathrm{\Delta }E(nl)=\frac{2}{3}\frac{(Z\alpha )^4}{n^3}m^3R_N^2\delta _{l0},$$
(2)
where $`R_N`$ is a mean–squared nuclear charge radius. The position of theoretical values depends on the accepted value for the proton charge radius. We label three theoretical values with the proton radius (0.847 fm – Mainz dispersion analysis paper, 0.805 fm – Stanford scattering experiment, 0.862 fm – Mainz scattering). More values for the proton radius are collected in Fig. 7 (see CJP for references). To discuss the discrepancy let us look at the most important data on electron–proton elastic scattering presented in Fig. 8. One can see that the Mainz experiment is more appropriate to precisely determine the proton radius containing more points at lower momentum transfer and with a higher precision. Due to this any compilation containing the Mainz data has to lead to a result close to the Mainz result, because the Mainz scattering points must be statistically responsible for the final result and, in particular, the dispersion analysis performed by Mainz theorists led to such a result. However it ($`R_p=0.847(9)`$ fm) differs from the empirical value ($`R_p=0.862(12)`$ fm). One problem in evaluating the data is their normalization. One can write a low momentum expansion of the form factor
$$G(q^2)=a_0+a_1q^2+a_2q^4+\mathrm{}$$
(3)
From a theoretical point of view $`G(0)=1`$ indeed. However, the normalization measurement was accurate not enough (an in particular in the Mainz case it is about 0.5%) and that means that a value tabulated from the data, as being the form factor, differed from it with some normalization. Three different fitting were performed by Wong Wong (see Fig. 8). The free fittings of $`a_0`$ led to a larger uncertainty (Wong–Mainz value in Fig. 7). Even this result must be treated as a preliminary value. It is necessary to take into account some higher–order corrections and that is not possible because of the absence of any complete description of the experiment. The reasonable estimate of the theory is presented in Fig. 5 as a filled area. All experimental values are consistent with the theory exept the corrected value from the Sokolov and Yakovlev experiment. The present status is that the computation uncertainty is about 2 ppm, the measurement inaccuracy of the grand average value is 3 ppm, while the uncertainty due to the proton size is about 10 ppm.
The problem of the nuclear size is not only a problem of the hydrogen Lamb shift: a similar situation arises with the helium–4 ion Lamb shift, where uncertainties resulting from the QED computation and the nuclear size are about the same. The comparison of theory to experiment is presented in Fig. 9. The evolution of the measured value has been due to a study of possible systematic sources (Drake and van Wijngaarden).
## Proton–free hydrogen physics
A more difficult problem is that of the hyperfine structure, which is more sensitive to the nuclear structure. While the experimental uncertainty is below $`10^{12}`$, the theoretical inaccuracy is about 10 ppm. The main problem is a distribution of the magnetic moment inside the proton. It seems the scattering cannot provide accurate enough data and we need to discuss how to manage the problem of nuclear structure by means of the atomic physics. We consider here three ways to do that:
* one is based on study of muonic atoms (the muon is a lepton with a lifetime of about 2 $`\mu `$s and a mass of about $`207m_e`$);
* another deals with the special difference $`\mathrm{\Delta }_{hfs}=E_{hfs}(1s)E_{hfs}(2s)`$ (cf. in Eq. (1)), which can be precisily measured;
* the third is for atoms without nuclear structure. In such an atom one must substitute the proton by some more appropriate positive particle (muon or positron).
A promising way is to determine the nuclear structure with muonic atoms and in particular with muonic hydrogen. The muon orbit lies lower than the electron one. Since $`m_\mu 207m_e`$ the muon hydrogen Bohr radius is about 200 times smaller than that in hydrogen and, hence, the former is more sensitive to nuclear effects. A scheme of an experiment running now at PSI (Pohl et al) is presented in Fig. 10. The experiment consists of the following steps: creating a metastable 2s state, exiting it to the 2p state by a laser, measuring the intensity of the X–ray decay 2p–1s. A similar scheme was used for muonic helium muhelium , however a recent experiment by PSI PSIhelium showed no appropriate signal (Fig. 11). A study of the helium experiment revealed a crucial point: the creation of a metastable state, which can be destroyed by collisions. The collision rate is proportional to the target gas density as well as the rate of creating the muonic atom, and so the density cannot be varied arbitrarily. The slow muon beam at PSI allow one to use a low density gas target and creation of the 2s state has been detected. In case of success, the PSI experiment will give us the charge radius of the proton and the so–called Zemach correction to the 2s state of muonic hydrogen. Comparison of the muonic hydrogen hfs and hydrogen hfs will allow us to go farther with the study of the proton structure.
Another way to manage the problem of nuclear structure is to compare the 1s and 2s hfs. The experiments were performed for hydrogen (recently by Rothery and Hessels), deuterium and helium ion. The recent hydrogen experiment has attracted our attention to the problem of $`\mathrm{\Delta }_{hfs}`$ and it was discovered (Karshenboim and Sapirstein) that the results (and primarily those for the helium ion) are quite sensitive to higher order corrections. All value for the hfs used to be presented in units of Fermi energy ($`\nu _F`$), which is the result of a nonrelativistic interaction of the magnetic moment of an electron in the 1s state and the nucleus. The accuracy of the difference allows one to detect the fourth order corrections, namely, $`\alpha (Z\alpha )^3`$, $`\alpha ^2(Z\alpha )^2`$, $`\alpha (Z\alpha )^2m/M`$, and $`(Z\alpha )^3m/M`$.
The same fourth order corrections are now a subject of study in the muonium ground state hfs (see Table 2). A muonium atom is a kind of hydrogen without the proton: instead the proton the nucleus is a positive muon. The present status of the muonium hyperfine splitting is as follows: the experiment at LAMPF gave 4 463 302 765(53) Hz , while the theoretical prediction is consistent with the experiment but less accurate. A computational part of the uncertainty is about 200 Hz, while a hidden experimental uncertainty in the theory is about 500 Hz. It is due to a calculation of the Fermi energy, which is proportional to the muon magnetic moment, determined from the same LAMPF experiment. Possible progress is considered by Jungmann.
Another proton–free simple atom is positronium. Its lifetime is much shorter than that of muonium, but it can be more easily produced. Different measurements in positronium are summarized in Fig. 12–16. Energy levels in positronium can be presented in the form
$$R_{\mathrm{}}\times \left\{C_{20}+C_{40}\alpha ^2+C_{50}\alpha ^2+\left(C_{61}\mathrm{ln}(1/\alpha )+C_{60}\right)\alpha ^2+\left(C_{72}\mathrm{ln}^2(1/\alpha )+\mathrm{}\right)\alpha ^3\right\}.$$
After two decades of intensive theoretical study we know the coefficients in the above expression up to $`C_{60}`$ (Adkins , Czarnecki et al.) and $`C_{72}`$ (JETP93 ; Pachucki ). The decay width in positronium is known up to fractinal order $`\alpha ^2`$ (Czarnecki et al. for parapositronium) (and for orthopositronium Adkins) and $`\alpha ^3\mathrm{ln}^2\alpha `$ JETP93 . Since the Adkins’s result is a preliminary one we do not include it in Fig. 12. The Adkins’s result led to a fractional correction about $`4\alpha ^2`$ and so the theory is in contradiction to the Ann Arbor experiment. Some progress in the study of positronium is expected in the near future (Conti).
## The status of bound state QED
After briefly reviewing the studies for hydrogen, muonium and positronium, let us discuss a problem and current trend of bound state QED. First of all we need to mention that, in our mind, the QED theory as a theory is well established. As a pure theory it is absolutely correct and absolutely useless:
* The QED theory is a theory of interaction between leptons (electrons and muons) and photons only. We need to include hadrons (such as proton) into consideration. Even for the case of pure leptonic values (like for muonium) we need to calculate a hadronic vacuum polarization contribution. So the QED theory is incomplete.
* The QED theory cannot predict anything exactly, but only in terms of expansion and the uncertainty can be presented in terms like $`O(\alpha ^7m)`$. It is necessary to develop an effective approach to estimate uncalculated corrections quantitively in Hz and eV.
* The QED deals rather with free particles and it is necessary to develop an effective approach to solve a bound state problem for two bodies.
These three problems lie beyond QED as a mathematical theory, but are an essential part of any real QED calculations. A test of different effective approaches is a real problem, as is evaluating of hadronic contributions needed for the precision theory of simple atoms.
The bound state problem has mainly three small parameters: $`\alpha `$ (associated with QED effects), $`Z\alpha `$ (due to binding effects) and $`m/M`$ (the recoil parameter). Now, for the first time, it is necessary to try to really study effects which involve essential QED, two–body and binding effects simultaneously (the $`\alpha (Z\alpha )^2m/M`$ corrections). That is a problem for the hyperfine structure of muonium, the 1s/2s hyperfine structure in hydrogen, deuterium and He<sup>+</sup> and for the positronium spectrum. Another crucial problem is that of the Lamb shift in hydrogen and light ions: this is a higher order two–loop corrections ($`\alpha ^2(Z\alpha )^6m`$) already known in part. We summarize all crucial terms in Table 3.
The three parameters we mention generate different expansions and it is found that all three kind of expansions are not free from problems.
* The QED expansion over $`\alpha `$ is an asymptotic one and the value of terms will decrease to some $`n_c`$ and increase after it. Fortunately, that is not important for $`n=13`$, which are only actual for the bound state QED.
* The $`Z\alpha `$ expansion involves another problem. One knows that high $`Z`$ is a bad limit (strong coupling) and it is believed that low $`Z`$ is a good limit. The latter is wrong. It is clear that at $`Z=0`$ there is no bound system at all and the behavior of any expansion in the limit of low $`Z`$ is not an analytical one. This eventually leads to logarithmic contributions. Even a cube of logarithm ($`\mathrm{ln}^3(1/\alpha )120`$ at $`Z=1`$) appears. The imaginary part of the logarithm is $`\pi `$ and the non–leading terms have often large coefficients because of this. There is another mechanism for large coefficients. As it is well–known the Bethe logarithm ($`\mathrm{ln}\left(k_0(ns)\right)3`$) is a logarithm of an effective energy (in atomic units) of an intermediate state in a calculation of the electron self–energy. A logarithm equal to 3 corresponds to a quite relativistic intermediate p–state ($`v/c4.5(Z\alpha )`$) and that also leads to large coefficients because of relativistic corrections for intermediate states.
* The recoil effect with the $`m/M`$ expansion also involves a non–analytical behavior. It is correct that the limit $`m_1=m_2`$ (positronium) is rather complicated for a calculation, however at the opposite limit ($`m/M0`$) there is no bound state. Hopefully, often the logarithmic recoil corrections are not quite important numerically. Since most of the recoil effects are relativistic ones, the exchange loop generates the effective parameter $`(Z\alpha )/\pi `$ rather than $`Z\alpha `$.
Due to the increasing number of logarithmic contributions we end up with a problem of large higher–order corrections. Some higher–order logarithmic terms are compatible in comparison to a constant part of some lower–order terms.
* For the hydrogen 2s Lamb shift, non–logarithmic parts of the fifth order corrections (in unit of the Rydberg contributions) lie from 164 kHz ($`\alpha (Z\alpha )^6m`$) and 37 kHz ($`\alpha ^2(Z\alpha )^6m`$) to a few kHz for recoil terms. The leading logarithmic term in the next order is $`\alpha ^2(Z\alpha )^6\mathrm{ln}^3(Z\alpha )`$, which contributs 3.6 kHz.
* The non–logarithmic parts of the third order correction (in unit of $`\nu _F`$) for the ground state muonium hyperfine splitting varies from 8.8 kHz ($`\alpha (Z\alpha )^2`$) to 2 kHz for recoil and radiative recoil terms and to 0.4 kHz for $`\alpha ^2(Z\alpha )`$. Three leading logarithmic corrections are slightly below 1 kHz (see Table 2): $`\alpha (Z\alpha )^3`$, $`\alpha (Z\alpha )^2(m/M)`$ and $`(Z\alpha )^3(m/M)`$.
* A number of positronium levels are under study. The non-logarithmic $`\alpha ^6m`$ term (7.2 MHz) for the hyperfine structure is bigger than the logarithmic part of the $`\alpha ^7m`$ contribution (0.9 MHz), while for the 1s-2s interval the situation is different: the non-logarithmic $`\alpha ^6m`$ term is only 0.5 MHz and that is less than 1.2 MHZ of the $`\alpha ^7m`$.
These example show that an estimation of the higher order terms is extremely important and we hope that a calculation of leading logarithmic contributions provides a reasonable way to estimate uncalculated terms. We estimate the non–leading term within a half–value of the leading logarithmic contributions.
Estimation of uncalculated terms is a crucial problem in any QED calculations. Let us now mention the case of moderate Z. Study of these ions provides a unique possibility of measuring higher order corrections. In particular an experimental study of helium (Drake, Burrows et al.) and nitrogen (Myers) hydrogen–like ions will allow us to extract information on higher–order two–loop contributions with the help of a theoretical study of all other terms (Ivanov and Karshenboim). Moderate–$`Z`$ few–electron atoms allow us to test our understanding of higher–order electron–electron interactions which is important for high $`Z`$ spectroscopy.
Large values of the higher–order terms imply a calculation without expansion. This is only possible for one parameter, either $`Z\alpha `$ or $`m/M`$. For the simplest corrections (like e. g. the vacuum polarization) it is possible to calculate analytically to any order of $`Z\alpha `$, otherwise only numerical results are possible. It is unlikely that a complete exact calculation of the two-loop self-energy can be performed soon and that means that expansion techniques is still the main approach in calculating the higher-order terms, perhaps in combination with experiments.
A new opprortunity appears due to recent measurements of a bound electron $`g`$–factor in a hydrogen–like atom (Häffner et al., Quint ). A recent result on the carbon ion is useful for indirect determination of the electron mass (Karshenboim). Our theoretical prediction
$$g_b(e)=2\left(1+520795(1)10^9\right)$$
(4)
mainly based on Beier calculation has a smaller uncertainty in part because of taking into account a known $`\alpha ^2(Z\alpha )^2`$ term. Studies of the $`g`$–factor will be very different from the Lamb shift and the hyperfine structure. In contrast to spectroscopic studies it is possible to go through all $`Z`$ and to determine some unknown coefficients of the theoretical expansion if we can fix its shape (we call this weak theory in contrast to a real theory which can give direct numerical predictions).
## Summary
Concluding the paper we wish to mention briefly different applications for the study of simple atoms. These studies are important for different field of physics:
* Determination of fundamental constants ($`R_{\mathrm{}}`$, $`\alpha `$, $`m_e/m_p`$): some of these are important for other application, like e. g. the fine structure constant is necessary to reproduce the value of the Ohm from the quantum Hall effect.
* Development of new optical standard and tool for same: like e. g. the new frequency chain designed recently (see Diddam et al., Udem et al. and Holzwarth ).
* New physics: a study of muonium-to-antimuonium conversion (Jungmann) and the exotic decay of positronium (Conti) provide a possibility of looking for new particles, while the antihydrogen project and some others are expected to test some symmerties or possible variations of fundamental constants.
* Particle physics: a recent study of hydrogen is rather important to learn the proton structure, than to test the bound state QED. The theoretical methods developped recently are of use both for atoms and two-quark particles (mesons) and QED is a good opprotunity to test the methods. Exotic atoms give us accurate information on hadron-hadron interactions.
* Nuclear physics: the situation is similar to that for particles physics: a study of light atoms offers information on structure of their nuclei. One the other hand, two– and three–body atoms are an appropriate problem to testing different effective methods before to applying them to light nuclei.
Most of these questions and a more broad range of problems in physics of simple atoms were considered at a Satellite meeting to the ICAP (Hydrogen Atom, 2: Precise Physics of Simple Atomic System). The Proceeding will be published by Springer in 2001.
## Aknowledgement
I am grateful to T. Hänsch, K. Jungmann, G. Werth, J. Sapirstein and E. Myers for useful and stimulating discussions. I would like to thank S. Nic Chormaic and F. Cataliotti for useful remarks on manuscript. The work was supported in part by RFBR (grant 00-02-16718), NATO (CRG 960003) and Russian State Program Fundamental Metrology.
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# “Classical” flag varieties for quantum groups: the standard quantum SL(𝑛,𝐂)
## Introduction
The study of quantum analogues of flag varieties, first suggested by Manin , has been undertaken during the past decade by several authors, from various points of view; see e.g. . Around the same time, an approach to noncommutative projective algebraic geometry was initiated by Artin, Tate, and Van den Bergh , and considerably developed since (see e.g. ). One attractive feature of their approach is the association of actual *geometric* data to certain classes of graded *non*commutative algebras.
The present work is an attempt to study quantum flag varieties from this point of view. As a consequence, our “quantum flag varieties” will be actual varieties (with some bells and whistles).
Recall the original idea of . If $`A`$ is the homogeneous coordinate ring of a projective scheme $`E`$, then the points of $`E`$ are in one-to-one correspondence with the isomorphism classes of so-called *point modules* of $`A`$, i.e. $`𝐍`$-graded cyclic $`A`$-modules $`P`$ such that $`dimP_n=1`$ for all $`n`$. Now if $`A`$ is an $`𝐍`$-graded *non*commutative algebra, one may still try to parametrize its point modules by the points of some projective scheme $`E`$. Of course, one cannot hope to reconstruct $`A`$ from $`E`$ alone, but there is now an additional ingredient: the *shift* operation $`\sigma :PP[1]`$, where $`P[1]`$ is the $`𝐍`$-graded $`A`$-module defined by $`P[1]_n:=P_{n+1}`$. (When $`A`$ is commutative, this shift is trivial: $`P[1]P`$ for every point module $`P`$.) Assume that $`\sigma `$ may be viewed as an automorphism of $`E`$: one may then hope, at least in “good” cases, to recover $`A`$ from the triple $`(E,\sigma ,)`$, where $``$ is the line bundle over $`E`$ defined by its embedding into a projective space. The first step of this recovery is the construction of the *twisted homogeneous coordinate ring* $`B(E,\sigma ,)`$ of a triple $`(E,\sigma ,)`$, defined in as follows:
$$B(E,\sigma ,)=\underset{n𝐍}{}B_n,B_n:=\mathrm{H}^0(E,^\sigma \mathrm{}^{\sigma ^{n1}})$$
(where $`^\sigma `$ denotes the pullback of $``$ along $`\sigma `$), the multiplication being given by $`\alpha \beta :=\alpha \beta ^{\sigma ^m}`$ for all $`\alpha B_m`$, $`\beta B_n`$. (When $`\sigma `$ is the identity, this algebra is the (commutative) homogeneous coordinate ring of $`E`$ w.r.t. the polarization $``$.) If the triple $`(E,\sigma ,)`$ comes from an algebra $`A`$ as above, the second step then consists in analysing the canonical morphism $`AB(E,\sigma ,)`$. The initial success of this method has been a complete study of all regular algebras of dimension three (where the kernel of $`AB(E,\sigma ,)`$ turns out to be generated by a single element of degree three).
The present paper is organized as follows.
In Part I, we give a general outline of a possible theory of flag varieties for quantum groups, using a multigraded version of the ideas of recalled above, some of which have already been introduced by Chan . This Part is largely conjectural and contains no (significant) new results; its purpose is rather to set up a framework that will be used in Parts II and III.
More specifically, we proceed as follows. Let $`G`$ be a simple complex group; our interest in flag varieties allows us to assume without harm that $`G`$ is simply connected. Let $`P^+`$ be the monoid of dominant integral weights of $`G`$ (w.r.t. some Borel subgroup $`BG`$): the *shape algebra* $`M`$ of $`G`$ is a $`P^+`$-graded $`G`$-algebra whose term of degree $`\lambda `$ is the irreducible representation of $`G`$ of highest weight $`\lambda `$. Now consider the definition of a point module (see above), but with $`𝐍`$-gradings replaced by $`P^+`$-gradings: we obtain the notion of a *flag module* of $`M`$ (Definition 2.1); this terminology is justified by the fact that the isomorphism classes of such modules are indeed parametrized by the points of the flag variety $`G/B`$ (Proposition 2.2).
If a quantum group has the “same” representation theory as $`G`$ (in the sense of Definition 1.1), then the we may still define a ($`P^+`$-graded) shape algebra. We then discuss the possibility to parametrize the latter’s flag modules (up to isomorphism) by the points of some scheme $`E`$, and to realize shifts $`FF[\lambda ]`$ ($`\lambda P^+`$) as automorphisms $`\sigma _\lambda `$ of $`E`$. It will of course be sufficient to know the automorphisms $`\sigma _1,\mathrm{},\sigma _{\mathrm{}}`$ associated to the fundamental weights $`\varpi _1,\mathrm{},\varpi _{\mathrm{}}`$, which freely generate $`P^+`$. Moreover, since we are in a *multi*graded situation, it will be more natural to view $`E`$ as a subscheme of a *product* of $`\mathrm{}`$ projective spaces, corresponding to $`\mathrm{}`$ line bundles $`_1,\mathrm{},_{\mathrm{}}`$ over $`E`$.
We then consider the converse problem of reconstructing the shape algebra from $`E`$, the $`\sigma _i`$, and the $`_i`$, using Chan’s construction of a twisted *multi*homogeneous coordinate ring: this is the $`P^+`$-graded algebra
$$B(E,\sigma _1,\mathrm{},\sigma _{\mathrm{}},_1,\mathrm{},_{\mathrm{}}):=\underset{\lambda P^+}{}\mathrm{H}^0(E,_\lambda ),$$
where the line bundles $`_\lambda `$ are constructed inductively from the rules $`_{\varpi _i}=_i`$, $`_{\varpi _i+\lambda }=_i_\lambda ^{\sigma _i}`$. (Again, if $`E=G/B`$, if each $`\sigma _i`$ is the identity, and if the $`_i`$ are the line bundles associated to the fundamental $`G`$-modules $`V^1,\mathrm{},V^{\mathrm{}}`$, then this algebra is the (commutative) multihomogeneous coordinate ring of $`G/B(V^1)\times \mathrm{}\times (V^{\mathrm{}})`$, which in turn is equal to the shape algebra $`𝒪(\overline{G/U})`$, $`U`$ the unipotent radical of $`B`$.)
We stress that the ideas developed in this Part are *not* restricted to the standard quantum groups of Drinfel and Jimbo , but could, in principle, be applied to other quantum groups as well, as long as they have the “same” representation theory as a given simple complex group. (Potential other examples include the multiparameter quantum groups of Artin, Schelter, and Tate , the quantum $`\mathrm{SL}(n)`$ of Cremmer and Gervais , or the quantum $`\mathrm{SL}(3)`$’s classified in .)
In Parts II and III, we *do* restrict ourselves to a standard Drinfeld-Jimbo quantum group $`𝒪_q^{\mathrm{DJ}}(G)`$, with $`q`$ not a root of unity. Thanks to the results of Lusztig and Rosso , $`𝒪_q^{\mathrm{DJ}}(G)`$ has the “same” representation theory as the group $`G`$, so one can define a shape algebra $`M^{\mathrm{DJ}}`$.
In Part II, we construct geometric data $`E^{\mathrm{DJ}}`$, $`\sigma _i`$, and $`_i`$, and we conjecture that these data indeed correspond to $`M^{\mathrm{DJ}}`$ as described above (Conjectures 9.1 and 9.2). The scheme $`E^{\mathrm{DJ}}`$ will actually be a union of certain $`T`$-stable subvarieties of the (ordinary) flag variety $`G/B`$. Since the latter may be of independent interest to algebraic geometers, we have decided to describe them in a separate note (but we recall their construction here, without proofs).
In Part III, we prove Conjecture 9.1 for $`G=\mathrm{SL}(n)`$, thus obtaining a “flag variety” for the standard Drinfeld-Jimbo quantum $`\mathrm{SL}(n)`$. The proof uses special features of the group $`\mathrm{SL}(n)`$ (the Weyl group is the symmetric group, all fundamental representations are minuscule, etc.) and is essentially combinatorial; it is therefore not likely to be extendable to an arbitrary $`G`$.
Acknowledgement. The author would like to thank the Université de Reims for granting a sabbatical leave during the year 1999–2000, when part of this work has been done.
Conventions. All vector spaces, dimensions, algebras, tensor products, varieties, schemes, etc. will be over the field $`𝐂`$ of complex numbers. If $`G`$ is a linear algebraic group, we denote by $`𝒪(G)`$ the Hopf algebra of polynomial functions on $`G`$. If $`A`$ is a (co)algebra, then the dual of a left $`A`$-(co)module is a right $`A`$-(co)module, and vice-versa; morphisms of $`A`$-(co)modules will simply be called *$`A`$-morphisms*. When $`V`$ is a vector space and $`vV`$ is nonzero, we will sometimes still denote by $`v`$ the corresponding point in the projective space $`(V)`$.
## Part I An approach to quantum flag varieties: general outline
This Part contains no (significant) new results. It rather discusses a possible theory of flag varieties for simple quantum groups, asking several questions along the way (as well as two ambitious problems at the end).
Most of what we will say here is a multigraded version of some of the main ideas of , applied in a Lie-theoretic setting.
### 1. Simple quantum groups and their shape algebras
Let $`G`$ be a simply connected simple complex group, $`BG`$ a Borel subgroup, and $`P^+`$ the set of dominant integral weights of $`G`$ w.r.t. $`B`$. For each $`\lambda ,\mu ,\nu P^+`$, denote by
* $`d^\lambda `$ the dimension of the simple $`G`$-module of highest weight $`\lambda `$, and by
* $`c_\nu ^{\lambda \mu }`$ the multiplicity of the simple $`G`$-module of highest weight $`\nu `$ inside the tensor product of those of highest weights $`\lambda `$ and $`\mu `$.
Bearing in mind that the algebra $`𝒪(G)`$ of polynomial functions on $`G`$ is a commutative Hopf algebra, and that (finite-dimensional) left $`G`$-modules correspond to right $`𝒪(G)`$-comodules, recall the following definition from .
###### Definition 1.1.
We call a *quantum $`G`$* any (not necessarily commutative) Hopf algebra $`A`$ (over $`𝐂`$) such that
1. there is a family $`\{V^\lambda \lambda P^+\}`$ of simple and pairwise nonisomorphic (right) $`A`$-comodules, with $`dimV^\lambda =d^\lambda `$,
2. every $`A`$-comodule is isomorphic to a direct sum of these,
3. for every $`\lambda ,\mu P^+`$, $`V^\lambda V^\mu `$ is isomorphic to $`_\nu c_\nu ^{\lambda \mu }V^\nu `$.
For convenience, we will write
$$V_\lambda :=(V^\lambda )^{}.$$
For every $`\lambda ,\mu P^+`$, Definition 1.1(c) yields an injective $`A`$-morphism $`V^{\lambda +\mu }V^\lambda V^\mu `$ that is unique up to scalars. Denote by
$$m_{\lambda \mu }:V_\lambda V_\mu V_{\lambda +\mu }$$
the corresponding projection. Gluing these together on
$$M_A:=\underset{\lambda }{}V_\lambda ,$$
we get a (not necessarily associative) multiplication $`m:M_AM_AM_A`$.
###### Definition 1.2.
The algebra $`M_A`$ is called the *shape algebra* of $`A`$.
###### Question A.
Is it possible to renormalize the $`m_{\lambda \mu }`$ in such a way that the multiplication $`m`$ becomes associative?
Recall that this Question has a positive answer in the commutative case $`A=𝒪(G)`$: if $`U`$ is a maximal unipotent subgroup, then by the Borel-Weil theorem, we may set
$$M_{𝒪(G)}=𝒪(\overline{G/U}):=\{f𝒪(G)f(gu)=f(g)gG,uU\}.$$
The next Proposition provides a criterion for a positive answer to Question A. We first introduce some more notation: let $`\mathrm{}`$ be the rank of $`G`$, denote by $`\varpi _1,\mathrm{},\varpi _{\mathrm{}}`$ the fundamental weights, and let us use the shorthand notation
$$V_i:=V_{\varpi _i},V^i:=V^{\varpi _i},1i\mathrm{}.$$
For every $`1i,j,k\mathrm{}`$, Definition 1.1(c) implies that $`V_iV_jV_k`$ contains a unique subcomodule isomorphic to $`V_{\varpi _i+\varpi _j+\varpi _k}`$; denote this subcomodule by $`W_{ijk}`$.
###### Proposition 1.3.
Question *A* has a positive answer (for a given $`A`$) if and only if there exist $`A`$-isomorphisms $`R_{ij}:V_iV_jV_jV_i`$ for all $`i>j`$, such that the braid relation
(1.1)
$$(R_{jk}\mathrm{id})(\mathrm{id}R_{ik})(R_{ij}\mathrm{id})_{W_{ijk}}=(\mathrm{id}R_{ij})(R_{ik}\mathrm{id})(\mathrm{id}R_{jk})_{W_{ijk}}$$
holds for all $`i>j>k`$.
We defer the proof to Appendix A.
###### Corollary 1.4.
Question *A* has a positive answer in each of the following situations:
* when $`G`$ is of rank $`2`$,
* when $`A`$ is dual quasitriangular,
* when $`G=\mathrm{SL}(n)`$ (by the main result of ).
Since $`\varpi _1,\mathrm{},\varpi _{\mathrm{}}`$ generate the monoid $`P^+`$, and since the $`m_{\lambda \mu }`$ are surjective, the algebra $`M_A`$ is generated by
$$M_1:=V_1\mathrm{}V_{\mathrm{}}.$$
In this way, $`M_A`$ may be viewed as an $`𝐍`$-graded algebra. More explicitly, if $`\lambda P^+`$ decomposes as $`_ia_i\varpi _i`$ (with each $`a_i𝐍`$), and if we write $`h(\lambda ):=a_i`$ for the *height* of $`\lambda `$, then the $`𝐍`$-grading on $`M_A`$ is given by $`M_k:=_{h(\lambda )=k}V_\lambda `$.
###### Question B.
Is the shape algebra $`M_A`$ quadratic (as an $`𝐍`$-graded algebra)?
In the commutative case $`A=𝒪(G)`$, the shape algebra $`𝒪(\overline{G/U})`$ is indeed quadratic by a well known theorem of Kostant (see \[27, Theorem 1.1\] for a proof). This remains true for the standard Drinfeld-Jimbo quantum $`\mathrm{SL}(n)`$: a presentation of the corresponding shape algebra by generators and (quadratic) relations has been given by Taft and Towber .
###### Question C.
Is the shape algebra $`M_A`$ a Koszul algebra?
To finish this Section, let us take a closer look at the quadratic relations in $`M_A`$. For every $`1i,j\mathrm{}`$, let $`K_{ij}`$ be the kernel of the multiplication $`V_iV_jV_{\varpi _i+\varpi _j}`$. By Definition 1.1(c), the $`A`$-comodules $`V_iV_j`$ and $`V_jV_i`$ are isomorphic, and, rescaling the $`A`$-isomorphism $`R_{ij}:V_iV_jV_jV_i`$ of Proposition 1.3 if necessary, we may assume that the diagram (A.1) (in Appendix A) commutes. Using Definition 1.1(c), we see that the quadratic relations in $`M_A`$ of degree $`\varpi _i+\varpi _j`$ are of two kinds:
1. $`\xi =0`$, for $`\xi K_{ij}`$;
2. $`\xi =R_{ij}(\xi )`$, for $`\xi V_iV_j`$.
###### Remark 1.5.
By Definition 1.1(c), relations (I)<sub>ij</sub> and (II)<sub>ij</sub> for arbitrary $`i,j`$ are consequences of relations (I)<sub>ij</sub> for $`ij`$ only and relations (II)<sub>ij</sub> for $`i>j`$ only.
### 2. The scheme of flag modules
Assume that Question A has a positive answer. The following definition is a multigraded analogue of the point modules introduced in .
###### Definition 2.1.
A *flag module* is a $`P^+`$-graded right $`M_A`$-module $`F`$ such that
1. $`dimF_\lambda =1`$ for each $`\lambda P^+`$,
2. $`F`$ is cyclic.
The terminology is justified by the commutative case. Indeed, let $`BG`$ be a Borel subgroup and $`U`$ the unipotent radical of $`B`$. Then we have the following
###### Proposition 2.2.
The isomorphism classes of flag modules of $`M_{𝒪(G)}=𝒪(\overline{G/U})`$ are parametrized by the points of the flag variety $`G/B`$.
###### Proof.
First, recall from the Borel-Weil theorem that the decomposition $`𝒪(\overline{G/U})=_{\lambda P^+}V_\lambda `$ is given by
$$V_\lambda =\{f𝒪(G)f(gb)=\lambda (b)f(g)gG,bB\}.$$
Now fix $`gG`$ and endow a vector space $`F=_{\lambda P^+}𝐂e_\lambda `$ with the flag module structure defined by
$$e_\lambda .f=f(g)e_{\lambda +\mu }\text{for all }fV_\mu \text{.}$$
If we replace $`g`$ by $`gb`$ for some $`bB`$, the expression for $`e_\lambda .f`$ is just multiplied by $`\mu (b)`$, so up to isomorphism (of graded modules), the flag module thus obtained only depends on the class $`gBG/B`$.
Conversely, assume that $`F`$ is a flag module of $`𝒪(\overline{G/U})`$, and choose a graded basis $`\{e_\lambda \lambda P^+\}`$ of $`F`$. For each $`\lambda ,\mu P^+`$, let $`v_\lambda ^\mu V^\mu `$ be defined by
$$e_\lambda .f=f,v_\lambda ^\mu e_{\lambda +\mu }\text{for all }fV_\mu \text{.}$$
Since the algebra $`𝒪(\overline{G/U})`$ is commutative, we have $`(e_0.f).f^{}=(e_0.f^{}).f`$ for every $`fV_\lambda `$, $`f^{}V_\mu `$, hence
(2.1)
$$v_0^\lambda v_\lambda ^\mu =v_\mu ^\lambda v_0^\mu .$$
It follows that $`v_\lambda ^\mu `$ is a multiple of $`v_0^\mu =:v^\mu `$, say $`v_\lambda ^\mu =a_\lambda v^\mu `$. Inserting back into (2.1) yields $`a_\lambda =a_\mu `$, for all $`\lambda ,\mu P^+`$. Since $`a_0=1`$, we get $`a_\lambda =1`$ for all $`\lambda P^+`$. Therefore,
$$e_0.f=f,v^\mu e_\lambda \text{for all }fV_\lambda \text{.}$$
The collection $`\{v^\lambda \lambda P^+\}`$ defines a linear form $`v`$ on $`𝒪(\overline{G/U})`$. Furthermore, we have $`e_0.(ff^{})=(e_0.f).f^{}`$ for all $`fV_\lambda `$, $`f^{}V_\mu `$, so $`ff^{},v^{\lambda +\mu }=f,v^\lambda f^{},v^\mu `$, which shows that the linear form $`v`$ is a character on $`𝒪(\overline{G/U})`$, corresponding to a point $`x`$ of the affine variety $`\overline{G/U}`$. Moreover, since $`F`$ is cyclic, each $`v^\lambda `$ must be nonzero, so $`x`$ actually lies in $`G/U`$, say $`x=gU`$. This yields an element $`gBG/B`$.
It is clear that these two constructions are inverse to each other. ∎
We will now discuss a possible picture of this kind in the noncommutative situation: if $`A`$ is a quantum $`G`$, we would like to parametrize the isomorphism classes of flag modules over the shape algebra $`M_A`$ by the (closed) points of some scheme $`E`$.
Moreover, given a flag module $`F`$ and a weight $`\lambda P^+`$, consider the *shifted* flag module $`F[\lambda ]`$, defined as the $`P^+`$-graded module such that $`F[\lambda ]_\mu :=F_{\lambda +\mu }`$. We would then like that, for each $`\lambda `$, the shift operation $`FF[\lambda ]`$ corresponds to an automorphism of schemes $`\sigma _\lambda :EE`$.
To achieve this, let us encode the structure of a flag module more geometrically, as follows. If $`F`$ is a flag module with basis $`\{e_\lambda \lambda P^+\}`$, then for each $`\lambda P^+`$ and each $`1i\mathrm{}`$, let $`v_\lambda ^iV^i`$ be defined by
(2.2)
$$e_\lambda .f=f,v_\lambda ^ie_{\lambda +\varpi _i}\text{for all }fV_i\text{.}$$
Replacing $`F`$ by an isomorphic flag module (i.e. rescaling the $`e_\lambda `$) only multiplies each $`v_\lambda ^i`$ by a scalar, so let $`p_\lambda ^i`$ be the corresponding point in $`(V^i)`$. To simplify notation, let us write
$$^1\mathrm{}\mathrm{}:=(V^1)\times \mathrm{}\times (V^{\mathrm{}})$$
and denote by $`\mathrm{pr}^i:^1\mathrm{}\mathrm{}(V^i)`$ the natural projection. For any point $`p^1\mathrm{}\mathrm{}`$, we use the shorthand notation $`p^i:=\mathrm{pr}^i(p)`$. Thus, to an isomorphism class of flag modules, we associate a collection of points $`\{p_\lambda \lambda P^+\}`$ in $`^1\mathrm{}\mathrm{}`$.
From now on, we assume that Question B has a positive answer. The quadratic relations (I) and (II) in $`M_A`$ (see the end of Section 1) impose some conditions on this collection of points, which we now analyse.
For relations of type (I), identify $`(V^i)\times (V^j)`$ with its image in $`(V^iV^j)`$ under the Segre embedding. Relations (I)<sub>ij</sub> may be viewed as equations defining a subscheme $`\mathrm{\Gamma }^{ij}`$ of $`(V^i)\times (V^j)`$. We then have
(2.3)
$$(p_\lambda ^i,p_{\lambda +\varpi _i}^j)\mathrm{\Gamma }^{ij}$$
for all $`\lambda P^+`$ and all $`1i,j\mathrm{}`$.
Similarly, for relations of type (II), we consider the map $`(R^{ji}):(V^jV^i)(V^iV^j)`$, where $`R^{ji}`$ denotes the transpose of $`R_{ij}`$. Then we must have
(2.4)
$$(p_\lambda ^i,p_{\lambda +\varpi _i}^j)=(R^{ji})(p_\lambda ^j,p_{\lambda +\varpi _j}^i)$$
(again identifying $`(V^i)\times (V^j)`$ with its image under the Segre embedding).
Gluing together conditions (2.3) and (2.4) for all $`i,j`$, we are led to consider the subscheme $`\mathrm{\Gamma }(^1\mathrm{}\mathrm{})^{\mathrm{}+1}`$ of all $`(\mathrm{}+1)`$-tuples $`(p_0,p_1,\mathrm{},p_{\mathrm{}})`$ satisfying
$$(p_0^i,p_i^j)\mathrm{\Gamma }^{ij},$$
$$(p_0^i,p_i^j)=(R^{ji})(p_0^j,p_j^i)$$
for all $`1i,j\mathrm{}`$. We may now rephrase conditions (2.3) and (2.4) by saying that the collection $`\{p_\lambda \lambda P^+\}`$ satisfies
(2.5)
$$(p_\lambda ,p_{\lambda +\varpi _1},\mathrm{},p_{\lambda +\varpi _{\mathrm{}}})\mathrm{\Gamma }\text{for all }\lambda P^+\text{.}$$
###### Proposition 2.3.
Assume that $`M_A`$ is quadratic (as an $`𝐍`$-graded algebra). Then there is a one-to-one correspondence between isomorphism classes of flag modules over $`M_A`$ and families $`\{p_\lambda \lambda P^+\}`$ of points of $`^1\mathrm{}\mathrm{}`$ satisfying (2.5).
###### Proof.
It remains to show that the above construction can be reversed, so assume that $`\{p_\lambda \lambda P^+\}`$ is a collection of points in $`^1\mathrm{}\mathrm{}`$ satisfying (2.5). Choose a (nonzero) representative $`v_\lambda ^iV^i`$ for each $`p_\lambda ^i`$, and endow a vector space $`_{\lambda P^+}𝐂e_\lambda `$ with the flag module structure defined by the rule (2.2). By (2.3), this rule is compatible with relations of type (I) in $`M_A`$. By (2.4), it is also compatible with relations of type (II), *provided* that, for each $`\lambda P^+`$ and each $`1i,j\mathrm{}`$, we suitably rescale one of $`v_\lambda ^i`$, $`v_{\lambda +\varpi _i}^j`$, $`v_\lambda ^j`$, $`v_{\lambda +\varpi _j}^i`$. Proceeding by induction over the height $`h(\lambda )`$, we may perform this rescaling in a consistent way.
It is clear that the two constructions are inverse to each other. ∎
###### Remark 2.4.
Rescaling the $`m_{\lambda \mu }`$ only multiplies the $`R_{ij}`$ by scalars. Therefore, the scheme $`\mathrm{\Gamma }`$ does not depend on the normalizations of the multiplication in $`M_A`$, but only on $`A`$ itself.
The following Question is inspired by the description given in the Introduction of .
###### Question D.
Do there exist a subscheme $`E`$ of $`^1\mathrm{}\mathrm{}`$ and $`\mathrm{}`$ pairwise commuting automorphisms $`\sigma _1,\mathrm{},\sigma _{\mathrm{}}:EE`$ such that the scheme $`\mathrm{\Gamma }`$ is given by
(2.6)
$$\mathrm{\Gamma }=\{(p,\sigma _1(p),\mathrm{},\sigma _{\mathrm{}}(p))pE\}\mathrm{?}$$
A positive answer to this Question would fulfill the aim of parametrizing flag modules, as expressed at the beginning of this Section. Indeed, assume that $`E`$ and $`\sigma _1,\mathrm{},\sigma _{\mathrm{}}`$ as in Question D do exist. For each weight $`\lambda =a_i\varpi _i`$, define $`\sigma _\lambda :=\sigma _1^{a_1}\mathrm{}\sigma _{\mathrm{}}^a_{\mathrm{}}`$; since the $`\sigma _i`$ commute, we have $`\sigma _{\lambda +\mu }=\sigma _\lambda \sigma _\mu `$. Then for every family $`\{p_\lambda \lambda P^+\}`$ of points in $`^1\mathrm{}\mathrm{}`$ satisfying (2.5), the realization (2.6) shows that $`p_\lambda =\sigma _\lambda (p_0)`$ for all $`\lambda P^+`$, with $`p_0E`$. Conversely, for every $`pE`$, the family $`\{\sigma _\lambda (p)\lambda P^+\}`$ satisfies (2.5) and thus defines an isomorphism class of flag modules by Proposition 2.3. Therefore, if Question D had a positive answer, flag modules (up to isomorphism) would be parametrized by the points of $`E`$, with $`\sigma _\lambda `$ corresponding to the shift operation $`FF[\lambda ]`$.
Finally, for future reference, we define, for each $`1i\mathrm{}`$, the line bundle $`_i`$ over $`E`$ as the pullback of $`𝒪_{(V^i)}(1)`$ along $`\mathrm{pr}^i`$ (restricted to $`E`$), and we call the tuple
$$\mathrm{T}(M_A):=(E,\sigma _1,\mathrm{},\sigma _{\mathrm{}},_1,\mathrm{},_{\mathrm{}})$$
the *flag tuple* associated to $`A`$.
### 3. Braided tuples and reconstruction of shape algebras
Chan has given a construction in the opposite direction, starting from a scheme $`E`$, automorphisms $`\sigma _1,\mathrm{},\sigma _{\mathrm{}}`$ of $`E`$, and line bundles $`_1,\mathrm{},_{\mathrm{}}`$ over $`E`$ (satisfying some compatibility conditions; see Definition 3.1), and building a $`P^+`$-graded algebra from these data. Let us briefly recall his construction. To improve legibility, we will write $`^\sigma `$ for the pullback of a line bundle $``$ along a map $`\sigma `$.
###### Definition 3.1.
We call a tuple $`T=(E,\sigma _1,\mathrm{},\sigma _{\mathrm{}},_1,\mathrm{},_{\mathrm{}})`$ as above a *braided tuple* if
1. the $`\sigma _i`$ pairwise commute,
2. for every $`i>j`$, there exists an equivalence $`R_{ij}:_i_j^{\sigma _i}\stackrel{}{}_j_i^{\sigma _j}`$ of line bundles such that the braid relation
(3.1)
$$(R_{jk}\mathrm{id})(\mathrm{id}R_{ik}^{\sigma _j})(R_{ij}\mathrm{id})=(\mathrm{id}R_{ij}^{\sigma _k})(R_{ik}\mathrm{id})(\mathrm{id}R_{jk}^{\sigma _i})$$
holds for every $`i>j>k`$ (both sides being equivalences $`_i_j^{\sigma _i}_k^{\sigma _i\sigma _j}\stackrel{}{}_k_j^{\sigma _k}_i^{\sigma _k\sigma _j}`$).
Note that if we set $`R_{ii}:=\mathrm{id}`$ for all $`i`$ and $`R_{ji}:=R_{ij}^1`$ for all $`i>j`$, then (3.1) becomes true for all $`i,j,k`$.
If $`\lambda P^+`$ decomposes as $`\lambda =a_i\varpi _i`$, then define $`\sigma _\lambda :=\sigma _1^{a_1}\mathrm{}\sigma _{\mathrm{}}^a_{\mathrm{}}`$, as before (so $`\sigma _{\lambda +\mu }=\sigma _\lambda \sigma _\mu `$). Define a line bundle $`_\lambda `$ over $`E`$ by the following inductive rules (with $`_0`$ the trivial bundle):
(3.2)
$$\begin{array}{c}_{\varpi _i}=_i,1i\mathrm{},\\ _{\lambda +\mu }=_\lambda _\mu ^{\sigma _\lambda }.\end{array}$$
As is shown in , this procedure is, thanks to (3.1), well defined up to unique equivalences of line bundles built from the $`R_{ij}`$ (cf. also the proof of Proposition 1.3). Now define the product of two sections $`\alpha \mathrm{H}^0(E,_\lambda )`$ and $`\beta \mathrm{H}^0(E,_\mu )`$ by
(3.3)
$$\alpha \beta :=\alpha \beta ^{\sigma _\lambda }\mathrm{H}^0(E,_{\lambda +\mu }).$$
###### Theorem 3.2 (Chan ).
The product rule (3.3) turns the direct sum
$$B(T):=\underset{\lambda P^+}{}\mathrm{H}^0(E,_\lambda )$$
into an associative $`P^+`$-graded algebra.
The algebra $`B(T)`$ is not quadratic in general, so we consider its quadratic cover
$$M(T):=B(T)^{(2)}.$$
(If $`B`$ is any $`𝐍`$-graded algebra, we define its *quadratic cover* $`B^{(2)}`$ as follows: consider the canonical homomorphism $`\mathrm{T}(B_1)B`$ and its kernel $`J=_{k2}J_k`$, then set $`B^{(2)}:=\mathrm{T}(B_1)/(J_2)`$. Here we view $`B(T)`$ as an $`𝐍`$-graded algebra via the height function $`h(\lambda )`$.)
The quadratic algebra $`M(T)`$ may also be described more directly in terms of the braided tuple $`T`$, as follows. For each $`1i\mathrm{}`$, set $`V_i:=\mathrm{H}^0(E,_i)`$, denote by $`V^i`$ the dual of $`V_i`$, and consider the map $`\mathrm{Pl}^i:E(V^i)`$ corresponding to the line bundle $`_i`$. For every $`1i,j\mathrm{}`$, the map $`\mathrm{Pl}^i(\mathrm{Pl}^j)^{\sigma _i}`$ corresponding to the line bundle $`_i_j^{\sigma _i}`$ is then given by the composite
(3.4)
$$E\stackrel{\text{diag.}}{}E\times E\stackrel{\mathrm{id}\times \sigma _i}{}E\times E\stackrel{\mathrm{Pl}^i\times \mathrm{Pl}^j}{}(V^i)\times (V^j)\stackrel{\text{Segre}}{}(V^iV^j).$$
Denote by $`\mathrm{\Gamma }^{ij}`$ the image of this map and by $`K_{ij}V_iV_j`$ the subspace of linear forms on $`V^iV^j`$ vanishing on $`\mathrm{\Gamma }^{ij}`$.
For every $`1i,j\mathrm{}`$, Definition 3.1 implies that there exists a linear isomorphism $`R^{ji}:V^jV^iV^iV^j`$ such that the following diagram commutes:
(3.5)
Let $`R_{ij}:V_iV_jV_jV_i`$ be the transpose of $`R^{ji}`$. It is clear that modulo $`K_{ij}`$ and $`K_{ji}`$, the map $`R_{ij}`$ is unique up to a scalar.
The algebra $`M(T)`$ is then generated by $`V_1\mathrm{}V_{\mathrm{}}`$, with relations given by (I)<sub>ij</sub> and (II)<sub>ij</sub> for all $`1i,j\mathrm{}`$ (see the end of Section 1; Remark 1.5 still applies).
###### Question E.
What is the kernel of the canonical morphism $`M(T)B(T)`$?
Having constructed the algebra $`M(T)`$ from a braided tuple $`T`$, we may formulate a converse to Question D:
###### Question F.
Assume that $`A`$ is a quantum $`G`$ such that the shape algebra $`M_A`$ is quadratic. Does there exist a braided tuple $`T`$ such that $`M_A=M(T)`$?
This Question is *a priori* weaker than Question D, for the following reason. If $`M_A`$ is quadratic and does admit a flag tuple $`T`$ as in Question D, then the reconstructed algebra $`M(T)`$ is canonically isomorphic to $`M_A`$. However, we might also have $`M(T^{})=M_A`$ for some subtuple $`T^{}`$ of $`T`$ (i.e. a subscheme $`E^{}`$ of $`E`$ stabilized by each $`\sigma _i`$, with $`\sigma _i^{}`$ and $`_i^{}`$ the obvious restrictions).
###### Problem G.
Given a simple complex group $`G`$, characterize the flag tuples of all quantum $`G`$’s intrinsically (i.e. as braided tuples).
For $`G=\mathrm{SL}(2)`$, this is elementary: $`E`$ must be the projective line $`^1`$, $`\sigma `$ can be an arbitrary automorphism of infinite order, and $`=𝒪_^1(1)`$. The three possible forms of $`\sigma `$ correspond to three different quantum $`\mathrm{SL}(2)`$’s, namely, $`𝒪(\mathrm{SL}(2))`$ (when $`\sigma =\mathrm{id}`$), the standard Drinfeld-Jimbo quantum $`\mathrm{SL}(2)`$ for $`q`$ not a root of unity (when $`\sigma `$ has two fixed points), and the Jordanian quantum $`\mathrm{SL}(2)`$ (when $`\sigma `$ has one fixed point). These are known to be the only quantum $`\mathrm{SL}(2)`$’s (in the sense of Definition 1.1). The associated shape algebras are $`𝐂x,y/(xyyx)`$, $`𝐂x,y/(xyqyx)`$, and $`𝐂x,y/(xyyxy^2)`$, respectively.
###### Problem H.
Reconstruct not only a shape algebra, but a quantum $`G`$ itself from a braided tuple satisfying the conditions found in Problem G.
## Part II A conjectural flag tuple for the standard Drinfeld-Jimbo quantum groups
In this Part, we describe ingredients for a potential braided tuple, and we conjecture that these geometric data provide positive answers to Questions F and D for the standard quantum groups of Drinfeld and Jimbo. (The conjecture concerning Question F will be proved for $`\mathrm{SL}(n)`$ in Part III.)
Again, $`G`$ will denote a simply-connected simple complex group.
### 4. Recollections on $`\mathrm{U}_q^{\mathrm{DJ}}(𝔤)`$ and $`𝒪_q^{\mathrm{DJ}}(G)`$
Let $`𝔤`$ be the Lie algebra of $`G`$. Drinfel and Jimbo have defined (independently) a Hopf algebra $`\mathrm{U}_q^{\mathrm{DJ}}(𝔤)`$ that depends on a parameter $`q𝐂^{}`$ and that is a “quantum analogue” of the universal enveloping algebra $`\mathrm{U}(𝔤)`$ (in the sense that its comultiplication is no longer cocommutative). When $`q`$ is not a root of unity, finite-dimensional $`\mathrm{U}_q^{\mathrm{DJ}}(𝔤)`$-modules have been studied (independently) by Lusztig and by Rosso : in particular, discarding unwanted nontrivial one-dimensional modules, there still exists a family $`\{V_\lambda \lambda P^+\}`$ of $`\mathrm{U}_q(𝔤)`$-modules satisfying conditions (a) and (c) of Definition 1.1 (condition (c) follows e.g. from Theorem 4.12(b) of ).
Therefore, if $`𝒪_q^{\mathrm{DJ}}(G)`$ denotes the subspace of $`\mathrm{U}_q^{\mathrm{DJ}}(𝔤)^{}`$ spanned by the matrix coefficients of the modules $`V_\lambda `$, then $`𝒪_q^{\mathrm{DJ}}(G)`$ (for $`q`$ not a root of unity) is a quantum $`G`$ in the sense of Definition 1.1. (If $`G`$ is not assumed to be simply connected, then $`𝒪_q^{\mathrm{DJ}}(G)`$ may still be defined in this way, provided $`P^+`$ is replaced by the appropriate submonoid.) We call $`𝒪_q^{\mathrm{DJ}}(G)`$ the *standard* quantum $`G`$. When $`G=\mathrm{SL}(n)`$, $`\mathrm{SO}(n)`$, or $`\mathrm{Sp}(n)`$, a presentation of $`𝒪_q^{\mathrm{DJ}}(G)`$ by generators and relations has been given by Faddeev, Reshetikhin, and Takhtajan .
### 5. Recollections from
Choose a Borel subgroup $`BG`$ and a maximal torus $`TB`$, and let $`W:=\mathrm{N}_G(T)/T`$ be the associated Weyl group. Denote by $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^+`$ the root system and the set of positive roots, respectively. To each $`\alpha \mathrm{\Phi }`$ are associated a reflection $`s_\alpha W`$, a root group $`U_\alpha `$, and a copy $`L_\alpha =U_\alpha ,U_\alpha `$ of $`(\mathrm{P})\mathrm{SL}(2)`$ in $`G`$.
Recall the following construction from : an *orthocell* (of rank $`d`$) is a left coset in $`W`$ of the form
$$C=C(w;\alpha _1,\mathrm{},\alpha _d):=ws_{\alpha _1},\mathrm{},s_{\alpha _d},$$
where $`wW`$ and $`\alpha _1,\mathrm{},\alpha _d`$ are positive and pairwise orthogonal roots.
*Warning:* the $`\alpha _k`$ are *not* assumed to be *strongly* orthogonal, i.e. the sum of two of them may well be a root.
By orthogonality, the reflections $`s_{\alpha _1},\mathrm{},s_{\alpha _d}`$ pairwise commute. Therefore, the following notation makes sense, and we will use it frequently:
$$s_L:=\underset{kL}{}s_{\alpha _k},L\{1,\mathrm{},d\}.$$
(Note that the elements of $`C`$ are those of the form $`ws_L`$.)
Reordering the sequence $`\alpha _1,\mathrm{},\alpha _d`$ if necessary, assume that it is *nonincreasing*, in the sense that $`\alpha _k\alpha _k^{}`$ for all $`k<k^{}`$; then define
$$E(C):=\{\dot{w}g_1\mathrm{}g_dBg_kL_{\alpha _k}k\}G/B,$$
where $`\dot{w}\mathrm{N}_G(T)`$ is some representative of $`w`$. In , we show that $`E(C)`$ only depends on $`C`$ as a coset (and not on the choice of $`w`$ in $`C`$, nor of its representative $`\dot{w}`$, nor on the chosen nonincreasing ordering of the $`\alpha _k`$). Furthermore, we show that $`E(C)`$ is a $`T`$-stable subvariety of $`G/B`$, isomorphic to the product $`^1\times \mathrm{}\times ^1`$ of $`d`$ projective lines.
###### Remark 5.1.
Orthocells may also be defined in terms of *right* cosets: if we set
$$C(\alpha _1,\mathrm{},\alpha _d;w):=s_{\alpha _1},\mathrm{},s_{\alpha _d}w,$$
then $`C(\alpha _1,\mathrm{},\alpha _d;w)=C(w;w^1\alpha _1,\mathrm{},w^1\alpha _d)`$. Moreover, recall (see , end of §9.2.1) that for each $`wW`$ and each root $`\alpha `$, we have $`wU_\alpha w^1=U_{w\alpha }`$, and hence $`wL_\alpha w^1=L_{w\alpha }`$. It follows that for $`C=C(\alpha _1,\mathrm{},\alpha _d;w)`$, we have
$$E(C)=\{g_1\mathrm{}g_dwBg_kL_{\alpha _k}k\}.$$
### 6. Monogressive orthocells and the variety $`E^{\mathrm{DJ}}`$
Denote by $`<`$ the Bruhat order on $`W`$ and by $``$ the associated cover relation (i.e. $`ww^{}`$ if $`w<w^{}`$ and if no element of $`W`$ lies between $`w`$ and $`w^{}`$). Denote also by $`\mathrm{}(w)`$ the length of an element $`wW`$. Recall the following combinatorial characterization (see e.g. Sections 5.9 and 5.11 of ):
$$ww^{}\mathrm{}(w^{})=\mathrm{}(w)+1\text{ and }w^{}=ws\text{ for some reflection }s\text{.}$$
Assume that $`wC`$ has been chosen of minimal length.
###### Definition 6.1.
An orthocell $`C=C(w;\alpha _1,\mathrm{},\alpha _d)`$ will be called *monogressive* if
$$ws_Lws_Ls_{\alpha _k}L\{1,\mathrm{},d\},kL,$$
or, equivalently, if
$$\mathrm{}(ws_L)=\mathrm{}(w)+|L|L\{1,\mathrm{},d\}.$$
We then define the variety $`E^{\mathrm{DJ}}G/B`$ by
$$E^{\mathrm{DJ}}:=\underset{C\text{ monogressive}}{}E(C).$$
### 7. The automorphisms $`\sigma _1,\mathrm{},\sigma _{\mathrm{}}`$
Let $`\beta _1,\mathrm{},\beta _{\mathrm{}}`$ be the simple roots. Then the morphism
$$T(𝐂^{})^{\mathrm{}}:t(\beta _1(t),\mathrm{},\beta _{\mathrm{}}(t))$$
is surjective, so we may choose, for each $`1i\mathrm{}`$, an element $`t_iT`$ such that
$$\beta _j(t_i)=q^{(\varpi _i|\beta _j)}=\{\begin{array}{cc}q^{(\beta _j|\beta _j)/2}\hfill & \text{if }j=i,\hfill \\ 1\hfill & \text{if }ji\text{.}\hfill \end{array}$$
By (multiplicative) linearity, it then follows that $`\alpha (t_i)=q^{(\varpi _i|\alpha )}`$ for every root $`\alpha \mathrm{\Phi }`$.
Now let $`C=C(w;\alpha _1,\mathrm{},\alpha _d)`$ be a monogressive orthocell. Since $`E(C)`$ is $`T`$-stable in $`G/B`$, the automorphism
$$\sigma _{i,C}:E(C)E(C):gBwt_iw^1gB$$
is well defined, and it is independent of the choice of $`t_i`$ because the kernel of the above morphism $`T(𝐂^{})^{\mathrm{}}`$ is equal to the centre of $`G`$ (see e.g. \[36, Proposition 8.1.1\]).
###### Proposition 7.1.
For each $`i`$, the automorphisms $`\sigma _{i,C}`$ glue together to form a well defined automorphism $`\sigma _i`$ of $`E^{\mathrm{DJ}}`$.
We defer the proof to Appendix B.
### 8. The line bundles $`_1,\mathrm{},_{\mathrm{}}`$
Recall that for each $`\lambda P^+`$, the highest weight point in $`(V^\lambda )`$ is fixed by $`B`$, hence we get a well defined *Plücker map*
$$\mathrm{Pl}^\lambda :G/B(V^\lambda ).$$
Let $`\varpi _1,\mathrm{},\varpi _{\mathrm{}}`$ be the fundamental weights and write $`\mathrm{Pl}^i:=\mathrm{Pl}^{\varpi _i}`$ for each $`i`$. We then define the a line bundle $`_i`$ as the pullback of $`𝒪_{(V^i)}(1)`$ along $`\mathrm{Pl}^i`$, restricted to $`E^{\mathrm{DJ}}`$.
*Warning.* We may not define $`_\lambda `$ to be the pullback of $`𝒪_{(V^\lambda )}(1)`$ for all $`\lambda P^+`$: this would cause a conflict with the recursion rule (3.2).
### 9. Main conjectures and result
###### Conjecture 9.1 (Positive answer to Question F).
Assume that $`q𝐂^{}`$ is not a root of unity. The tuple $`T^{\mathrm{DJ}}=(E^{\mathrm{DJ}},\sigma _1,\mathrm{},\sigma _{\mathrm{}},_1,\mathrm{},_{\mathrm{}})`$ defined in Sections 68 is a braided tuple (see Definition 3.1), and the associated quadratic algebra $`M(T^{\mathrm{DJ}})`$ is the shape algebra of the standard quantum group $`𝒪_q^{\mathrm{DJ}}(G)`$.
###### Conjecture 9.2 (Positive answer to Question D).
Moreover, the same tuple $`T^{\mathrm{DJ}}`$ is the flag tuple associated to $`𝒪_q^{\mathrm{DJ}}(G)`$ (i.e. $`E^{\mathrm{DJ}}`$ parametrizes *all* flag modules of the shape algebra of $`𝒪_q^{\mathrm{DJ}}(G)`$).
###### Theorem 9.3.
Conjecture 9.1 is true for $`G=\mathrm{SL}(n)`$.
## Part III The standard quantum $`\mathrm{SL}(n)`$
In this Part, we will describe the objects of Sections 57 more explicitely when $`G=\mathrm{SL}(n)`$, and we prove Conjecture 9.1 in that case.
### 10. The varieties $`E(C)`$
From now on, it will be more convenient to view orthocells as *right* cosets (see Remark 5.1).
Let us first recall the usual realization of the flag variety $`\mathrm{SL}(n)/B`$, of the Plücker maps $`\mathrm{Pl}^i`$, and of the subgroups $`L_\alpha `$.
We let $`B\mathrm{SL}(n)`$ be the subgroup of all upper triangular matrices, i.e. the stabilizer of the flag
$$𝐂e_1𝐂e_1𝐂e_2\mathrm{}𝐂e_1\mathrm{}𝐂e_{n1},$$
where $`e_1,\mathrm{},e_n`$ denotes the canonical basis of $`𝐂^n`$. This identifies $`\mathrm{SL}(n)/B`$ with the set of all (full) flags in $`𝐂^n`$ (or in $`^{n1}:=(𝐂^n)`$).
We also let $`TB`$ be the subgroup of all diagonal matrices: the Weyl group $`W`$ then identifies with the symmetric group $`S_n`$, and the reflections correspond exactly to the transpositions.
For each $`1in1`$, recall that the fundamental representation $`V^i:=V^{\varpi _i}`$ is given by the exterior power $`\mathrm{\Lambda }^i𝐂^n`$, and that the map $`\mathrm{Pl}^i:\mathrm{SL}(n)/B(\mathrm{\Lambda }^i𝐂^n)`$ may be described as follows: given a flag $`F\mathrm{SL}(n)/B`$, choose a basis $`f_1,\mathrm{},f_i`$ of its component $`F_i`$ of dimension $`i`$, then send $`F`$ to the point $`f_1\mathrm{}f_i(\mathrm{\Lambda }^i𝐂^n)`$ (which is independent of the choice of the basis). Moreover, the elements
$$e_w^i:=e_{w(1)}\mathrm{}e_{w(i)},wS_n,$$
form a basis of $`\mathrm{\Lambda }^i𝐂^n`$ (up to obvious redundancies).
Let $`\alpha \mathrm{\Phi }^+`$ and write $`s_\alpha =(ab)`$, $`1a<bn`$. Then the subgroup $`L_\alpha \mathrm{SL}(n)`$ is the group $`\mathrm{SL}(2)`$ acting naturally on $`𝐂e_a𝐂e_b`$ and trivially on all other $`e_c`$. Clearly, if $`s_\alpha ,s_\beta `$ commute (i.e. if $`\alpha ,\beta `$ are orthogonal), then so do $`L_\alpha `$ and $`L_\beta `$. (This is not true for arbitrary $`G`$.)
Now fix an orthocell $`C=C(\alpha _1,\mathrm{},\alpha _d;w)`$ and let us describe the variety $`E(C)`$, or rather, its images under the maps $`\mathrm{Pl}^i`$, $`1in1`$.
###### Remark 10.1.
For each $`1kd`$, the following conditions are equivalent:
* $`s_{\alpha _k}w\varpi _i=w\varpi _i`$,
* the transposition $`s_{\alpha _k}`$ leaves the set $`\{w(1),\mathrm{},w(i)\}`$ invariant,
* $`e_{s_{\alpha _k}w}^i=\pm e_w^i`$.
Number the $`\alpha _k`$ in such a way that for some $`1ad`$, the above conditions hold for $`1ka`$ and do not hold for $`a+1kd`$. For each $`k`$, write $`s_{\alpha _k}=(a_kb_k)`$, $`a_k<b_k`$, and pick an element $`g_kL_{\alpha _k}`$ acting as $`\left(\begin{array}{cc}x_k& \\ y_k& \end{array}\right)`$ on $`𝐂e_{a_k}+𝐂e_{b_k}`$ (and trivially on the other $`e_c`$). For any subset $`L\{1,\mathrm{},a\}`$, write
$$\overline{L}:=\{1,\mathrm{},a\}L,x_L:=\underset{kL}{}x_k,y_L:=\underset{kL}{}y_k.$$
The above description of the map $`\mathrm{Pl}^i`$ and of the subgroups $`L_{\alpha _k}`$ now imply that
(10.1)
$$\mathrm{Pl}^i(g_1\mathrm{}g_dwB)=\underset{L\{1,\mathrm{},a\}}{}x_{\overline{L}}y_Le_{s_Lw}^i(\mathrm{\Lambda }^i𝐂^n).$$
###### Remark 10.2.
The variety $`E(C)`$ being a product of $`d`$ projective lines, we may view $`(x_1:y_1),\mathrm{},(x_d:y_d)`$ as homogeneous coordinates on these lines.
A more geometric description of the varieties $`E(C)`$ (not needed here) can be found in \[33, Example 5.1\].
### 11. Monogressivity
Let us first recall a more explicit description of the Bruhat cover relation in $`S_n`$. Write a permutation $`wS_n`$ as an array $`[w(1)\mathrm{}w(n)]`$, and write e.g. $`w=[abc]`$ to signify that in the array $`w`$, $`a`$ appears to the left of $`b`$ and $`b`$ appears to the left of $`c`$.
If $`sS_n`$ is a transposition, say $`s=(ab)`$ with $`a<b`$, then $`wsw`$ if and only if (i) $`w=[ab]`$ and (ii) whenever $`w=[acb]`$, $`c`$ is outside of the (numerical) interval $`[a,b]`$. For example, if $`n=7`$ and $`s=(\mathrm{4\hspace{0.17em}6})`$, then we have $`[3472651][3672451]`$, but $`[7415623]\overline{)}[7615423]`$ because the subarray $`[4156]`$ contains $`5`$.
Now let $`C=C(\alpha _1,\mathrm{},\alpha _d;w)`$ be an orthocell, and write again $`s_{\alpha _k}=(a_kb_k)`$, $`a_k<b_k`$, for all $`k`$. The above description of the Bruhat cover relation shows the following
###### Criterion 11.1.
With the above notation, the orthocell $`C`$ is monogressive if and only if the following conditions hold for all $`k`$:
* $`w=[a_kb_k]`$,
* whenever $`w=[a_kcb_k]`$, $`c`$ is outside of the interval $`[a_k,b_k]`$,
* whenever $`w=[a_ka_k^{}b_k]`$ for some $`k^{}k`$, both $`a_k^{}`$ and $`b_k^{}`$ are outside of the interval $`[a_k,b_k]`$, and similarly whenever $`w=[a_kb_k^{}b_k]`$.
###### Example 11.2 ($`n=4`$).
There are fifty-eight monogressive orthocells of rank $`1`$, viz. those of one of the following forms:
* $`\{[ijkl],[jikl]\}`$, $`\{[kijl],[kjil]\}`$, or $`\{[klij],[klji]\}`$ ($`i<j`$),
* $`\{[ikjl],[jkil]\}`$ or $`\{[likj],[ljki]\}`$ ($`i<j`$; $`k[i,j]`$),
* $`\{[iklj],[jkli]\}`$ ($`i<j`$; $`k,l[i,j]`$).
There are eleven monogressive orthocells $`C(\alpha _1,\alpha _2;w)`$ of rank $`2`$, given by
* $`s_{\alpha _1}=(\mathrm{1\hspace{0.17em}2})`$, $`s_{\alpha _2}=(\mathrm{3\hspace{0.17em}4})`$, $`w=[1234]`$, $`[3412]`$, $`[1324]`$, $`[3142]`$, $`[1342]`$, or $`[3124]`$;
* $`s_{\alpha _1}=(\mathrm{1\hspace{0.17em}3})`$, $`s_{\alpha _2}=(\mathrm{2\hspace{0.17em}4})`$, $`w=[1324]`$ or $`[2413]`$;
* $`s_{\alpha _1}=(\mathrm{1\hspace{0.17em}4})`$, $`s_{\alpha _2}=(\mathrm{2\hspace{0.17em}3})`$, $`w=[1423]`$, $`[2314]`$, or $`[2143]`$.
(Pictures for the corresponding varieties $`E(C)`$ may be found in \[33, Example 5.1\].)
###### Example 11.3 ($`n=3`$).
View $`\mathrm{SL}(3)/B`$ as the set of flags $`(p,l)`$ in $`^2=(𝐂^3)`$. Consider $`e_1,e_2,e_3`$ as points in $`^2`$ and let $`e_{ab}^2`$ be the line through $`e_a`$ and $`e_b`$. Then $`E^{\mathrm{DJ}}`$ is the union of the following eight curves in $`\mathrm{SL}(3)/B`$:
$`\{(e_a,l)e_al\},a=1,2,3,`$
$`\{(p,e_{ab})pe_{ab}\},ab=12,13,23,`$
$`\{(p,l)e_{12}ple_3\},`$
$`\{(p,l)e_{23}ple_1\}.`$
See Figure 1.
### 12. The automorphisms $`\sigma _i`$
Denote again the simple roots by $`\beta _1,\mathrm{},\beta _{n1}`$. If $`t=\mathrm{diag}(x_1,\mathrm{},x_n)T`$ (with $`_jx_j=1`$), then recall that $`\beta _i(t)=x_ix_{i+1}^1`$. Therefore, $`t_iT`$ is equal, up to a factor, to the matrix $`\mathrm{diag}(1,\mathrm{},1,q,\mathrm{},q)`$ ($`i`$ times $`1`$ and $`ni`$ times $`q`$).
If $`C=C(\alpha _1,\mathrm{},\alpha _d;w)`$, with $`s_{\alpha _k}=(a_kb_k)`$ as before, then the action of the associated automorphism $`\sigma _i:gBwt_iw^1gB`$ on $`E(C)`$ may be described more explicitly using the homogeneous coordinates $`(x_1:y_1),\mathrm{},(x_d:y_d)`$ of Remark 10.2:
$$\sigma _i:(x_k:y_k)\{\begin{array}{cc}(x_k:qy_k)\hfill & \text{if }s_{\alpha _k}w\varpi _iw\varpi _i\text{,}\hfill \\ (x_k:y_k)\hfill & \text{if }s_{\alpha _k}w\varpi _i=w\varpi _i\text{.}\hfill \end{array}$$
###### Example 12.1 ($`n=3`$).
On each of the eight components of $`E^{\mathrm{DJ}}`$ (see Example 11.3), $`\sigma _1,\sigma _2`$ act as homotheties (viewing the two $`T`$-stable points on this component as $`0`$ and $`\mathrm{}`$). The ratios for $`\sigma _1`$ are, respectively, $`1,1,1,q,q,q,q,q`$, and those for $`\sigma _2`$ are $`q,q,q,1,1,1,q,q`$.
### 13. Proof of Conjecture 9.1 for $`G=\mathrm{SL}(n)`$
For each $`1in1`$, consider the vector space $`V^i:=\mathrm{\Lambda }^i𝐂^n`$. On one hand, $`\mathrm{SL}(n)`$ acts on it naturally, and the corresponding map $`\mathrm{SL}(n)/B(V^i)`$ induces a line bundle $`_i`$ on $`E^{\mathrm{DJ}}\mathrm{SL}(n)/B`$ (see Section 8). On the other hand, we will make $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$ act on $`V^i`$ (see below, before Lemma 13.4), turning $`V^i`$ into the simple $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$-module of highest weight $`\varpi _i`$.
Both the algebra $`M(T^{\mathrm{DJ}})`$ and the shape algebra $`M^{\mathrm{DJ}}`$ thus become quotients of the tensor algebra $`T(V_1\mathrm{}V_{n1})`$ (where $`V_i:=(V^i)^{}`$). Note also that $`M(T^{\mathrm{DJ}})`$ is quadratic by definition, and $`M^{\mathrm{DJ}}`$ is quadratic by . So we need to show that relations of types (I) and (II) (see end of Section 1) agree for both algebras (and, of course, that the tuple $`T^{\mathrm{DJ}}`$ is braided in the first place).
We will break down the proof into several lemmas.
###### Definition 13.1.
Let $`1i,jn1`$. An orthocell $`C=C(\alpha _1,\mathrm{},\alpha _d;w)`$ will be called *$`ij`$-effective* if, for every $`1kd`$, we have both $`s_{\alpha _k}w\varpi _iw\varpi _i`$ and $`s_{\alpha _k}w\varpi _jw\varpi _j`$. In this case, we define the following element of $`V^iV^j`$:
$$e_C^{ij}:=\underset{L\{1,\mathrm{},d\}}{}q^{|L|}e_{s_{\overline{L}}w}^ie_{s_Lw}^j,$$
where, as before, $`s_L:=_{kL}s_{\alpha _k}`$ and $`\overline{L}:=\{1,\mathrm{},d\}L`$.
We denote by $`V^{ij}V^iV^j`$ the linear span of the image of the map $`\mathrm{Pl}^i(\mathrm{Pl}^j)^{\sigma _i}`$ (see (3.4)).
###### Lemma 13.2.
The subspace $`V^{ij}`$ is linearly spanned by the $`e_C^{ij}`$ for $`C`$ monogressive and $`ij`$-effective.
###### Proof.
First, let $`C=C(\alpha _1,\mathrm{},\alpha _d;w)`$ be monogressive and $`ij`$-effective. If a point $`pE(C)`$ has homogeneous coordinates $`(x_1:y_1),\mathrm{},(x_d:y_d)`$ (see Remark 10.2), then by $`ij`$-effectiveness, the coordinates of $`\sigma _i(p)`$ are $`(x_1:qy_1),\mathrm{},(x_d:qy_d)`$. Using (10.1), we therefore see that $`\mathrm{Pl}^i(\mathrm{Pl}^j)^{\sigma _i}`$ sends $`p`$ to the following point in $`(V^iV^j)`$:
$$\underset{L,M\{1,\mathrm{},d\}}{}q^{|M|}x_{\overline{L}}x_{\overline{M}}y_Ly_Me_{s_Lw}^ie_{s_Mw}^j.$$
Using the change of “variables” $`I:=L\mathrm{}M`$ (symmetric difference), $`J:=LM`$, and $`N:=ML=MJI`$, this expression may be rewritten as
$`{\displaystyle \underset{\begin{array}{c}I,J\{1,\mathrm{},d\}\\ IJ=\mathrm{}\end{array}}{}}q^{|J|}(x_{\overline{IJ}})^2x_Iy_I(y_J)^2\left({\displaystyle \underset{NI}{}}q^{|N|}e_{s_{\overline{N}}s_Jw}^ie_{s_Ns_Jw}^j\right)`$
$`={\displaystyle \underset{\begin{array}{c}I,J\{1,\mathrm{},d\}\\ IJ=\mathrm{}\end{array}}{}}q^{|J|}(x_{\overline{IJ}})^2x_Iy_I(y_J)^2e_{C(\alpha _I;s_Jw)}^{ij},`$
where $`\alpha _I`$ is shorthand for the set $`\{\alpha _kkI\}`$. By induction over $`d`$, we may assume that $`e_C^{}^{ij}V^{ij}`$ for all monogressive $`ij`$-effective orthocells $`C^{}`$ of rank smaller than $`d`$ (the case $`d=0`$ being trivial). Since the above sum is in $`V^{ij}`$ by definition, the only remaining term, namely $`e_C^{ij}`$, is in $`V^{ij}`$ as well.
We still need to show that the image of a point $`pE(C)`$ is in the span of the $`e_C^{}^{ij}`$ (for $`C^{}`$ monogressive and $`ij`$-effective) even if $`C`$ is not $`ij`$-effective (but still monogressive). Reordering the $`\alpha _k`$ if necessary, we may assume that, for some $`1abcd`$, they satisfy
$`s_{\alpha _k}w\varpi _i`$ $`w\varpi _i,`$ $`s_{\alpha _k}w\varpi _j`$ $`w\varpi _j`$ if $`1ka`$;
$`s_{\alpha _k}w\varpi _i`$ $`w\varpi _i,`$ $`s_{\alpha _k}w\varpi _j`$ $`=w\varpi _j`$ if $`a+1kb`$;
$`s_{\alpha _k}w\varpi _i`$ $`=w\varpi _i,`$ $`s_{\alpha _k}w\varpi _j`$ $`w\varpi _j`$ if $`b+1kc`$;
$`s_{\alpha _k}w\varpi _i`$ $`=w\varpi _i,`$ $`s_{\alpha _k}w\varpi _j`$ $`=w\varpi _j`$ if $`c+1kd`$.
Let again $`(x_1:y_1),\mathrm{},(x_d:y_d)`$ be homogeneous coordinates for a point $`pE(C)`$. This time, the coordinates for $`\sigma _i(p)`$ are obtained by multiplying $`y_k`$ by $`q`$ only for $`1kb`$. Furthermore, we have
$$\mathrm{Pl}^i(p)=\underset{\begin{array}{c}L\{1,\mathrm{},a\}\\ L^{}\{a+1,\mathrm{},b\}\end{array}}{}x_{\overline{L}}x_{\overline{L}^{}}y_Ly_L^{}e_{s_Ls_L^{}w}^i,$$
and a similar expression for $`\mathrm{Pl}^j(p)`$, with $`\{a+1,\mathrm{},b\}`$ replaced by $`\{b+1,\mathrm{},c\}`$. A computation similar to the one above shows that $`\mathrm{Pl}^i(\mathrm{Pl}^j)^{\sigma _i}`$ now sends $`p`$ to
(13.1)
$$\underset{\begin{array}{c}I,J\{1,\mathrm{},a\}\\ IJ=\mathrm{}\\ L^{}\{a+1,\mathrm{},b\}\\ M^{}\{b+1,\mathrm{},c\}\end{array}}{}q^{|J|}(x_{\overline{IJ}})^2x_Ix_{\overline{L}^{}}x_{\overline{M}^{}}y_I(y_J)^2y_L^{}y_M^{}e_{C(\alpha _I;s_Js_L^{}s_M^{}w)}^{ij}$$
(where we have used the fact that $`e_w^i=e_{s_M^{}w}^i`$ and $`e_w^j=e_{s_L^{}w}^j`$). ∎
###### Lemma 13.3.
Let $`1i,jn1`$, with, say, $`i<j`$. Then
$$dimV^{ij}D_{n;i,j}:=\left(\genfrac{}{}{0pt}{}{n}{i}\right)\left(\genfrac{}{}{0pt}{}{n}{j}\right)\left(\genfrac{}{}{0pt}{}{n}{i1}\right)\left(\genfrac{}{}{0pt}{}{n}{j+1}\right).$$
###### Proof.
Consider a monogressive $`ij`$-effective cell $`C=C(\alpha _1,\mathrm{},\alpha _d;w)`$. For each $`1kd`$, write again $`s_{\alpha _k}=(a_kb_k)`$, $`a_k<b_k`$ (hence $`w=[a_kb_k]`$ by monogressivity). Reorder the $`\alpha _k`$ in such a way that $`b_1<\mathrm{}<b_d`$. By $`ij`$-effectiveness, each $`a_k`$ must appear in the subarray $`[w(1)\mathrm{}w(i)]`$, and each $`b_k`$ in the subarray $`[w(j+1)\mathrm{}w(n)]`$. Now let $`S_{ij}:=S_i\times S_{ji}\times S_{n(i+j)}S_n`$, and note that for each $`\pi S_{ij}`$, replacing $`w`$ by $`w\pi `$ in $`C`$ leaves $`e_C^{ij}`$ invariant up to a sign. Choosing $`\pi `$ appropriately, we may assume that $`w`$ takes the following form:
$$[w(1)\mathrm{}w(id)a_1\mathrm{}a_dw(i+1)\mathrm{}w(j)b_d\mathrm{}b_1w(j+d+1)\mathrm{}w(n)],$$
with, say, the following orderings:
$$w(1)>\mathrm{}>w(id),w(i+1)>\mathrm{}>w(j),w(j+d+1)>\mathrm{}>w(n).$$
This rearrangement does not affect the monogressivity of $`C`$ (nor, for that matter, its $`ij`$-effectiveness). Indeed, the only nonobvious point here is the relative ordering of the $`a_k`$ and the $`b_k`$: by monogressivity, we have $`ws_{\alpha _k}ws_{\alpha _k}s_{\alpha _k^{}}w`$ and $`ws_{\alpha _k^{}}ws_{\alpha _k}s_{\alpha _k^{}}w`$, so if $`k<k^{}`$, then, whatever the order in which $`a_k,a_k^{},b_k,b_k^{}`$ appear in the original array $`[w(1)\mathrm{}w(n)]`$, we must have either $`a_k<b_k<a_k^{}<b_k^{}`$, or $`a_k^{}<a_k<b_k<b_k^{}`$. In both cases, $`a_k^{}`$ and $`b_k^{}`$ are outside of the (numerical) interval $`[a_k,b_k]`$, so they may indeed appear between $`a_k`$ and $`b_k`$ in the new array without affecting monogressivity.
An orthocell thus modified will be called *$`ij`$-normal*. The proof will be finished if we show that there are $`D_{n;i,j}`$ $`ij`$-normal orthocells in $`S_n`$. Since we have the recursion rule
$$D_{n+1;i,j}=D_{n;i1,j1}+D_{n;i1,j}+D_{n;i,j1}+D_{n;i,j},$$
it is enough to show that the number of $`ij`$-normal orthocells in $`S_{n+1}`$ satisfies the same recursion rule. If $`C`$ is such an orthocell, there are two possibilities.
* Either each $`s_{\alpha _k}`$ fixes $`n+1`$. Removing $`n+1`$ from the array $`[w(1)\mathrm{}w(n+1)]`$, we then obtain an orthocell in $`S_n`$, which is $`(i1)(j1)`$-normal, $`i(j1)`$-normal, or $`ij`$-normal, according to the position of $`n+1`$ in the array, relative to $`w(i)`$ and $`w(j)`$.
* Or some $`s_{\alpha _k}`$ involves $`n+1`$: necessarily, $`k=d`$ and $`b_d=n+1`$. Removing again $`n+1`$ from the array, and discarding $`\alpha _d`$ from $`C`$, we then obtain an $`(i1)j`$-normal orthocell in $`S_n`$.
Clearly, this procedure may be reversed, starting from a normal orthocell in $`S_n`$ and inserting $`n+1`$ at all possible places in the corresponding array. Hence the desired recursion rule. ∎
Now let us recall a presentation for the quantized enveloping algebra $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$, following e.g. : it is generated by $`4(n1)`$ elements $`K_\beta `$, $`K_\beta ^1`$, $`X_\beta `$, $`Y_\beta `$ ($`\beta `$ a simple root), subject to the commutation relations
$`K_\beta K_\beta ^1=1`$ $`=K_\beta ^1K_\beta ,K_\beta K_\gamma =K_\gamma K_\beta ,`$
$`K_\beta X_\gamma K_\beta ^1`$ $`=q^{(\beta |\gamma )}X_\gamma ,`$
$`K_\beta Y_\gamma K_\beta ^1`$ $`=q^{(\beta |\gamma )}Y_\gamma ,`$
$`X_\beta Y_\gamma Y_\gamma X_\beta `$ $`=\delta _{\beta \gamma }{\displaystyle \frac{K_\beta K_\beta ^1}{qq^1}},`$
as well as the quantized Serre relations
$`X_\beta ^2X_\gamma (q+q^1)X_\beta X_\gamma X_\beta +X_\gamma X_\beta ^2`$ $`=0\text{if }\beta ,\gamma \text{ adjacent,}`$
$`X_\beta X_\gamma X_\gamma X_\beta `$ $`=0\text{if }\beta ,\gamma \text{ not adjacent,}`$
$`Y_\beta ^2Y_\gamma (q+q^1)Y_\beta Y_\gamma Y_\beta +Y_\gamma Y_\beta ^2`$ $`=0\text{if }\beta ,\gamma \text{ adjacent,}`$
$`Y_\beta Y_\gamma Y_\gamma Y_\beta `$ $`=0\text{if }\beta ,\gamma \text{ not adjacent.}`$
Moreover, $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$ is a Hopf algebra whose comultiplication is given on the generators by
(13.2) $`\mathrm{\Delta }K_\beta ^{\pm 1}`$ $`=K_\beta ^{\pm 1}K_\beta ^{\pm 1},`$
$`\mathrm{\Delta }X_\beta `$ $`=X_\beta 1+K_\beta X_\beta ,`$
$`\mathrm{\Delta }Y_\beta `$ $`=Y_\beta K_\beta ^1+1Y_\beta .`$
We then define a $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$-module structure on $`V^i`$ as follows. For every $`wW`$ and every simple root $`\beta `$, we set
$$K_\beta e_w^i=q^{(w\varpi _i|\beta )}e_w^i,K_\beta ^1e_w^i=q^{(w\varpi _i|\beta )}e_w^i,$$
and
$`X_\beta e_w^i`$ $`=0,`$ $`Y_\beta e_w^i`$ $`=e_{s_\beta w}^i`$ if $`(w\varpi _i|\beta )=1`$;
$`X_\beta e_w^i`$ $`=0,`$ $`Y_\beta e_w^i`$ $`=0`$ if $`(w\varpi _i|\beta )=0`$;
$`X_\beta e_w^i`$ $`=e_{s_\beta w}^i,`$ $`Y_\beta e_w^i`$ $`=0`$ if $`(w\varpi _i|\beta )=1`$.
(These are the only possible values for $`(w\varpi _i|\beta )`$, because $`\varpi _i`$ is minuscule.) It is straightforward to check that this module structure is well defined, and that it is the simple $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$-module of highest weight $`\varpi _i`$.
###### Lemma 13.4.
The subspace $`V^{ij}`$ is a $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$-submodule of $`V^iV^j`$.
By Lemma 13.2, the statement means that the action of a generator of $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$ on a vector $`e_C^{ij}`$, $`C`$ monogressive and $`ij`$-effective, must again be a linear combination of such vectors. We postpone these rather tedious computations to Appendix C.
###### Corollary 13.5.
The subspace $`V^{ij}`$ is equal to the (unique) $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$-submodule of $`V^iV^j`$ of highest weight $`\varpi _i+\varpi _j`$, and the $`e_C^{ij}`$ (for $`C`$ monogressive and $`ij`$-effective) are linearly independent.
###### Proof.
The vector $`e_{C(\text{};1)}^{ij}=e_1^ie_1^j`$ is a highest weight vector, of weight $`\varpi _i+\varpi _j`$. Now apply Lemmas 13.2 and 13.3, noting that the dimension of the simple module of highest weight $`\varpi _i+\varpi _j`$ is precisely $`D_{n;i,j}`$. ∎
###### Lemma 13.6.
The linear map $`R^{ji}:V^{ji}V^{ij}`$ defined by
(13.3)
$$R^{ji}(e_C^{ji})=e_C^{ij}\text{for all monogressive }ij\text{-effective }C$$
is an isomorphism of $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$-modules.
###### Proof.
This is immediate from the action of the generators of $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$ on the basis elements of $`V^{ij}`$ and $`V^{ji}`$, as described in Appendix C: the formulas obtained there are symmetric in $`i`$ and $`j`$. ∎
Extend $`R^{ji}`$ to an isomorphism $`V^jV^i\stackrel{}{}V^iV^j`$ of $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$-modules (in an arbitrary way).
###### Lemma 13.7.
The maps $`R^{ji}`$ induce isomorphisms $`R_{ij}:_i_j^{\sigma _i}\stackrel{}{}_j_i^{\sigma _j}`$ of line bundles over $`E^{\mathrm{DJ}}`$, and the latter satisfy (3.1) for all $`i,j,k`$.
###### Proof.
The first statement amounts to the commutativity of the diagram (3.5), which immediately follows from (13.1) and (13.3).
For the second statement, consider the composite map
(13.4)
$$\begin{array}{c}E^{\mathrm{DJ}}\stackrel{\text{diag.}}{}E^{\mathrm{DJ}}\times E^{\mathrm{DJ}}\times E^{\mathrm{DJ}}\stackrel{\mathrm{id}\times \sigma _i\times \sigma _i\sigma _j}{}E^{\mathrm{DJ}}\times E^{\mathrm{DJ}}\times E^{\mathrm{DJ}}\hfill \\ \hfill \stackrel{\mathrm{Pl}^i\times \mathrm{Pl}^j\times \mathrm{Pl}^k}{}(V^i)\times (V^j)\times (V^k)\stackrel{\text{Segre}}{}(V^iV^jV^k)\end{array}$$
corresponding to the line bundle $`_i_j^{\sigma _i}_k^{\sigma _i\sigma _j}`$, and denote by $`V^{ijk}V^iV^jV^k`$ the linear span of the image of this map.
Claim A. *The subspace $`V^{ijk}`$ is contained in the unique simple $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$-submodule $`W^{ijk}`$ of $`V^iV^jV^k`$ of highest weight $`\varpi _i+\varpi _j+\varpi _k`$.*
Indeed, let $`K_{ij}`$ be the kernel of a projection $`V_iV_jV_{\varpi _i+\varpi _j}`$, and define similarly $`K_{jk}`$, $`K_{ijk}`$. Since $`M^{\mathrm{DJ}}`$ is quadratic (cf. ), we have $`K_{ijk}=K_{ij}V_k+V_iK_{jk}`$, so dually, $`W^{ijk}=V^{ij}V^kV^iV^{jk}`$, and $`V^{ijk}`$ is clearly contained in the right hand side. This shows Claim A.
The proof will be finished if we show the following Claim (from which (3.1) follows):
Claim B. *Consider the maps $`(R^{ji}\mathrm{id})(\mathrm{id}R^{ki})(R^{kj}\mathrm{id})`$ and $`(\mathrm{id}R^{kj})(R^{ki}\mathrm{id})(\mathrm{id}R^{ji})`$ from $`V^kV^jV^i`$ to $`V^iV^jV^k`$. Their restrictions to $`V^{kji}`$ agree.*
By Claim A, it will be enough to show that the restrictions to $`W^{kji}`$ agree. Since both maps are morphisms between the simple $`\mathrm{U}_q^{\mathrm{DJ}}(𝔰𝔩(n))`$-modules $`W^{kji}`$ and $`W^{ijk}`$, they must be equal up to a constant. But they both send the (highest weight) vector $`e_1^ke_1^je_1^i`$ to $`e_1^ie_1^je_1^k`$, so this constant is equal to $`1`$. This shows Claim B. ∎
It now follows from Lemma 13.7 that the tuple $`T^{\mathrm{DJ}}=(E^{\mathrm{DJ}},\sigma _1,\mathrm{},\sigma _{\mathrm{}},_1,\mathrm{},_{\mathrm{}})`$ is braided. It also follows from Corollary 13.5 and Lemma 13.6 that the quadratic algebras $`M(T^{\mathrm{DJ}})`$ and $`M^{\mathrm{DJ}}`$ agree (as quotients of $`T(V_1\mathrm{}V_{n1})`$). Conjecture 9.1 is thus proved for $`G=\mathrm{SL}(n)`$.
### Appendix A Proof of Proposition 1.3
Assume that $`M_A`$ is associative. By Definition 1.1(c), there exists an $`A`$-isomorphism $`R_{ij}:V_iV_jV_jV_i`$. Rescaling $`R_{ij}`$ if necessary, we may assume that the diagram
(A.1)
commutes. Now consider the following diagram:
where we have omitted all tensor product symbols and written $`V_{i+j}`$ instead of $`V_{\varpi _i+\varpi _j}`$, etc. (The arrows are the obvious ones, coming either from the multiplication $`m`$ or from the $`R_{ij}`$.) All diamonds commute by associativity, and all triangles commute, being instances of (A.1). Moreover, each object in the diagram contains a unique copy of $`V_{\varpi _i+\varpi _j+\varpi _k}`$, and when all arrows are restricted to these subcomodules, they become isomorphisms. Therefore, the outer rim commutes, i.e. (1.1) holds.
Conversely, assume that (1.1) holds for all $`i>j>k`$. We first extend the definition of the $`R_{ij}`$ by setting $`R_{ii}:=\mathrm{id}`$ for all $`i`$ and $`R_{ji}:=R_{ij}^1`$ for all $`i>j`$. Relation (1.1) then holds for all $`i,j,k`$.
We will realize $`M_A`$ as a quotient of the tensor algebra $`\mathrm{T}(V_1\mathrm{}V_{\mathrm{}})`$. Let $`\mathrm{\Gamma }`$ be the free monoid on $`\{1,\mathrm{},\mathrm{}\}`$. For every $`I=i_1\mathrm{}i_r\mathrm{\Gamma }`$, let
$$\varpi _I:=\varpi _{i_1}+\mathrm{}+\varpi _{i_r},V^I:=V_{i_1}\mathrm{}V_{i_r},$$
so $`\mathrm{T}(V_1\mathrm{}V_{\mathrm{}})=_{I\mathrm{\Gamma }}V_I`$. By Definition 1.1(c), $`V^I`$ contains a unique copy of $`V_{\varpi _I}`$; let $`K_IV^I`$ be its unique invariant supplement. Then the direct sum $`K=_{I\mathrm{\Gamma }}K_I`$ is a (two-sided) ideal in $`\mathrm{T}(V_1\mathrm{}V_{\mathrm{}})`$ (again by Definition 1.1(c)), so we get a quotient algebra
$$\mathrm{T}(V_1\mathrm{}V_{\mathrm{}})/K=:\underset{I\mathrm{\Gamma }}{}V_I.$$
We still need to identify $`V_I`$ with $`V_J`$ whenever $`\varpi _I=\varpi _J`$. The argument will be fairly standard: use the $`R_{ij}`$ to exchange generators from different $`V_i`$’s, and check that this is consistent, using (1.1). More explicitly, denote by $`S_r`$ the symmetric group. If $`\pi S_r`$ and $`I=i_1\mathrm{}i_r\mathrm{\Gamma }`$, define $`\pi I:=i_{\pi (1)}\mathrm{}i_{\pi (r)}`$. Denote the usual generators of $`S_r`$ by $`s_j:=(j,j+1)`$, $`1jr1`$, and define an $`A`$-isomorphism $`R_{I,j}:V^IV^{s_jI}`$ by $`R_{i_ji_{j+1}}`$ on $`V_{i_j}V_{i_{j+1}}`$ and by $`\mathrm{id}`$ on all other $`V_{i_k}`$. If $`\pi S_r`$ decomposes as $`\pi =s_{j_1}\mathrm{}s_{j_t}`$ (*not* necessarily in a reduced way), define $`R_{I,\pi }:V^IV^{\pi I}`$ by $`R_{I,\pi }:=R_{I,j_1}\mathrm{}R_{I,j_t}`$. Since this is an $`A`$-isomorphism, it restricts to $`R_{I,\pi }:V_IV_{\pi I}`$. Note that this restriction does not depend on the chosen decomposition of $`\pi `$, thanks to $`R_{ii}=\mathrm{id}`$, to $`R_{ij}R_{ji}=\mathrm{id}`$, and to (1.1).
Now consider $`\mathrm{T}(V_1\mathrm{}V_{\mathrm{}})/K`$ as a $`P^+`$-graded algebra, the term of degree $`\lambda P^+`$ being $`U_\lambda :=_{I:\varpi _I=\lambda }V_I`$. Let $`S_\lambda U_\lambda `$ be the span of all elements $`xR_{I,\pi }(x)`$, $`xV_I`$, where $`I`$ runs over all elements of $`\mathrm{\Gamma }`$ such that $`\varpi _I=\lambda `$. Then $`U_\lambda /S_\lambda `$ consists of just one copy of $`V_\lambda `$: indeed, on one hand, $`\varpi _I=\varpi _J`$ if and only if $`J=\pi I`$ for some $`\pi S_r`$, and on the other hand, the construction of the $`R_{I,\pi }`$ implies that
$$R_{I,\pi \pi ^{}}=R_{\pi ^{}I,\pi }R_{I,\pi ^{}}.$$
Moreover, the direct sum $`S:=_{\lambda P^+}S_\lambda `$ is an ideal in $`\mathrm{T}(V_1\mathrm{}V_{\mathrm{}})/K`$, so the corresponding quotient yields the desired associative realization of the shape algebra $`M_A`$.
### Appendix B Proof of Proposition 7.1
We need to show that for any monogressive orthocells $`C_1,C_2`$, the automorphisms $`\sigma _{i,C_1}`$ and $`\sigma _{i,C_2}`$ agree on $`E(C_1)E(C_2)`$. Since this intersection is clearly $`T`$-stable and closed, it is a union of $`T`$-orbit closures in each of $`E(C_1)`$ and $`E(C_2)`$, hence \[33, Corollary 6.2\] a union of $`E(C^{})`$ with $`C^{}`$ a common subcell of $`C_1`$ and $`C_2`$.
It is therefore enough to show that the $`\sigma _{i,C}`$ are compatible with restriction to subcells. Consider a monogressive orthocell $`C=C(w;\alpha _1,\mathrm{},\alpha _d)`$ and a subcell, say, $`C^{}=C(ws_L;\alpha _1^{},\mathrm{},\alpha _e^{})`$, with $`L\{1,\mathrm{},d\}`$, $`\{\alpha _1^{},\mathrm{},\alpha _e^{}\}\{\alpha _1,\mathrm{},\alpha _d\}`$, and $`\alpha _k\{\alpha _1^{},\mathrm{},\alpha _e^{}\}`$ for all $`kL`$ (so that $`ws_L`$ is again of minimal length in $`C^{}`$). For each $`\alpha \mathrm{\Phi }`$, fix an isomorphism $`u_\alpha :(𝐂,+)U_\alpha `$ such that $`tu_\alpha (z)t^1=u_\alpha (\alpha (t)z)`$ for all $`tT`$ and all $`z𝐂`$ \[36, Proposition 8.1.1(i)\]. Then the set of all $`\dot{w}\dot{s}_Lu_{\alpha _1^{}}(z_1)\mathrm{}u_{\alpha _e^{}}(z_e)B`$, $`(z_1,\mathrm{},z_e)𝐂^e`$, is an open dense subset of $`E(C^{})`$ (cf. , proof of Theorem 4.1), and the action of $`\sigma _{i,C^{}}`$ on such an element is given by
$`ws_Lt_is_L^1w^1\dot{w}\dot{s}_Lu_{\alpha _1}(z_1)\mathrm{}u_{\alpha _d}(z_d)B`$
$`=\dot{w}\dot{s}_Lt_iu_{\alpha _1}(z_1)\mathrm{}u_{\alpha _d}(z_d)B`$
$`=\dot{w}\dot{s}_Lu_{\alpha _1}(\alpha _1(t_i)^1z_1)\mathrm{}u_{\alpha _d}(\alpha _d(t_i)^1z_d)B,`$
whereas the action of $`\sigma _{i,C}`$ (i.e. multiplication by $`wt_iw^1`$ instead of $`ws_Lt_is_L^1w^1`$) is given by the same expression, with $`t_i`$ replaced by $`s_L^1t_is_L`$. But since $`s_L`$ is a product of reflections w.r.t. roots orthogonal to each $`\alpha _k^{}`$, we have $`\alpha _k^{}(s_L^1ts_L)=(s_L\alpha _k^{})(t)=\alpha _k^{}(t)`$. Thus, the restriction of $`\sigma _{i,C}`$ to $`C^{}`$ coincides with $`\sigma _{i,C^{}}`$, and the result follows.
### Appendix C Proof of Lemma 13.4
We begin by collecting some more explicit information on the root system of $`\mathrm{SL}(n)`$. First, $`(\alpha |\alpha )=2`$ for every root $`\alpha `$, so in particular,
$$s_\alpha (\lambda )=\lambda (\lambda |\alpha )\alpha $$
for any weight $`\lambda `$. Recall also that all fundamental weights $`\varpi _1,\mathrm{},\varpi _{n1}`$ are minuscule, so for any $`wS_n`$ and any root $`\alpha `$, $`(w\varpi _i|\alpha )=0`$, $`1`$, or $`1`$. Moreover, if $`\alpha >0`$, then
$`w<s_\alpha w`$ $`(w\varpi _i|\alpha )=0\text{ or }1\text{,}`$
$`w>s_\alpha w`$ $`(w\varpi _i|\alpha )=0\text{ or }1\text{.}`$
Now let $`\alpha \alpha ^{}`$ be two positive roots and $`s_\alpha =(ab)`$, $`s_\alpha ^{}=(a^{}b^{})`$, with $`a<b`$ and $`a^{}<b^{}`$. Then
$$(\alpha |\alpha ^{})=\{\begin{array}{cc}1\hfill & \text{if }a=a^{}\text{ or }b=b^{}\text{ (but not both),}\hfill \\ 0\hfill & \text{if }\{a,b\}\{a^{},b^{}\}=\mathrm{}\text{,}\hfill \\ 1\hfill & \text{if }a=b^{}\text{ or }b=a^{}\text{.}\hfill \end{array}$$
The preceding information will be used freely in the sequel, without explicit reference.
We fix a simple root $`\beta `$. Consider first the action of the generator $`K_\beta `$ on a vector $`e_C^{ij}`$, where $`C=C(\alpha _1,\mathrm{},\alpha _d;w)`$ is monogressive and $`ij`$-effective. Recalling the expression (13.2) for $`\mathrm{\Delta }K_\beta `$, we get
$$K_\beta e_C^{ij}=\underset{L\{1,\mathrm{},d\}}{}q^{(s_{\overline{L}}w\varpi _i|\beta )+(s_Lw\varpi _j|\beta )}q^{|L|}e_{s_{\overline{L}}w}^ie_{s_Lw}^j.$$
For each $`1kd`$, we have $`s_{\alpha _k}w\varpi _j=w\varpi _j\alpha _k`$. More generally, $`s_Lw\varpi _j=w\varpi _j_{kL}\alpha _k`$, and similarly for $`s_{\overline{L}}w\varpi _i`$; therefore,
$$K_\beta e_C^{ij}=q^{(w(\varpi _i+\varpi _j)|\beta )_{k=1}^d(\alpha _k|\beta )}e_C^{ij}.$$
A similar formula holds for $`K_\beta ^1e_C^{ij}`$.
Now we study the action of $`X_\beta `$ and of $`Y_\beta `$ on a vector $`e_C^{ij}`$. Note that the root $`\beta `$ will be orthogonal to all defining roots of the orthocell $`C`$, except at most two. Thus, there are four cases to consider:
$`C=C(\alpha ,\alpha ^{},\alpha _1,\mathrm{},\alpha _d;w)`$,
$`C=C(\alpha ,\alpha _1,\mathrm{},\alpha _d;w)`$,
$`C=C(\alpha _1,\mathrm{},\alpha _d;w)`$,
$`C=C(\beta ,\alpha _1,\mathrm{},\alpha _d;w)`$,
where, in all cases, $`\alpha ,\alpha ^{},\alpha _1,\mathrm{},\alpha _d`$ are pairwise orthogonal, $`(\beta |\alpha )=\pm 1`$, $`(\beta |\alpha ^{})=\pm 1`$, and $`(\beta |\alpha _k)=0`$ for all $`k`$.
We will first treat these four cases when $`d=0`$, and then describe how to deduce results for arbitrary $`d`$ from this particular case.
Let us use the notation $`cc^{}`$ even when $`c,c^{}`$ are integers, meaning that $`c^{}=c+1`$. We will also use the following notation throughout:
$$s:=s_\alpha =:(ab),s^{}:=s_\alpha ^{}=:(a^{}b^{}),t:=s_\beta ,$$
with $`a<b`$ and $`a^{}<b^{}`$. Note that monogressivity and $`ij`$-effectiveness *exclude* the orderings $`a<a^{}<b<b^{}`$ and $`a^{}<a<b^{}<b`$.
Case I: $`C=C(\alpha ,\alpha ^{};w)`$.
*Subcase I.1: $`(\beta |\alpha )=(\beta |\alpha ^{})=1`$:* Exchanging $`\alpha ,\alpha ^{}`$ if necessary, we may assume that $`a<a^{}`$. Then we must have $`a<ba^{}<b^{}`$ and $`t=(ba^{})`$. Furthermore, $`(sw\varpi _i|\beta )=(s^{}w\varpi _i|\beta )=(w\varpi _i|\beta )1`$ and $`(ss^{}w\varpi _i|\beta )=(w\varpi _i|\beta )2`$. Since all these inner products must be equal to $`0`$, $`1`$, or $`1`$, we get
$$(w\varpi _i|\beta )=1,(sw\varpi _i|\beta )=(s^{}w\varpi _i|\beta )=0,(ss^{}w\varpi _i|\beta )=1,$$
and similarly for $`\varpi _i`$ replaced by $`\varpi _j`$. Recalling the expression (13.2) for $`\mathrm{\Delta }X_\beta `$ and $`\mathrm{\Delta }Y_\beta `$, we obtain
$`X_\beta e_C^{ij}`$ $`=X_\beta \left(q^2e_w^ie_{ss^{}w}^j+qe_{sw}^ie_{s^{}w}^j+qe_{s^{}w}^ie_{sw}^j+e_{ss^{}w}^ie_w^j\right)`$
$`=q^2X_\beta e_w^ie_{ss^{}w}^j+K_\beta e_{ss^{}w}^iX_\beta e_w^j`$
$`=q^2e_{tw}^ie_{ss^{}w}^j+qe_{ss^{}w}^ie_{tw}^j,`$
and, by a similar computation,
$$Y_\beta e_C^{ij}=q^2e_w^ie_{tss^{}w}^j+qe_{tss^{}w}^ie_w^j.$$
The vanishing inner products obtained above imply that $`tsw\varpi _i=sw\varpi _i`$ and $`ts^{}w\varpi _i=s^{}w\varpi _i`$. Using the Coxeter relations $`(ts)^3=(ts^{})^3=1`$, we then also have the equalities
$`tw\varpi _i`$ $`=stw\varpi _i=s^{}tw\varpi _i=ss^{}tw\varpi _i,`$
$`ss^{}w\varpi _i`$ $`=sts^{}w\varpi _i=s^{}tsw\varpi _i=ss^{}tss^{}tw\varpi _i,`$
$`w\varpi _i`$ $`=tstw\varpi _i=ts^{}tw\varpi _i=tss^{}tw\varpi _i,`$
$`tss^{}w\varpi _i`$ $`=stss^{}w\varpi _i=s^{}tss^{}w\varpi _i=ss^{}tss^{}w\varpi _i`$
(and similarly for $`\varpi _i`$ replaced by $`\varpi _j`$), which may be used in the above expressions for $`X_\beta e_C^{ij}`$ and $`Y_\beta e_C^{ij}`$; cf. Remark 10.1. Since $`a<b<a^{}<b^{}`$, monogressivity implies that there are four possible relative positions of $`a,b,a^{},b^{}`$ inside the array $`w`$, yielding the following expressions for $`X_\beta e_C^{ij}`$ and $`Y_\beta e_C^{ij}`$:
| $`w`$ | $`X_\beta e_C^{ij}`$ | $`Y_\beta e_C^{ij}`$: |
| --- | --- | --- |
| $`[aa^{}bb^{}]`$ | $`qe_{C(ss^{}(\beta );ss^{}tw)}^{ij}`$ | $`qe_{C(ss^{}(\beta );tss^{}tw)}^{ij}`$, |
| $`[a^{}ab^{}b]`$ | $`qe_{C(ss^{}(\beta );tw)}^{ij}`$ | $`qe_{C(ss^{}(\beta );w)}^{ij}`$, |
| $`[a^{}abb^{}]`$ | $`qe_{C(ss^{}(\beta );s^{}tw)}^{ij}`$ | $`qe_{C(ss^{}(\beta );ts^{}tw)}^{ij}`$, |
| $`[aa^{}b^{}b]`$ | $`qe_{C(ss^{}(\beta );stw)}^{ij}`$ | $`qe_{C(ss^{}(\beta );tstw)}^{ij}`$. |
These results are valid provided all orthocells involved are monogressive (their $`ij`$-effectiveness being clear). In each case, this may easily be checked using Criterion 11.1. For example, in the first line, we have $`w=[aa^{}bb^{}]`$, $`ss^{}tw=[bab^{}a^{}]`$, and $`s_{ss^{}(\beta )}=(ab^{})`$: since $`C=C(\alpha ,\alpha ^{};w)`$ is monogressive by assumption, the subarray $`[a^{}\mathrm{}b]`$ of $`w`$ contains no numbers in the (numerical) interval $`[a,b]`$, nor in $`[a^{},b^{}]`$, and therefore not in $`[a,b^{}]`$ (because $`a<ba^{}<b^{}`$); thus, the orthocell $`C(ss^{}(\beta );ss^{}tw)`$ is again monogressive.
*Subcase I.2: $`(\beta |\alpha )=1`$ and $`(\beta |\alpha ^{})=1`$:* Here we must have either $`aa^{}<b^{}<b`$ and $`t=(aa^{})`$, or $`a<a^{}<b^{}b`$ and $`t=(b^{}b)`$. Furthermore,
$$(w\varpi _i|\beta )=(ss^{}w\varpi _i|\beta )=0,(sw\varpi _i|\beta )=1,(s^{}w\varpi _i|\beta )=1,$$
and similarly for $`\varpi _i`$ replaced by $`\varpi _j`$. It follows that
$`X_\beta e_C^{ij}`$ $`=qe_{tsw}^ie_{s^{}w}^j+q^2e_{s^{}w}^ie_{tsw}^j,`$
$`Y_\beta e_C^{ij}`$ $`=qe_{sw}^ie_{ts^{}w}^j+q^2e_{ts^{}w}^ie_{sw}^j.`$
Arguments similar to those of Subcase I.1 then show (whether $`t=(aa^{})`$ or $`t=(b^{}b)`$) that the orthocells $`C(ss^{}(\beta );s^{}tw)`$ and $`C(ss^{}(\beta );ts^{}tw)`$ are monogressive and that
$$X_\beta e_C^{ij}=qe_{C(ss^{}(\beta );s^{}tw)}^{ij},Y_\beta e_C^{ij}=qe_{C(ss^{}(\beta );ts^{}tw)}^{ij}.$$
(We omit the details.)
*Subcase I.3: $`(\beta |\alpha )=(\beta |\alpha ^{})=1`$:* These inner products force $`a<a^{}<b<b^{}`$ or $`a^{}<a<b^{}<b`$, contradicting the fact that $`C(\alpha ,\alpha ^{};w)`$ is monogressive and $`ij`$-effective. Therefore, this Subcase is impossible.
Case II: $`C=C(\alpha ;w)`$.
*Subcase II.1: $`(\beta |\alpha )=1`$:* We must have either $`ca<b`$ and $`t=(ca)`$, or $`a<bc`$ and $`t=(bc)`$. Moreover, since $`(sw\varpi _i|\beta )=(w\varpi _i|\beta )+1`$, and similarly for $`\varpi _j`$, we obtain the following cases.
* If $`(w\varpi _i|\beta )=(w\varpi _j|\beta )=1`$, then arguments similar to those of Case I show that
$$Y_\beta e_C^{ij}=0$$
and that $`X_\beta e_C^{ij}=qe_{tw}^ie_{sw}^j+e_{sw}^ie_{tw}^j`$ is given by the following table:
| $`t`$ | $`w`$ | $`X_\beta e_C^{ij}`$: |
| --- | --- | --- |
| $`(ca)`$ | $`[abc]`$ | $`e_{C(s(\beta );tw)}^{ij}`$, |
| $`(ca)`$ | $`[acb]`$ | $`e_{C(s(\beta );stw)}^{ij}`$, |
| $`(bc)`$ | $`[acb]`$ | $`e_{C(s(\beta );stw)}^{ij}`$, |
| $`(bc)`$ | $`[cab]`$ | $`e_{C(s(\beta );tw)}^{ij}`$. |
* If $`(w\varpi _i|\beta )=(w\varpi _j|\beta )=0`$, then
$$X_\beta e_C^{ij}=0$$
and $`Y_\beta e_C^{ij}=qe_w^ie_{tsw}^j+e_{tsw}^ie_w^j`$ is given by the following table:
| $`t`$ | $`w`$ | $`Y_\beta e_C^{ij}`$: |
| --- | --- | --- |
| $`(ca)`$ | $`[cab]`$ | $`e_{C(s(\beta );tw)}^{ij}`$, |
| $`(ca)`$ | $`[acb]`$ | $`e_{C(s(\beta );w)}^{ij}`$, |
| $`(bc)`$ | $`[acb]`$ | $`e_{C(s(\beta );w)}^{ij}`$, |
| $`(bc)`$ | $`[abc]`$ | $`e_{C(s(\beta );tw)}^{ij}`$. |
* If $`(w\varpi _i|\beta )=1`$ and $`(w\varpi _j|\beta )=0`$, or vice versa, then
$$X_\beta e_C^{ij}=qe_{C(\text{};stw)}^{ij},Y_\beta e_C^{ij}=qe_{C(\text{};stsw)}^{ij}.$$
*Subcase II.2: $`(\beta |\alpha )=1`$:* Here we have either $`ac<b`$ and $`t=(ac)`$, or $`a<cb`$ and $`t=(cb)`$. Moreover, since $`(sw\varpi _i|\beta )=(w\varpi _i|\beta )1`$, and similarly for $`\varpi _j`$, we obtain the following cases.
* If $`(w\varpi _i|\beta )=(w\varpi _j|\beta )=1`$, then
$$X_\beta e_C^{ij}=0,Y_\beta e_C^{ij}=\pm e_{C(s(\beta );tw)}^{ij}.$$
* If $`(w\varpi _i|\beta )=(w\varpi _j|\beta )=0`$, then
$$X_\beta e_C^{ij}=\pm e_{C(s(\beta );tw)}^{ij},Y_\beta e_C^{ij}=0.$$
* The case $`(w\varpi _i|\beta )=1`$ and $`(w\varpi _j|\beta )=0`$, or vice versa, contradicts the $`ij`$-effectiveness of $`C`$ and is therefore impossible.
Case III: $`C=C(\text{};w)`$. We obtain the following table:
| $`(w\varpi _i|\beta )`$ | $`(w\varpi _j|\beta )`$ | $`X_\beta e_C^{ij}`$ | $`Y_\beta e_C^{ij}`$: |
| --- | --- | --- | --- |
| $`1`$ | $`1`$ | $`0`$ | $`q^1e_{C(\beta ;w)}^{ij}`$, |
| $`1`$ | $`0`$ | $`0`$ | $`e_{C(\text{};tw)}^{ij}`$, |
| $`0`$ | $`1`$ | $`0`$ | $`e_{C(\text{};tw)}^{ij}`$, |
| $`0`$ | $`0`$ | $`0`$ | $`0`$, |
| $`0`$ | $`1`$ | $`e_{C(\text{};tw)}^{ij}`$ | $`0`$, |
| $`1`$ | $`0`$ | $`e_{C(\text{};tw)}^{ij}`$ | $`0`$, |
| $`1`$ | $`1`$ | $`q^1e_{C(\beta ;tw)}^{ij}`$ | $`0`$. |
Case IV: $`C=C(\beta ;w)`$. Since $`C`$ is monogressive and $`ij`$-effective, we must have $`(w\varpi _i|\beta )=(w\varpi _j|\beta )=1`$, hence
$$X_\beta e_C^{ij}=(q^2+1)e_{C(\text{};w)}^{ij},Y_\beta e_C^{ij}=(q^2+1)e_{C(\text{};tw)}^{ij}.$$
Finally, we show how, in the preceding four cases, one can deduce the action of $`X_\beta `$ and $`Y_\beta `$ for arbitrary $`d`$ from that for $`d=0`$. The idea is that $`\alpha _1,\mathrm{},\alpha _d`$, being orthogonal to $`\beta ,\alpha ,\alpha ^{}`$, do not “interfere” with the computations done above. To make this idea precise, we will restrict ourselves to the very first case treated above (all other cases being similar), namely, the action of $`X_\beta `$ on $`e_C^{ij}`$ when $`C=C(\alpha ,\alpha ^{},\alpha _1,\mathrm{},\alpha _d;w)`$, $`(\beta |\alpha )=(\beta |\alpha ^{})=1`$, $`(\beta |\alpha _k)=0`$ for all $`k`$, $`s_\alpha =(ab)`$ and $`s_\alpha ^{}=(a^{}b^{})`$ with $`a<ba^{}<b^{}`$ (so $`t:=s_\beta =(ba^{})`$), and $`w=[aa^{}bb^{}]`$.
Since $`\beta `$ is orthogonal to each $`\alpha _k`$, we have $`(s_L\lambda |\beta )=(\lambda |\beta )`$ for any weight $`\lambda `$ and any $`L\{1,\mathrm{},d\}`$, so we still get
$$(s_Lw\varpi _i|\beta )=1,(ss_Lw\varpi _i|\beta )=(s^{}s_Lw\varpi _i|\beta )=0,(ss^{}s_Lw\varpi _i|\beta )=1.$$
Therefore, the action of $`X_\beta `$ on each term of
$`e_{C(\alpha ,\alpha ^{},\alpha _1,\mathrm{},\alpha _d;w)}^{ij}={\displaystyle \underset{L\{1,\mathrm{},d\}}{}}q^{|L|}`$ $`(q^2e_{s_{\overline{L}}w}^ie_{ss^{}s_Lw}^j+qe_{ss_{\overline{L}}w}^ie_{s^{}s_Lw}^j`$
$`+qe_{s^{}s_{\overline{L}}w}^ie_{ss_Lw}^j+e_{ss^{}s_{\overline{L}}w}^ie_{s_Lw}^j)`$
is still computed in a similar way to that on $`e_{C(\alpha ,\alpha ^{};w)}^{ij}`$, viz.
$$X_\beta e_{C(\alpha ,\alpha ^{},\alpha _1,\mathrm{},\alpha _d;w)}^{ij}=\underset{L\{1,\mathrm{},d\}}{}q^{|L|}\left(q^2e_{ts_{\overline{L}}w}^ie_{ss^{}s_Lw}^j+qe_{ss^{}s_{\overline{L}}w}^ie_{ts_Lw}^j\right).$$
It follows that
$$X_\beta e_{C(\alpha ,\alpha ^{},\alpha _1,\mathrm{},\alpha _d;w)}^{ij}=qe_{C(ss^{}(\beta ),\alpha _1,\mathrm{},\alpha _d;ss^{}tw)}^{ij},$$
provided the orthocell $`C(ss^{}(\beta ),\alpha _1,\mathrm{},\alpha _d;ss^{}tw)`$ is monogressive. But it is easy to see that the analysis of the monogressivity of $`C(ss^{}(\beta );ss^{}tw)`$ made earlier, using Criterion 11.1, remains valid if $`\alpha _1,\mathrm{},\alpha _d`$ are added to the orthocells $`C(\alpha ,\alpha ^{};w)`$ and $`C(ss^{}(\beta );ss^{}tw)`$.
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# Search for light-to-heavy quark flavor changing neutral currents in 𝜈_𝜇𝑁 and 𝜈̄_𝜇𝑁 scattering at the Tevatron
## I Introduction
Flavor-changing neutral current (FCNC) interactions of $`c`$\- and $`b`$-quarks appear in a number of extensions to the Standard Model (SM) of particle physics including extra quark generations , technicolor, multiple Higgs sectors (as in supersymmetry), left-right symmetric models, and leptoquarks. Evidence for FCNC effects in the heavy quark sector beyond higher order SM processes has not yet been observed. Present limits on FCNC result from searches for rare decays of charm and beauty mesons, in particular decays of the type $`D\mathrm{}^+\mathrm{}^{}X`$ and $`B\mathrm{}^+\mathrm{}^{}X`$, where $`\mathrm{}=e`$ or $`\mu `$, $`D=(D^+,D^0,D_S^+)`$, $`B=(B^+,B^0,B_S^0)`$; and $`X`$ is nothing, a pseudoscalar, or a vector meson. The $`cu`$ transitions are particularly sensitive to new physics since loop level SM-FCNC decays are severely suppressed by the Cabibbo-Kobayashi-Maskawa (CKM) matrix. While experimental signatures for FCNC in $`D`$ and $`B`$ decays are clear, their interpretation is ambiguous. Meson decay rates depend on one or more incalculable hadronic form factors. In addition, experimentally attractive final states such as $`D^0e^+e^{}`$ and $`B^0\mu ^+\mu ^{}`$are helicity-suppressed, which obscures dynamical roles played by particular FCNC models.
This article presents an alternative search for FCNC processes in the DIS data of the NuTeV experiment; where either neutrinos or anti-neutrinos interact with a massive iron target. Flavor changing effects can be sought in the reactions
$`\nu _\mu N`$ $`\nu _\mu cX\text{}c\mu ^+X^{},`$ (1)
$`\nu _\mu N`$ $`\nu _\mu \overline{b}X\text{}\overline{b}\mu ^+X^{},`$ (2)
$`\nu _\mu N`$ $`\nu _\mu bX\text{}bcX^{},c\mu ^+X^{\prime \prime },`$ (3)
and their charge-conjugates. The experimental signature in the detector is a muon of opposite lepton number from the beam neutrino. It is possible to isolate this final state because NuTeV ran with a high purity sign-selected beam in which the $`\overline{\nu }_\mu /\nu _\mu `$ event ratio in neutrino mode and the $`\nu _\mu /\overline{\nu }_\mu `$ ratio in anti-neutrino mode were $`0.8\times 10^3`$ and $`4.8\times 10^3`$, respectively. Because of the semi-inclusive character of the measurement, FCNC effects in neutrino scattering can be probed at the quark, rather than the hadron level. Furthermore, neutrino scattering is particularly sensitive to any FCNC process mediated by an intermediate neutral object that couples more strongly to neutrinos than to charged leptons (e.g., a $`Z^0`$-like coupling).
## II Experimental Apparatus and Beam
The NuTeV (Fermilab-E815) neutrino experiment collected data during the 1996-97 fixed target run with the refurbished Lab E neutrino detector and a newly installed Sign-Selected Quadrupole Train (SSQT) neutrino beamline. The sign-selection optics of the SSQT pick the charge of secondary pions and kaons which determines whether $`\nu _\mu `$ or $`\overline{\nu }_\mu `$ are predominantly produced. During NuTeV’s run the primary production target received $`1.13\times 10^{18}`$ and $`1.41\times 10^{18}`$ protons-on-target in neutrino and anti-neutrino modes, respectively. The SSQT and its performance are described in detail elsewhere.
The Lab E detector, consists of two major parts; a target calorimeter and an iron toroid spectrometer. The target calorimeter contains 690 tons of steel sampled at 10 cm intervals by 84 3m$`\times `$3m scintillator counters and at 20 cm intervals by 42 3m$`\times `$3m drift chambers. The toroid spectrometer consists of four stations of drift chambers separated by iron toroid magnets. Precision hadron and muon calibration beams monitored the calorimeter and spectrometer performance throughout the course of data taking. The calorimeter achieves a sampling-dominated hadronic energy resolution of $`\sigma _{E_{HAD}}/E_{HAD}=2.4\%87\%/\sqrt{E_{HAD}}`$ and an absolute scale uncertainty of $`\delta E_{HAD}/E_{HAD}=0.5\%`$. The spectrometer’s multiple Coulomb scattering dominated muon energy resolution is $`\sigma _{E_\mu }/E_\mu =`$ $`11\%`$ and the muon momentum scale is known to $`\delta E_\mu /E_\mu =1.0\%`$. With the selection criteria used in this analysis, the muon charge mis-identification probability in the spectrometer is $`2\times 10^5`$. This latter rate is confirmed by measurement the muon calibration beam.
## III Analysis Procedure
### A Introduction and Data Selection
The analysis technique consists of comparing the distributions of $`y_{VIS}=E_{HAD}/\left(E_{HAD}+E_\mu \right)`$ measured in the $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ wrong sign muon (WSM) data samples to a Monte Carlo(MC) simulation containing all known conventional WSM sources and a possible FCNC signal. The FCNC signal peaks at high $`y_{VIS}`$ because the decay muon from the heavy flavor hadron is usually much less energetic than the hadron energy produced in the NC interaction. The largest backgrounds, from beam impurities, are concentrated at low $`y_{VIS}`$ in $`\nu _\mu `$ and distributed evenly across $`y_{VIS}`$ in $`\overline{\nu }_\mu `$ mode due to the respective $`\left(1y\right)^2`$ and uniform-in-$`y`$ characteristics of the CC interactions of wrong-flavor beam backgrounds.
Events in the WSM sample must satisfy a number of selection criteria(“cuts”). The fiducial volume cut requires that event vertices be reconstructed at least 25 cm-Fe (cm of iron-equivalent) from the outer edges of the detector in the transverse directions, at least 35 cm-Fe downstream of the upstream face of the detector, and at least 200 cm-Fe upstream of the toroid. Events must possess a hadronic energy of at least 10 GeV, and exactly one track (the muon) must be found. The muon is required to be well-reconstructed and to pass within the understood regions of the toroid’s magnetic field. The muon’s energy must be between 10 and 150 GeV, and its charge must be consistent with having the opposite lepton number as the primary beam component. Requiring that the muon energy reconstructed in different longitudinal sections of the toroid agree within 25% of the value measured using the full toroid reduces charge mis-identification backgrounds to the $`2\times 10^5`$ level. Finally, for the purposes of the final FCNC fit, the reconstructed $`y_{VIS}`$ is required to be larger than 0.5. With these cuts there are 207 $`\nu `$-mode and 127 $`\overline{\nu }`$-mode WSM events remaining in NuTeV’s nearly 2 million single muon sample.
### B Source and Background Simulations
Conventional WSM sources arise from beam impurities, right-flavor charged current (CC) events where the charge of the muon is mis-reconstructed, CC and NC events where a $`\pi `$ or $`K`$ decays in the hadron shower, charged current (CC) charm production where the primary muon is not reconstructed or the charm quark is produced via a $`\nu _e`$ interaction, and neutral current (NC) $`c\overline{c}`$ pair production. Single charm CC production and NC $`c\overline{c}`$ pair production background sources produce broad peaks at high $`y_{VIS}`$ and must be handled with care. Table I gives the fractional contribution of each background component. The relatively large beam impurity background consists of contributions from hadrons (including charm) that decay before the sign-selecting dipoles in the SSQT, neutral kaon decays, muon decays, decay of hadrons produced by secondary interactions in the SSQT (“scraping”), and from decay of wrong-sign pions produced in kaon decays. Table II summarizes the relative contributions of each beam source.
A complete GEANT simulation of the SSQT is used to model beam impurities. This simulation uses Malensek’s parameterization for hadron production from the primary target. Scraping contributions are modeled by GHEISHA. Production of $`K_L^0`$ is handled by extending Malensek’s charged kaon parameterizations using the quark counting relation $`K_L^0=\left(3K^{}+K^+\right)/4`$. Charm production is parametrized using available data from 800 GeV proton beams . GEANT properly handles cascade decays such as $`K^\pm \pi ^\pm \pi ^\pm \pi ^{},\pi ^{}\mu ^{}\overline{\nu }_\mu (\nu _\mu )`$ and $`\pi ^\pm \mu ^\pm \overline{\nu }_\mu (\nu _\mu ),\mu ^\pm e^\pm \overline{\nu }_\mu (\nu _\mu )\nu _e(\overline{\nu }_e)`$. The NuTeV detector is likewise modeled with a GEANT-based hit-level MC simulation. Wrong-sign muons generated from the flux simulation are propagated through the detector MC and then reconstructed using the same package that is used for data reconstruction. A fast parametric MC is also used to compare the high statistics right-sign flux simulation to data in $`\nu _\mu `$ and $`\overline{\nu }_\mu `$. These comparisons showed that the SSQT dipoles required a downward shift of -2.5% from their nominal values. The right-sign comparisons after these shifts, are shown in Fig. 1 and indicate agreement between predicted flux and data at roughly the $`2\%`$ level.
The high density target-calorimeter suppresses WSM contributions from $`\pi /K`$ decay in the hadron shower; their contribution is estimated from a previous measurement of $`\mu `$-production in hadron showers using the same detector . The small charge mis-identification contribution is estimated by passing a large sample of simulated events through the full detector MC and event reconstruction.
After impurities, the next largest WSM source comes from CC production of charm in which the charm quark decays semi-muonically (dimuon) and its decay muon is picked up in the spectrometer while the primary lepton is either an electron or a muon which exits from or ranges out in the calorimeter. The $`\nu _e`$ beam fraction is $`1.9(1.3)\%`$ in $`\nu `$ ($`\overline{\nu }`$)-mode, and $`22\%`$ of the CC charm events which pass WSM cuts originate from a $`\nu _e`$. The CC charm background is simulated using a leading-order QCD charm production model with production, fragmentation, and charm decay parameters tuned on neutrino dimuon data collected by NuTeV and a previous experiment using the same detector. Overall normalization of the source is obtained from the measured charm-to-total CC cross section ratio and the single muon right-sign data sample. Simulated dimuon events are passed through the full GEANT simulation of the detector. Fig. 2 provides a check of the modeling of this source through a comparison of the distribution of $`y_{VIS}^{}=`$ $`E_{HAD}/\left(E_{HAD}+E_{\mu 2}\right)`$, where $`E_{\mu 2}`$ is the energy of the WSM in the event, between data and MC for dimuon events in which both muons are reconstructed by the spectrometer. This distribution should closely mimic the expected background to the $`y_{VIS}`$ distribution in the WSM sample. A $`\chi ^2`$ comparison test between data and model yields a value of 19 for 17 degrees of freedom.
Finally, NC $`c\overline{c}`$ production produces a WSM when the $`c\left(\overline{c}\right)`$ decays semi-muonically in $`\nu _\mu \left(\overline{\nu }_\mu \right)`$ mode. An excess over other sources at high $`y_{VIS}`$ indicates that this source is present in the data; its analysis will appear in a forthcoming publication. For the FCNC search, NC charm production is simulated at production level by a $`Z^0`$gluon fusion model with charm mass parameter $`m_c=1.70\pm 0.19`$ GeV$`/c^2`$ taken from a next-to-leading (NLO) order QCD analysis of CC charm production and using the GRV94HO gluon parton distribution function (PDF). The NLO charm mass is used because it is influenced in part by contributions from $`W`$gluon fusion diagrams similar to the $`Z^0`$gluon process. Note that the value of $`m_c`$ used is larger than that obtained in LO analyses of CC charm production. This choice tends to reduce the NC charm contribution to the WSM sample and results in more conservative limits on FCNC production. The NC charm quarks are fragmented and decayed using procedures adapted from the CC charm simulation, and the resulting WSM events are then simulated with the full MC.
## IV Results and Interpretation
### A FCNC Production
The neutrino FCNC $`uc`$ cross section can be parameterized to LO in QCD as
$$\frac{d\sigma \left(\nu _\mu u\nu _\mu c;c\mu ^+\right)}{d\xi dy}=\left|\frac{V_{uc}}{V_{cd}}\right|^2\left[\mathrm{cos}^2\beta +\mathrm{sin}^2\beta \frac{\left(1y\right)\left(1xy/\xi \right)}{1y+xy/\xi }\right]\frac{d\sigma \left(\nu _\mu d\mu ^{}c;c\mu ^+\right)}{d\xi dy}.$$
(4)
Here $`V_{cd}`$ is the $`cd`$ CKM matrix element, $`V_{uc}`$<sup>*</sup><sup>*</sup>*We use the notation $`V_{uc}`$, $`V_{bd}`$, and $`V_{sd}`$ in simple analogy to the CKM matrix in order to compare our results to those from FCNC decay searches.. We do not assume any constraints exist for this FCNC CKM-like matrix. In our notation, the FCNC left and right-handed couplings for charm are $`g_L^2=\left|V_{uc}\right|^2\mathrm{cos}^2\beta `$ and $`g_R^2=\left|V_{uc}\right|^2\mathrm{sin}^2\beta `$. represents a possible $`uc`$ coupling, $`\mathrm{sin}^2\beta `$ gives the fraction of right-handed coupling of the $`c`$quark to the FCNC, $`y`$ is the inelasticity, and $`\xi x\left(1+m_c^2/Q^2\right)`$ is the fraction of the nucleon’s momentum carried by the struck $`u`$quark, with $`x`$ the Bjorken scaling variable, $`Q^2`$ the squared momentum transfer, and $`m_c`$ the effective charm quark mass. The $`dc`$ charged current cross section $`d\sigma \left(\nu _\mu d\nu _\mu c;c\mu ^+\right)/d\xi dy`$ is measured in the same experiment . Since the $`u`$ and $`d`$ quark distributions are identical in an isoscalar target, the FCNC cross section should experience the same charm mass suppression as the analogous CC charm production. Fragmentation of subsequent semi-muonic decays of charmed mesons should also be identical for FCNC and CC-charm production. One therefore expects the extracted $`V_{uc}`$ to have little model dependence.
For FCNC bottom production there is as yet no measured CC analog final state. Therefore, the explicit LO QCD cross section,
$`{\displaystyle \frac{d\sigma \left(\nu _\mu N\nu _\mu \overline{b}X\right)}{d\xi ^{}dy}}`$ $`=`$ $`{\displaystyle \frac{G_F^2ME\left|V_{bd}\right|^2}{\pi }}\left[\mathrm{cos}^2\beta ^{}\left(1y\right)\left(1xy/\xi \right)+\mathrm{sin}^2\beta ^{}\left(1y+xy/\xi \right)\right]`$ (6)
$`\times \left(\overline{u}(\xi ^{},Q^2)+\overline{d}(\xi ^{},Q^2)\right),`$
where $`M`$ is the nucleon mass and $`E`$ is the neutrino energy, must be convolved with $`b`$-quark fragmentation functions for mesons of type $`B_i`$ $`\left(D_b^_i\right)`$ and $`B_i`$ meson decay distribution functions $`\left(\mathrm{\Delta }_B^i\right)`$ multiplied by appropriate branching fractions $`\left(F_B^i\right)`$ to yield a WSM cross section:
$$\frac{d\sigma \left(\nu _\mu N\nu _\mu \overline{b};\overline{b}\mu ^+\right)}{d\xi ^{}dy}=\underset{i}{}\frac{d\sigma \left(\nu _\mu N\nu _\mu \overline{b}X\right)F_B^iD_b^i\mathrm{\Delta }_B^i}{d\xi ^{}dy}.$$
(7)
The struck quark momentum fraction $`\xi ^{}`$ becomes $`\xi ^{}x\left(1+m_b^2/Q^2\right),`$ with $`m_b=4.8`$ GeV$`/c^2`$ the effective $`b`$-quark mass. It is also possible for FCNC $`b`$-production to form a WSM muon signal through the cascade $`bc\mu ^+`$. This mode offers the advantages of the larger and higher $`\xi `$ valence $`d`$-quark PDF at the cost of reduced acceptance for the softer $`c`$-decay muon. A similar expression holds for FCNC $`sb`$ transitions with the replacements $`u(\xi ^{},Q^2)+d(\xi ^{},Q^2)2s(\xi ^{},Q^2)`$, $`\left|V_{bd}\right|^2\left|V_{bs}\right|^2`$, and $`\mathrm{sin}^2\beta ^{}\mathrm{sin}^2\beta ^{^{\prime \prime }}`$.
Production cross sections for both $`c`$\- and $`b`$\- FCNC sources are computed from the GRV94LO PDF set for several choices of right-left coupling admixtures. Acceptance for a charm FCNC-WSM signal is calculated using a fragmentation-decay model tuned to NuTeV and CCFR dimuon data. For FCNC-WSM from $`b`$-quarks, fragmentation and decays are handled with the Lund string fragmentation model . Detector response is simulated with the full hit-level MC.
### B Fits to Data
Binned likelihood fits are performed to the $`y_{VIS}`$ distributions of the data using a model consisting of all conventional WSM sources described above and an FCNC source. The fit varies the level, but not the shape, of the FCNC signal contribution. The NC charm contribution is also varied in shape and level by allowing $`m_c`$ to float within its errors. The three FCNC sources ($`uc`$, $`db`$, and $`sb`$) are treated separately. Only neutrino data is used for the $`uc`$, but both modes are used for FCNC bottom production to exploit the possibility of a cascade decays to charm. A series of fits are performed for each FCNC source, corresponding to different mixtures of right and left-handed FCNC couplings to the quarks; a typical result is shown in Fig. 3.
In all cases, the signal for FCNC is within $`\pm `$ 2.0 $`\sigma `$ of zero, and limits are set accordingly. Since Gaussian statistics apply, the $`90\%`$ confidence level upper limit is set by adding $`1.64\sigma `$ to the best-fit value if the best-fit value is positive, or $`1.64\sigma `$ to zero if the best fit is negative. Here, $`\sigma `$ consists of the statistical error from the fit added in quadrature to the estimated systematic error described in the next section. Table III summarizes the fit results.
### C Systematic Errors
The dominant systematic errors result from modeling the rejection of CC charm events, and the overall normalization of CC charm events. Estimates of systematic uncertainties are obtained by varying the event selection procedure as well as parameters characterizing the detector response and physics models. Errors are assumed to be independent.
Charged current charm events are removed by requiring that exactly one track be found and reconstructed by the NuTeV tracking software. Another independent way to remove dimuons is to use calorimeter information. The stop parameter is the first of three consecutive counters downstream of the interaction, each with less than 1.5 MIPs. The stop cut requires that the distance between the interaction and the stop counter be less than 15 counters. Replacing the tracking cut with the stop cut gives the systematic errors listed in Table IV.
The next largest systematic error is due to the normalization of CC charm events. Normalization of these events is obtained from the right-sign muon CC sample. One can also normalize CC charm events with only one reconstructed track, to those with both tracks found. These normalizations disagree by 3% resulting in the systematic errors listed in Table IV. Systematic errors due to the beam normalization, detector calibration, and other sources are small.
### D Comparison to Limits from Decays
For comparison purposes, the following expressions are used to relate FCNC heavy flavor meson decay branching fractions $`\left(BF\right)`$ to the parameter $`V_{uc}`$:
$`BF\left(D^0\mathrm{}^+\mathrm{}^{}\right)`$ $`=2\left|{\displaystyle \frac{V_{uc}}{V_{cs}}}\right|^2{\displaystyle \frac{m_{\mathrm{}}^2}{m_\mu ^2}}BF\left(D_S^+\mu ^+\nu _\mu \right),`$ (8)
$`BF\left(D^+\pi ^+\mathrm{}^+\mathrm{}^{}\right)`$ $`=\left|{\displaystyle \frac{V_{uc}}{V_{cd}}}\right|^2BF\left(D^+\pi ^+\mathrm{}^+\nu _{\mathrm{}}\right),`$ (9)
$`BF\left(D_S^+K^+\mathrm{}^+\mathrm{}^{}\right)`$ $`=\left|{\displaystyle \frac{V_{uc}}{V_{cs}}}\right|^2BF\left(D_S^+\eta \mathrm{}^+\nu _{\mathrm{}}\right).`$ (10)
For estimates of $`V_{db}`$ and $`V_{sb}`$ from $`B`$ decays, it is assumed that
$`BF\left(B^0\mathrm{}^+\mathrm{}^{}\right)`$ $`=2\left|{\displaystyle \frac{V_{bd}}{V_{ub}}}\right|^2{\displaystyle \frac{m_{\mathrm{}}^2}{m_\mu ^2}}BF\left(B^+\mu ^+\nu _\mu \right),`$ (11)
$`BF\left(B^+\pi ^+\mathrm{}^+\mathrm{}^{}\right)`$ $`=\left|{\displaystyle \frac{V_{bd}}{V_{ub}}}\right|^2BF\left(B^0\pi ^{}\mathrm{}^+\nu _{\mathrm{}}\right),`$ (12)
$`BF\left(B_s^0\mathrm{}^+\mathrm{}^{}\right)`$ $`=2\left|{\displaystyle \frac{V_{bs}}{V_{ub}}}\right|^2{\displaystyle \frac{m_{\mathrm{}}^2}{m_\mu ^2}}BF\left(B^+\mu ^+\nu _\mu \right),`$ (13)
$`BF\left(B^+K^+\mathrm{}^+\mathrm{}^{}\right)`$ $`=\left|{\displaystyle \frac{V_{bs}}{V_{cb}}}\right|^2BF\left(B^+D^0\mathrm{}^+\nu _{\mathrm{}}\right).`$ (14)
Measured values are used for the branching fractions on the right hand side except for the leptonic decay $`B^+\mu ^+\nu _\mu ,`$ for which it is assumed that
$$BF\left(B^+\mu ^+\nu _\mu \right)=2.2\times 10^6\left(f_B/200\text{ MeV}\right)^2,$$
(15)
with $`f_B=200`$ MeV, the $`B`$ decay constant.
Table V summarizes the limits on $`\left|V_{uc}\right|^2`$, $`\left|V_{bd}\right|^2`$, and $`\left|V_{sb}\right|^2`$ from meson decays. We note that our overall limits from neutrino scattering, which would approximately correspond to decay searches of the type $`D\nu _\mu \overline{\nu }_\mu X`$ and $`B\nu _\mu \overline{\nu }_\mu X`$, are generally weaker than the decay search limits. Our result for $`V_{db}`$ is competitive, and we have effectively added new modes to the search that do not depend on specific mechanisms for heavy meson decay.
## V Conclusion
In this paper we have established a new method for probing FCNC processes in deep inelastic neutrino scattering. Our experiment tests for FCNC at the inclusive quark level, and we are particularly sensitive to any FCNC process in which the mediating field couples more strongly to neutrinos than to charged leptons. We observe no evidence for FCNC interactions, and we set limits on the effective mixing elements $`\left|V_{uc}\right|^2`$, $`\left|V_{bd}\right|^2`$, and $`\left|V_{bs}\right|^2`$ at the $`10^3`$ level.
###### Acknowledgements.
We would like to thank the staffs of the Fermilab Particle Physics and Beams Divisions for their contributions to the construction and operation of the NuTeV beamlines. We would also like to thank the staffs of our home institutions for their help throughout the running and analysis of NuTeV. This work has been supported by the U.S. Department of Energy and the National Science Foundation.
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# Dissipation and Quantization
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## Abstract
We show that the dissipation term in the Hamiltonian for a couple of classical damped-amplified oscillators manifests itself as a geometric phase and is actually responsible for the appearance of the zero point energy in the quantum spectrum of the 1D linear harmonic oscillator. We also discuss the thermodynamical features of the system. Our work has been inspired by ’t Hooft proposal according to which information loss in certain classical systems may lead to “an apparent quantization of the orbits which resembles the quantum structure seen in the real world”.
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Recently Gerard ’t Hooft has discussed classical, deterministic, dissipative models and has shown that constraints imposed on the solutions which introduce information loss provide bounded Hamiltonians and lead to “an apparent quantization of the orbits which resembles the quantum structure seen in the real world”. ’t Hooft’s conjecture is that the dissipation of information which would occur at Planck scale in a regime of completely deterministic dynamics would play a rôle in the quantum mechanical nature of our world.
The purpose of this Letter is to present an example of dissipation in a classical system which explicitly leads, under suitable conditions, to a quantum behavior. We show that the dissipation term in the Hamiltonian for a couple of classical damped-amplified oscillators manifests itself as a geometric phase and is actually responsible for the appearance of the zero point energy in the quantum spectrum of the 1D linear harmonic oscillator. We also discuss the thermodynamical features of the system.
The considered system of two oscillators, one damped, the other one amplified, has revealed to be an useful playground for the study of several problems of physical interest, such as the study of phase coherence in quantum Brownian motion , the study of topologically massive gauge theories in the infrared region in $`2+1`$ dimensions , the Chern-Simons-like dynamics of Bloch electrons in solids , and exhibits features also common to the structure of two-dimensional gravity models .
By imposing a condition of “adiabaticity” on such a classical system of two dissipative oscillators, we obtain a one–dimensional quantum oscillator with zero point energy originating in the geometric phase of the classical system. Thus our conclusions seem to support ’t Hooft’s analysis.
Let us first briefly outline ’t Hooft’s scenario. He considered Hamiltonians of the form $`H=_ip_if_i(q)`$, where $`f_i(q)`$ are non–singular functions of the canonical coordinates $`q_i`$. The crucial point is that equations for the $`q`$’s (i.e. $`\dot{q_i}=\{q_i,H\}=f_i(q)`$) are decoupled from the conjugate momenta $`p_i`$. This means that the system can be described deterministically even when expressed in terms of operators acting on some functional space of states, such as the Hilbert space. This is possible because it exists, for such systems, a complete set of observables which Poisson commute at all times. They are called beables . Of course, such a description in terms of operators and Hilbert space, does not implies per se quantization of the system. Quantization is in fact achieved only as a consequence of the dissipation of information, according to ’t Hooft’s scenario.
The above mentioned class of Hamiltonians is, however, not bounded from below. This might be cured by splitting $`H`$ as: $`H=H_1H_2`$, with
$`H_1={\displaystyle \frac{1}{4\rho }}\left(\rho +H\right)^2,H_2={\displaystyle \frac{1}{4\rho }}\left(\rho H\right)^2,`$ (1)
and $`\rho `$ a certain time–independent, positive function of $`q_i`$. As a result, $`H_1`$ and $`H_2`$ are positively (semi)definite and $`\{H_1,H_2\}=\{\rho ,H\}=0`$ .
The non–linear nature of the evolution could lead to integral curves (presumably even chaotic ) in the configuration space with one–dimensional attractive trajectories (e.g. limit cycles). States that are initially different may evolve into the same final state: attractive trajectories may then be treated as distinct equivalence classes. ‘t Hooft has observed that the dynamics of classical deterministic systems may be mathematically formulated in terms of the unitary evolution operator acting on the Hilbert space of states $`|\psi `$ and expressed in the form of a Schrödinger–like equation. Once again, we stress that this is yet not equivalent to quantization, but it is a pure mathematical construct allowing one to get a statistical inference about the system in question.
To get the lower bound for the Hamiltonian one thus imposes the constraint condition onto the Hilbert space:
$`H_2|\psi =0,`$ (2)
which projects out the states responsible for the negative part of the spectrum. In the deterministic language this means that one gets rid of the unstable trajectories.
What we present here is an explicit realization of ’t Hooft mechanism, with the further step of discovering a connection between the zero point energy and the geometric phase. The model we consider is, however, different from ’t Hooft model. Nevertheless, the Hamiltonian of our model belongs to the same class of the Hamiltonians considered by ’t Hooft, as we explicitly show below.
We start our discussion by considering a system of 1D damped and amplified harmonic oscillators:
$`m\ddot{x}+\gamma \dot{x}+\kappa x=0,`$ (3)
$`m\ddot{y}\gamma \dot{y}+\kappa y=0,`$ (4)
respectively. The $`y`$-oscillator is the time–reversed image of the $`x`$-oscillator. The corresponding Hamiltonian reads
$`H={\displaystyle \frac{1}{m}}p_xp_y+{\displaystyle \frac{1}{2m}}\gamma \left(yp_yxp_x\right)+\left(\kappa {\displaystyle \frac{\gamma ^2}{4m}}\right)xy,`$ (5)
with $`p_x=m\dot{y}\frac{1}{2}\gamma y`$ ; $`p_y=m\dot{x}+\frac{1}{2}\gamma x`$.
In order to show that $`H`$ of Eq.(5) belongs to the class of Hamiltonians above mentioned, we introduce the following notation: $`x_1=(x+y)/\sqrt{2}`$, $`x_2=(xy)/\sqrt{2}`$ and $`p_1=m\dot{x}_1+\frac{1}{2}\gamma x_2`$ ; $`p_2=m\dot{x}_2\frac{1}{2}\gamma x_1`$. We also put $`x_1=r\mathrm{cosh}u`$, $`x_2=r\mathrm{sinh}u`$, and
$`𝒞`$ $`=`$ $`{\displaystyle \frac{1}{4\mathrm{\Omega }m}}\left[\left(p_1^2p_2^2\right)+m^2\mathrm{\Omega }^2\left(x_1^2x_2^2\right)\right]`$ (6)
$`=`$ $`{\displaystyle \frac{1}{4\mathrm{\Omega }m}}\left[p_r^2{\displaystyle \frac{1}{r^2}}p_u^2+m^2\mathrm{\Omega }^2r^2\right],`$ (7)
$`J_2`$ $`=`$ $`{\displaystyle \frac{m}{2}}\left[\left(\dot{x}_1x_2\dot{x}_2x_1\right)\mathrm{\Gamma }r^2\right]={\displaystyle \frac{1}{2}}p_u,`$ (8)
where $`\mathrm{\Gamma }=\gamma /2m`$, $`\mathrm{\Omega }=\sqrt{\frac{1}{m}(\kappa \frac{\gamma ^2}{4m})}`$, with $`\kappa >\frac{\gamma ^2}{4m}`$.
Eq.(5) can be then rewritten as:
$`H`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}p_if_i(q),`$ (9)
with $`f_1(q)=2\mathrm{\Omega }`$, $`f_2(q)=2\mathrm{\Gamma }`$, provided we use the canonical transformation :
$`q_1={\displaystyle \frac{dzm\mathrm{\Omega }}{\sqrt{4J_2^2+4m\mathrm{\Omega }𝒞zm^2\mathrm{\Omega }^2z^2}}},`$ (10)
$`q_2=2u+{\displaystyle \frac{dz}{z}\frac{2J_2}{\sqrt{4J_2^2+4m\mathrm{\Omega }𝒞zm^2\mathrm{\Omega }^2z^2}}},`$ (11)
$`p_1=𝒞,p_2=J_2,`$ (12)
with $`z=r^2`$. One has $`\{q_i,p_i\}=1`$, and the other Poisson brackets vanishing. The above notation reminds one that the structure of the Hamiltonian is the one of $`su(1,1)`$ ($`𝒞`$ is the Casimir operator and $`J_2`$ one of its generators).
Thus from the very definition $`J_2`$ and $`𝒞`$ are beables. Yet also $`q_1`$ and $`q_2`$ are beables (although members of a different Lie sub-algebra than their conjugates) as it can be directly seen from the Hamiltonian (9). The fact that a given dynamics admits more than one Lie-subgroup of beables is by no means generic and here it is only due to a special structure of $`H`$. In such cases it is quite intriguing question which set of beables should be used. We intend here to exploit the outcomes of the $`J_2`$, $`C`$ set.
We put $`H=H__IH_{_{II}}`$, with
$`H__I={\displaystyle \frac{1}{2\mathrm{\Omega }𝒞}}(2\mathrm{\Omega }𝒞\mathrm{\Gamma }J_2)^2,H_{_{II}}={\displaystyle \frac{\mathrm{\Gamma }^2}{2\mathrm{\Omega }𝒞}}J_2^2.`$ (13)
Of course, only nonzero $`r^2`$ should be taken into account in order for $`𝒞`$ to be invertible. Note that $`𝒞`$ is a constant of motion (being the Casimir operator): this ensures that once it has been chosen to be positive, as we do from now on, it will remain such at all times.
We now implement the constraint
$`J_2|\psi =0,`$ (14)
which defines the physical states. Although the system (9), i.e.(5), is deterministic, $`|\psi `$ is not an eigenvector of $`u`$ ($`u`$ does not commute with $`p_u)`$. Of course, if one does not use the operatorial formalism to describe our system, then $`p_u=0`$ implies $`u=\frac{\gamma }{2m}t`$. This point illustrate the specific feature of the description of a classical deterministic system in terms of beables, i.e. of operators commuting at all times: $`u`$ is not a beable.
It should be also mentioned that the constraint (14) automatically implies that the corresponding conjugate variable $`q_2`$ is unbounded in the physical states $`|\psi `$ ( see Eq.(11) ). It should be born in mind that such a behavior of conjugate variables reflects into a typical feature of quantum theory as it assures a compatibility with the Heisenberg uncertainty relation. Eq.(14) implies
$`H|\psi =H__I|\psi =2\mathrm{\Omega }𝒞|\psi =\left({\displaystyle \frac{1}{2m}}p_r^2+{\displaystyle \frac{K}{2}}r^2\right)|\psi ,`$ (15)
where $`Km\mathrm{\Omega }^2`$. $`H__I`$ thus reduces to the Hamiltonian for the linear harmonic oscillator $`\ddot{r}+\mathrm{\Omega }^2r=0`$. The physical states are even with respect to time-reversal ($`|\psi (t)=|\psi (t)`$) and periodical with period $`\tau =\frac{2\pi }{\mathrm{\Omega }}`$.
Let us write the generic state $`|\psi (t)_H`$ as
$`|\psi (t)_H=\widehat{T}\left[\mathrm{exp}\left({\displaystyle \frac{i}{\mathrm{}}}{\displaystyle _{t_0}^t}2\mathrm{\Gamma }J_2𝑑t^{}\right)\right]|\psi (t)_{H__I},`$ (16)
where $`\widehat{T}`$ denotes time-ordering. We remark that for dimensional reasons, in Eq.(16) we need a constant $`\mathrm{}`$, with dimension of an action which from purely classical considerations cannot be fixed in magnitude. Quantum Mechanics tells us how to fix it. Exactly the same situation occurred in the classical statistical mechanics of Gibbs: the precise value to be chosen for the action quantum $`2\pi \mathrm{}`$ was evident only after quantum theory.
The states $`|\psi (t)_H`$ and $`|\psi (t)_{H__I}`$ satisfy the equations:
$`i\mathrm{}{\displaystyle \frac{d}{dt}}|\psi (t)_H`$ $`=`$ $`H|\psi (t)_H,`$ (17)
$`i\mathrm{}{\displaystyle \frac{d}{dt}}|\psi (t)_{H__I}`$ $`=`$ $`2\mathrm{\Omega }𝒞|\psi (t)_{H__I}.`$ (18)
Eq.(18) describes the 2D “isotropic” (or “radial”) harmonic oscillator. $`H__I=2\mathrm{\Omega }𝒞`$ has the spectrum $`__I^n=\mathrm{}\mathrm{\Omega }n`$, $`n=0,\pm 1,\pm 2,\mathrm{}`$. According to our choice for $`𝒞`$ to be positive, only positive values of $`n`$ will be considered.
We can conveniently choose $`t_0=t`$ if we consider the integration along the closed time path, say $`C_t`$, which goes from $`t+iϵ`$ to some arbitrary final time $`t_f`$ and back to $`tiϵ`$ . Since $`J_2`$ is time independent, we can drop the time ordering and write:
$$|\psi (t)_H=\mathrm{exp}\left(i_{C_t}A(t^{})𝑑t^{}\right)|\psi (t)_{H__I},$$
(19)
where $`A(t)\frac{\mathrm{\Gamma }m}{\mathrm{}}(\dot{x}_1x_2\dot{x}_2x_1)`$ and we used the fact that $`_{C_t}r^2=_{t_0}^tr^2+_t^{t_0}r^2=_{t_0}^tr^2_{t_0}^tr^2=0`$. Note that $`(\dot{x}_1x_2\dot{x}_2x_1)dt`$ is the area element in the $`(x_1,x_2)`$ plane enclosed by the trajectories (see Fig.1). Consider the above expression for $`t=\tau `$ and $`t=0`$:
$`|\psi (\tau )_H`$ $`=`$ $`\mathrm{exp}\left(i{\displaystyle _{C_\tau }}A(t^{})𝑑t^{}\right)|\psi (\tau )_{H__I},`$ (20)
$`|\psi (0)_H`$ $`=`$ $`\mathrm{exp}\left(i{\displaystyle _{C_0}}A(t^{})𝑑t^{}\right)|\psi (0)_{H__I},`$ (21)
where $`C_\tau `$ and $`C_0`$ are the time contours going from $`\tau `$ (or from $`0`$) to $`t_f`$ and back along the real line. We focus now on eigenstates of $`H`$ and $`H_I`$. We then get
$`{}_{H}{}^{}\psi (\tau )|\psi (0)_{H}^{}`$ (22)
$`=_{H__I}\psi (0)|\mathrm{exp}\left(i{\displaystyle _{C_{0\tau }}}A(t^{})𝑑t^{}\right)|\psi (0)_{H__I}`$ (23)
$`e^{i\varphi },`$ (24)
where the contour $`C_{0\tau }`$ is the one going from $`t^{}=0`$ to $`t^{}=\tau `$ and back. Notice that the dependence on the (arbitrary) final time $`t_f`$ has disappeared. One may observe that the evolution (or dynamical) part of the phase does not enter in $`\varphi `$, as the integral in Eq.(24) picks up a purely geometric contribution (Berry–Anandan–like phase). The integral in (24) can be calculated by rewriting it as a contour integral in the complex plane. If $`x_1`$ and $`x_2`$ are analytically continued into imaginary $`u`$, we may define $`z=x_1+x_2`$ and $`\overline{z}=x_1x_2`$. The hyperbola is thus mapped onto a circle in the Gaussian plane. We then obtain
$`{\displaystyle _{C_{0\tau }}}A(t^{})𝑑t^{}={\displaystyle \frac{\mathrm{\Gamma }m}{\mathrm{}}}\text{Im}\left[{\displaystyle _\mathrm{\Delta }}𝑑z\overline{z}\right]=\pi R^2{\displaystyle \frac{\gamma }{\mathrm{}}},`$ (25)
where $`\mathrm{\Delta }`$ is the clockwise contour in the complex plane with radius $`R`$ ($`z\overline{z}=R^2`$): this is equal to the hyperbolic radius for the envelope of the trajectories in the plane $`(x_1,x_2)`$. The value of $`R`$ is obtained by extremizing $`r(t)`$. We get
$`R=\sqrt{{\displaystyle \frac{2}{m}}}{\displaystyle \frac{1}{\mathrm{\Omega }}},`$ (26)
where $``$ is the initial energy given by $`=\frac{1}{2}mv_0^2+\frac{1}{2}m\mathrm{\Omega }^2r_0^2`$ (we use $`r(t)=r_0\mathrm{cos}\mathrm{\Omega }t+\frac{v_0}{\mathrm{\Omega }}\mathrm{sin}\mathrm{\Omega }t`$). Note that $`=0`$ is not allowed since it corresponds to the system confined to the origin, which we exclude. In conclusion, we get $`\varphi =\alpha \pi `$, with dimensionless $`\alpha R^2\frac{\gamma }{\mathrm{}}`$.
Because the physical states $`|\psi `$ are periodic ones, let us focus our attention on those. Following , one may generally write
$`|\psi (\tau )`$ $`=`$ $`e^{i\varphi \frac{i}{\mathrm{}}_0^\tau \psi (t)|H|\psi (t)𝑑t}|\psi (0)`$ (27)
$`=`$ $`e^{i2\pi n}|\psi (0),`$ (28)
i.e. $`\frac{\psi (\tau )|H|\psi (\tau )}{\mathrm{}}\tau \varphi =2\pi n`$, $`n=0,1,2,\mathrm{}`$, which by using $`\tau =\frac{2\pi }{\mathrm{\Omega }}`$ and $`\varphi =\alpha \pi `$, gives
$`_{{}_{I}{}^{},eff}^n\psi _n(\tau )|H|\psi _n(\tau )=\mathrm{}\mathrm{\Omega }\left(n+{\displaystyle \frac{\alpha }{2}}\right),`$ (29)
where the index $`n`$ has been introduced to exhibit the $`n`$ dependence of the state and the corresponding energy. $`_{{}_{I}{}^{},eff}^n`$ gives the effective $`n`$th energy level of the physical system, namely the energy given by $`__I^n`$ corrected by its interaction with the environment. We thus see that the dissipation term $`J_2`$ of the Hamiltonian is actually responsible for the “zero point energy” ($`n=0`$): $`E_0=\frac{\mathrm{}}{2}\mathrm{\Omega }\alpha `$.
As well known, the zero point energy is the “signature” of quantization since in Quantum Mechanics it is formally due to the non-zero commutator of the canonically conjugate $`q`$ and $`p`$ operators. Thus dissipation manifests itself as “quantization”. In other words, $`E_0`$, which appears as the “quantum contribution” to the spectrum of the conservative evolution of physical states, signals the underlying dissipative dynamics. If we really want to match the Quantum Mechanics zero point energy, we have to resort to experiment, which fixes $`\alpha =1`$, with $`\mathrm{}`$ being the Planck constant. In turn this fully exhibits the geometric nature of the phase $`\varphi `$ and gives $`R_0=\frac{\mathrm{}}{\gamma }`$. From Eq.(26) we then see that $`_0=\frac{1}{2}m\mathrm{\Omega }^2R_0^2`$ is the energy corresponding to $`E_0`$. We thus put $`_0=E_0`$, which gives $`\mathrm{\Omega }=\frac{\gamma }{m}`$, i.e. $`\kappa =5\frac{\gamma ^2}{4m}`$, consistent with the reality condition for $`\mathrm{\Omega }`$.
Notice that the only free parameter of the theory is the ratio $`\frac{\kappa }{m}`$. We also remark that in ref. the phase integral in Eq.(24), which plays there the rôle of dissipative interference phase, has proved to be in fact always non-zero in order to have quantum mechanical interference in the electron double slit experiment. Equivalently, in the present formalism we must always have $`x_1x_2`$ (i.e. $`r0`$) in order to have $`A(t)0`$ (cf. Eq.(19)) and thus a non-zero geometric phase.
In order to better understand the dynamical rôle of $`J_2`$ we rewrite Eq.(16) as follows
$`|\psi (t)_H=\widehat{T}\left[\mathrm{exp}\left(i{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle _{u(t_0)}^{u(t)}}2J_2𝑑u^{}\right)\right]|\psi (t)_{H__I},`$ (30)
by using $`u(t)=\mathrm{\Gamma }t`$. Accordingly, we have
$`i\mathrm{}{\displaystyle \frac{}{u}}|\psi (t)_H=2J_2|\psi (t)_H.`$ (31)
We thus see that $`2J_2`$ is responsible for shifts (translations) in the $`u`$ variable, as is to be expected since $`2J_2=p_u`$ (cf. Eq.(8)). In operatorial notation we can write indeed $`p_u=i\mathrm{}\frac{}{u}`$. Then, in full generality, Eq.(14) defines families of physical states, representing stable, periodic trajectories (cf. Eq.(15)). $`2J_2`$ implements transition from family to family, according to Eq.(31). Eq.(17) can be then rewritten as
$`i\mathrm{}{\displaystyle \frac{d}{dt}}|\psi (t)_H=i\mathrm{}{\displaystyle \frac{}{t}}|\psi (t)_H+i\mathrm{}{\displaystyle \frac{du}{dt}}{\displaystyle \frac{}{u}}|\psi (t)_H,`$ (32)
where the first term on the r.h.s. denotes of course derivative with respect to the explicit time dependence of the state. The dissipation contribution to the energy is thus described by the “translations” in the $`u`$ variable. It is then interesting to consider the derivative
$`{\displaystyle \frac{S}{U}}={\displaystyle \frac{1}{T}}.`$ (33)
¿From Eq.(9), by using $`S\frac{2J_2}{\mathrm{}}`$ and $`U2\mathrm{\Omega }𝒞`$, we obtain $`T=\mathrm{}\mathrm{\Gamma }`$. Eq. (33) is the defining relation for temperature in thermodynamics (with $`k_B=1`$) so that one could formally regard $`\mathrm{}\mathrm{\Gamma }`$ (which dimensionally is an energy) as the temperature, provided the dimensionless quantity $`S`$ is identified with the entropy. In such a case, the “full Hamiltonian” Eq.(9) plays the rôle of the free energy $``$: $`H=2\mathrm{\Omega }𝒞(\mathrm{}\mathrm{\Gamma })\frac{2J_2}{\mathrm{}}=UTS=`$. Thus $`2\mathrm{\Gamma }J_2`$ represents the heat contribution in $`H`$ (or $``$). Of course, consistently, $`\frac{}{T}|_\mathrm{\Omega }=\frac{2J_2}{\mathrm{}}`$. In conclusion $`\frac{2J_2}{\mathrm{}}`$ behaves as the entropy, which is not surprising since it controls the dissipative (thus irreversible) part of the dynamics.
We can also take the derivative of $``$ (keeping $`T`$ fixed) with respect to $`\mathrm{\Omega }`$. We then have
$`{\displaystyle \frac{}{\mathrm{\Omega }}}|_T={\displaystyle \frac{U}{\mathrm{\Omega }}}|_T=mr^2\mathrm{\Omega },`$ (34)
which is the angular momentum: this is to be expected since it is the conjugate variable of the angular velocity $`\mathrm{\Omega }`$. It is also suggestive that the temperature $`\mathrm{}\mathrm{\Gamma }`$ is actually given by the background zero point energy: $`\mathrm{}\mathrm{\Gamma }=\frac{\mathrm{}\mathrm{\Omega }}{2}`$.
In the light of the above results, the condition (14) can be then interpreted as a condition for an adiabatic physical system. $`\frac{2J_2}{\mathrm{}}`$ might be viewed as an analogue of the Kolmogorov–Sinai entropy for chaotic dynamical systems.
Finally, we remark that the thermodynamical picture above outlined is also consistent with the results on the canonical quantization of dissipative system in quantum field theory (QFT) presented in ref. . The transitions between unitarily inequivalent representations of the canonical commutation relations induced by the entropy operator in QFT correspond to the transitions between families of stable trajectories induced by the entropy $`\frac{2J_2}{\mathrm{}}`$ in the present paper.
G.V. is very grateful to Alan Widom for long, illuminating discussions and P.J. wishes to thank John C. Taylor for many useful comments.
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# 1 Introduction
## 1 Introduction
This report is devoted to well-known quantum mechanical problem to find the simultaneous eigenfunction of commuting set of Hamiltonians for the periodic Toda chain. The first important step in this direction has been done by Gutzwiller who solved the problem for the particular cases $`N=2,3`$ and 4 particles and found such important phenomena as quantization of spectrum and separation of multidimensional Baxter equation into the product of one dimension ones. In fact, he performed the quantization of periodic Toda chain in terms of separated variables introduced by Flaschka and McLaughlin . Next important step was taken by Sklyanin who constructed $`R`$-matrix formalism for both classical and quantum cases Toda chains and introduced the algebraic method of separation variables for an arbitrary number of particles. His approach drastically simplifies the derivation of the Baxter equation and works for wide spectrum of integrable models .
Our method to solve the spectral problem consists of analytical re-interpretation of Sklyanin’s algebraic ideas which allows to find the integral representation for the eigenfunctions of the periodic Toda chain as a kind of generalized Fourier transform with the eigenfunctions for the open Toda chain . In turn, this method can be treated as a natural generalization of original Gutzwiller’s approach. The explicit solution for the eigenfunctions of the open Toda chain plays a key role in this construction.
It has been discovered by Kostant that the commuting set of Hamiltonians of an open Toda chain coincides with the Whittaker model of the center of universal enveloping algebra. Hence, the Whittaker functions are in fact the eigenfunctions for the open Toda chain. In the usual group-theoretical way the Whittaker function is defined as a matrix element between compact and Whittaker vectors - in the principal series representation. There are many obstacles to generalize this approach to other quantum models or to loop groups.
The present approach to construct the eigenfunctions for both periodic and open chains is rather different : it bases on the Quantum Inverse Scattering Method for the periodic Toda chain . One of the interesting results of analytical calculations in the $`R`$-matrix framework is the revealing of recurrent relation between $`N`$ and $`N1`$ particles eigenfunctions for the open Toda chain (in fact the idea to use a recurrent relation was pointed out by Sklyanin in ; our recurrent relations is an explicit realization of such an idea). This naturally leads to new integral representation for the Weyl invariant Whittaker functions to compare with classical results -. This representation is quite explicit and very useful to investigate the different asymptotics. In particular, the Gindikin-Karpelevich formula for the Harish-Chandra function can be obtained in a very simple way for particular case of $`GL(N,)`$ group. The eigenfunction for the periodic Toda chain are constructed in a rather explicit form and have essentially the same form as the recurrent relation mentioned above. The integral formula for eigenfunctions can be considered as a representation of the Whittaker functions for $`\widehat{GL}(N)`$ group at the critical level.
The present approach can be generalized to other quantum integrable models. For example, the eigenfunctions for the relativistic Toda chain are calculated in using the same QISM ideology.
## 2 Quantum Toda chain: description of the model
### 2.1 Periodic spectral problem
The quantum $`N`$-periodic Toda chain is a multi-dimensional eigenvalue problem with $`N`$ mutually commuting Hamiltonians $`H_k(x_1,p_1;\mathrm{};x_N,p_N),(k=1,\mathrm{},N),`$ where the simplest Hamiltonians have the form
$$\begin{array}{c}H_1=_{k=1}^Np_k\\ H_2=_{k<m}p_kp_m_{k=1}^Ne^{x_kx_{k+1}}\\ H_3=_{k<m<n}p_kp_mp_n+\mathrm{},\end{array}$$
(2.1)
$`(x_{N+1}x_1)`$ etc. and the phase variables $`x_k,p_k`$ satisfy the standard commutation relations $`[x_k,p_m]=i\mathrm{}\delta _{km}`$. The main goal is to find the solution to the eigenvalue problem
$$\begin{array}{c}H_k\mathrm{\Psi }_E=E_k\mathrm{\Psi }_Ek=1,\mathrm{},N\end{array}$$
(2.2)
with fast decreasing wave function $`\mathrm{\Psi }_E`$. To be more precise, let us note that, due to translation invariance, the solution to (2.2) has the following structure:
$$\begin{array}{c}\mathrm{\Psi }_E(x_1,\mathrm{},x_N)=\stackrel{~}{\mathrm{\Psi }}_E(x_1x_2,\mathrm{},x_{N1}x_N)\mathrm{exp}\left\{\frac{i}{\mathrm{}}E_1_{k=1}^Nx_k\right\}\end{array}$$
(2.3)
One needs to find the solution to (4.11) such that $`\stackrel{~}{\mathrm{\Psi }}_EL^2(^{N1})`$. In equivalent terms, we impose the requirement
$$\begin{array}{c}f(E_1)\mathrm{\Psi }_E(x_1,\mathrm{},x_N)𝑑E_1L^2(^N)\end{array}$$
(2.4)
for any smooth function $`f(y),(y)`$ with finite support.
### 2.2 $`GL(N1,)`$ spectral problem
It turns out that solution to (2.2), (2.4) can be effectively written in terms of the wave functions corresponding to open $`N1`$-particle Toda chain (quantum $`GL(N1,)`$ chain). The Hamiltonians of the latter system can be (formally) derived from (2.1) by cancelling out all the operators containing $`p_N`$ and $`x_N`$ thus obtaining exactly $`N1`$ commuting Hamiltonians $`h_k(x_1,p_1;\mathrm{};x_{N1},p_{N1})`$ ($`k=1,\mathrm{},N1`$). Let $`\gamma =(\gamma _1,\mathrm{},\gamma _{N1})^{N1}`$, $`x=(x_1,\mathrm{},x_{N1})^{N1}`$. We consider $`GL(N1,)`$ spectral problem
$$\begin{array}{c}h_k\psi _\gamma (x)=\sigma _k(\gamma )\psi _\gamma (x)k=1,\mathrm{},N1\end{array}$$
(2.5)
where $`\sigma _k(\gamma )`$ are elementary symmetric functions.
Obviously, in the asymptotic region $`x_{k+1}x_k,(k=1,\mathrm{},N2)`$ all potentials vanish and the solution to (2.5) is a superposition of plane waves. The problem is to find a solution to (2.5) satisfying the following properties:
* The solution vanishes very rapidly
$$\begin{array}{c}\psi _\gamma (x)\mathrm{exp}\left\{\frac{2}{\mathrm{}}e^{(x_kx_{k+1})/2}\right\}x_kx_{k+1}\mathrm{}\end{array}$$
(2.6)
* The function $`\psi _\gamma `$ is Weyl-invariant, i.e. it is symmetric under any permutation
$$\begin{array}{c}\psi _{\mathrm{}\gamma _j\mathrm{}\gamma _k\mathrm{}}=\psi _{\mathrm{}\gamma _k\mathrm{}\gamma _j\mathrm{}}\end{array}$$
(2.7)
* $`\psi _\gamma `$ can be analytically continued to an entire function of $`\gamma ^{N1}`$ and the following asymptotics hold:
$$\begin{array}{c}\psi _\gamma |\gamma _j|^{\frac{2N}{2}}\mathrm{exp}\left\{\frac{\pi }{2\mathrm{}}(N2)|\gamma _j|\right\}\end{array}$$
(2.8)
as $`|\mathrm{Re}\gamma _j|\mathrm{}`$ in the finite strip of complex plane.
The properties (i)-(iii) define a unique solution to the spectral problem (2.2).
## 3 Main results
###### Theorem 3.1
The following statements hold :
* Let a set $`\gamma _{jk}`$ be the lower triangular $`(N1)\times (N1)`$ matrix. The solution to the spectral problem (2.5)-(2.8) can be written in the form of multiple Mellin-Barnes integrals <sup>4</sup><sup>4</sup>4We identify the set $`\gamma `$ with the last row $`(\gamma _{N1,1},\mathrm{},\gamma _{N1,N1})`$. :
$$\begin{array}{c}\psi _{\gamma _{_{N1,1}},\mathrm{},\gamma _{_{N1,N1}}}(x_1,\mathrm{},x_{N1})=\\ =\frac{(2\pi \mathrm{})^{\frac{(N1)(N2)}{2}}}{_{k=1}^{N2}k!}\underset{𝒞}{}_{n=1}^{N2}\frac{_{j=1}^n_{k=1}^{n+1}\mathrm{}^{\frac{\gamma _{nj}\gamma _{n+1,k}}{i\mathrm{}}}\mathrm{\Gamma }\left(\frac{\gamma _{nj}\gamma _{n+1,k}}{i\mathrm{}}\right)}{\underset{\stackrel{j,k=1}{j<k}}{\overset{n}{}}\mathrm{\Gamma }\left(\frac{\gamma _{nj}\gamma _{nk}}{i\mathrm{}}\right)\mathrm{\Gamma }\left(\frac{\gamma _{nk}\gamma _{nj}}{i\mathrm{}}\right)}\times \\ \mathrm{exp}\left\{\frac{i}{\mathrm{}}_{n,k=1}^{N1}x_n\left(\gamma _{nk}\gamma _{n1,k}\right)\right\}\underset{\stackrel{j,k=1}{jk}}{\overset{N2}{}}d\gamma _{jk}\end{array}$$
(3.1)
where the integral should be understand as follows: first we integrate on $`\gamma _{11}`$ over the line $`\mathrm{Im}\gamma _{11}>\mathrm{max}\{\mathrm{Im}\gamma _{21},\mathrm{Im}\gamma _{22}\}`$; then we integrate on the set $`(\gamma _{21},\gamma _{22})`$ over the lines $`\mathrm{Im}\gamma _{2j}>\mathrm{max}_m\{\mathrm{Im}\gamma _{3m}\}`$ and so on. The last integrations should be performed on the set of variables $`(\gamma _{N2,1}\mathrm{},\gamma _{N2,N2})`$ over the lines $`\mathrm{Im}\gamma _{N2,k}>\mathrm{max}_m\{\mathrm{Im}\gamma _{N1,m}\}`$.
* In the region $`x_kx_{k+1}(k=1,\mathrm{},N1)`$ the solution has the following asymptotics:
$$\begin{array}{c}\psi _\gamma (x)=_{sW}\varphi (s\gamma )e^{\frac{i}{\mathrm{}}(s\gamma ,x)}+O\left(\text{max}\left\{e^{x_kx_{k+1}}\right\}_{k=1}^{N1}\right)\end{array}$$
(3.2)
where $`(.,.)`$ is a scalar product in $`^{N1}`$ and the summation is performed over the permutation group; $`\varphi (\gamma )`$ is (renormalized) Harish-Chandra function
$$\begin{array}{c}\varphi (\gamma )=\mathrm{}^{2i(\gamma ,\rho )/\mathrm{}}_{j<k}\mathrm{\Gamma }\left(\frac{\gamma _j\gamma _k}{i\mathrm{}}\right)\end{array}$$
(3.3)
where $`(\gamma ,\rho )\frac{1}{2}\underset{m=1}{\overset{N1}{}}(N2m)\gamma _k`$.
* The functions (3.1) have the scalar product
$$\begin{array}{c}\underset{\text{}^{N1}}{}\overline{\psi }_\gamma ^{}(x)\psi _\gamma (x)𝑑x=\frac{\mu ^1(\gamma )}{(N1)!}_{sW}\delta (s\gamma \gamma ^{})(\gamma ,\gamma ^{}^{N1})\end{array}$$
(3.4)
and obey the completeness condition
$$\begin{array}{c}\underset{\text{}^{N1}}{}\mu (\gamma )\psi _\gamma (x)\overline{\psi }_\gamma (y)𝑑\gamma =\delta (xy)\end{array}$$
(3.5)
where
$$\begin{array}{c}\mu (\gamma )=\frac{(2\pi \mathrm{})^{1N}}{(N1)!}_{j<k}\left|\mathrm{\Gamma }\left(\frac{\gamma _j\gamma _k}{i\mathrm{}}\right)\right|^2\end{array}$$
(3.6)
is the Sklyanin measure .
The eigenfunctions for the periodic chain are constructed as a kind of Fourier transform with the function (3.1). Let
$$\begin{array}{c}t_N(\lambda ;E)=\underset{k=0}{\overset{N}{}}(1)^k\lambda ^{Nk}E_k\end{array}$$
(3.7)
and $`e_j`$ denotes $`j`$-th basis vector in $`^{N1}`$.
###### Theorem 3.2
The solution to the spectral problem (2.2), (2.4) can be represented as the integral over real variables $`\gamma =(\gamma _1,\mathrm{},\gamma _{N1})`$ in the following form:
$$\begin{array}{c}\mathrm{\Psi }_E(x,x_N)=\frac{1}{2\pi }\underset{\text{}^{N1}}{}\mu (\gamma )C(\gamma ;E)\mathrm{\Psi }_{\gamma ,E_1}(x,x_N)𝑑\gamma \end{array}$$
(3.8)
where
* The function $`\mathrm{\Psi }_{\gamma ,E_1}(x,x_N)`$ is defined in terms of solution (3.1) to the $`GL(N1,)`$ spectral problem:
$$\begin{array}{c}\mathrm{\Psi }_{\gamma ,E_1}(x,x_N)=\psi _\gamma (x)\mathrm{exp}\left\{\frac{i}{\mathrm{}}\left(E_1_{m=1}^{N1}\gamma _m\right)x_N\right\}\end{array}$$
(3.9)
* The function $`C(\gamma ;E)`$ is the solution of multi-dimensional Baxter equations
$$\begin{array}{c}t_N(\gamma _j;E)C(\gamma ;E)=i^NC(\gamma +i\mathrm{}e_j;E)+i^NC(\gamma i\mathrm{}e_j;E)\end{array}$$
(3.10)
which is symmetric entire function in $`\gamma `$-variables with the asymptotics
$$\begin{array}{c}C(\gamma ;E)|\gamma _k|^{N/2}\mathrm{exp}\left\{\frac{\pi N|\gamma _k|}{2\mathrm{}}\right\}\end{array}$$
(3.11)
as $`\mathrm{Re}\gamma _k\pm \mathrm{}`$ in the strip $`|\mathrm{Im}\gamma _k|\mathrm{}`$
The above restrictions imposed on solution to (3.10) are reformulation of the quantization condition (2.4) on the level of $`\gamma `$-representation. To obtain the explicit integral form for the eigenfunctions, we use the solution to (3.10), (3.11) in the Pasquier-Gaudin form (see sect.7 below)
$$\begin{array}{c}C(\gamma ;E)=_{j=1}^{N1}\frac{c_+(\gamma _j;E)\xi (E)c_{}(\gamma _j;E)}{\underset{k=1}{\overset{N}{}}\mathrm{sinh}\frac{\pi }{\mathrm{}}\left(\gamma _j\delta _k(E)\right)}\end{array}$$
(3.12)
where the entire functions $`c_\pm (\gamma )`$ are two Gutzwiller’s solutions of the one-dimensional Baxter equation
$$\begin{array}{c}t(\gamma ;E)c(\gamma ;E)=i^Nc(\gamma +i\mathrm{};E)+i^Nc(\gamma i\mathrm{};E)\end{array}$$
(3.13)
and the parameters $`\xi (E),\delta =(\delta _1(E),\mathrm{},\delta _N(E))`$ satisfy the Gutzwiller conditions (the energy quantization) (see sect.7 below). Then the multiple integral (3.8) can be explicitly evaluated. Let $`y=(y_1,\mathrm{},y_N)^N`$ be an arbitrary vector. We denote $`y^{(s)}(y_1,\mathrm{},y_{s1},y_{s+1},\mathrm{},y_N)`$ the corresponding vector in $`^{N1}`$.
###### Theorem 3.3
Assuming that $`\delta _j(E)\delta _k(E)`$, the solution (3.8) can be written (up to an inessential numerical factor) in the equivalent form
$$\begin{array}{c}\mathrm{\Psi }_E(x,x_N)=_{s=1}^N(1)^{Ns}_{n^{(s)}\text{}^{N1}}\mathrm{\Delta }(\delta ^{(s)}+i\mathrm{}n^{(s)})C_+(\delta ^{(s)}+i\mathrm{}n^{(s)})\mathrm{\Psi }_{\delta ^{(s)}+i\mathrm{}n^{(s)},E_1}(x,x_N)\end{array}$$
(3.14)
where
$$\begin{array}{c}C_+(\gamma )_{j=1}^{N1}c_+(\gamma _j;E)\end{array}$$
(3.15)
and $`\mathrm{\Delta }(\gamma )=\underset{j>k}{}(\gamma _j\gamma _k)`$ is the Vandermonde determinant.
###### Remark 3.1
For $`N=2,3`$ and $`N=4`$ formula (3.14) reproduces the results obtained by Gutzwiller .
## 4 $`R`$-matrix approach
The Toda chain can be nicely described using the $`R`$-matrix approach . It is well known that the Lax operator
$$\begin{array}{c}L_n(\lambda )=\left(\begin{array}{cc}\lambda p_n& e^{x_n}\\ e^{x_n}& 0\end{array}\right)\end{array}$$
(4.1)
satisfies the commutation relations
$$\begin{array}{c}R(\lambda \mu )(L_n(\lambda ))I)(IL_n(\mu ))=(IL_n(\mu ))(L_n(\lambda )I)R(\lambda \mu )\end{array}$$
(4.2)
where
$$\begin{array}{c}R(\lambda )=II+\frac{i\mathrm{}}{\lambda }P\end{array}$$
(4.3)
is a rational $`R`$-matrix. The monodromy matrix
$$\begin{array}{c}T__N(\lambda )\stackrel{\text{def}}{=}L_N(\lambda )\mathrm{}L_1(\lambda )\left(\begin{array}{cc}A_N(\lambda )& B_N(\lambda )\\ C_N(\lambda )& D_N(\lambda )\end{array}\right)\end{array}$$
(4.4)
satisfies the analogous equation
$$\begin{array}{c}R(\lambda \mu )(T(\lambda )I)(IT(\mu ))=(IT(\mu ))(T(\lambda )I)R(\lambda \mu )\end{array}$$
(4.5)
In particular, the following commutation relations hold:
$$\begin{array}{c}[A_N(\lambda ),A_N(\mu )]=[C_N(\lambda ),C_N(\mu )]=0\end{array}$$
(4.6)
$$\begin{array}{c}(\lambda \mu +i\mathrm{})A_N(\mu )C_N(\lambda )=(\lambda \mu )C_N(\lambda )A_N(\mu )+i\mathrm{}A_N(\lambda )C_N(\mu )\end{array}$$
(4.7)
$$\begin{array}{c}(\lambda \mu +i\mathrm{})D_N(\lambda )C_N(\mu )=(\lambda \mu )C_N(\mu )D_N(\lambda )+i\mathrm{}D_N(\mu )C_N(\lambda )\end{array}$$
(4.8)
From (4.5) it can be easily shown that the trace of the monodromy matrix
$$\begin{array}{c}\widehat{t}_N(\lambda )=A_N(\lambda )+D_N(\lambda )\end{array}$$
(4.9)
satisfies the commutation relations $`[\widehat{t}(\lambda ),\widehat{t}(\mu )]=0`$ and is a generating function for the Hamiltonians of the periodic Toda chain:
$$\begin{array}{c}\widehat{t}_N(\lambda )=_{k=0}^N(1)^k\lambda ^{Nk}H_k\end{array}$$
(4.10)
We reformulate the spectral equations (2.2) as follows:
$$\begin{array}{c}\widehat{t}_N(\lambda )\mathrm{\Psi }_E=t_N(\lambda ;E)\mathrm{\Psi }_E\end{array}$$
(4.11)
where
$$\begin{array}{c}t_N(\lambda ;E)=\underset{k=0}{\overset{N}{}}(1)^k\lambda ^{Nk}E_k\end{array}$$
(4.12)
On the other hand, it can be easily shown that the operator
$$\begin{array}{c}A_{N1}(\lambda )_{k=0}^{N1}(1)^k\lambda ^{Nk1}h_k(x_1,p_1;\mathrm{};x_{N1},p_{N1})\end{array}$$
(4.13)
is nothing but the generating function for the Hamiltonians $`h_k`$ of $`GL(N1)`$ Toda chain. Therefore, the $`GL(N1,)`$ spectral equations can be written in the form
$$\begin{array}{c}A_{N1}(\lambda )\psi _\gamma (x)=_{m=1}^{N1}(\lambda \gamma _m)\psi _\gamma (x)\end{array}$$
(4.14)
Using the obvious relation
$$\begin{array}{c}C_N(\lambda )=e^{x_N}A_{N1}(\lambda )\end{array}$$
(4.15)
one obtains, as a trivial corollary of (4.14),
$$\begin{array}{c}C_N(\gamma _j)\psi _\gamma (x)=0\gamma _j\gamma \end{array}$$
(4.16)
###### Remark 4.1
Equations (4.16) are an analytical analog of the notion of ”operator zeros” introduced by Sklyanin .
## 5 Eigenfunctions for the open Toda chain
Suppose that solution to (4.14) satisfying (2.6)-(2.8) is given. Using the commutation relations (4.7), (4.8) together with (4.16), it is easy to show that the following relations hold
$$\begin{array}{c}A_N(\gamma _j)\psi _\gamma =i^Ne^{x_N}\psi _{\gamma i\mathrm{}e_j}\end{array}$$
(5.1a)
$$\begin{array}{c}D_N(\gamma _j)\psi _\gamma =i^Ne^{x_N}\psi _{\gamma +i\mathrm{}e_j}\end{array}$$
(5.1b)
($`j=1,\mathrm{},N1`$) where $`e_j`$ is $`j`$-th basis vector in $`^{N1}`$. Note that (5.1b) is a corollary of (5.1a) since the quantum determinant of the monodromy matrix (4.4) is unity.
Let us introduce the key object - the auxiliary function
$$\begin{array}{c}\mathrm{\Psi }_{\gamma ,ϵ}(x_1,\mathrm{},x_N)\stackrel{\text{def}}{=}\psi _\gamma (x)\mathrm{exp}\left\{\frac{i}{\mathrm{}}\left(ϵ\underset{m=1}{\overset{N1}{}}\gamma _m\right)x_N\right\}\end{array}$$
(5.2)
where $`ϵ`$ is an arbitrary parameter. From (4.14), (4.15) and (5) it is readily seen that this function satisfies to equations
$$\begin{array}{c}C_N(\lambda )\mathrm{\Psi }_{\gamma ,ϵ}=e^{x_N}_{j=1}^{N1}(\lambda \gamma _j)\mathrm{\Psi }_{\gamma ,ϵ}\end{array}$$
(5.3a)
$$\begin{array}{c}A_N(\lambda )\mathrm{\Psi }_{\gamma ,ϵ}=\left(\lambda ϵ+_{m=1}^{N1}\gamma _m\right)_{j=1}^{N1}(\lambda \gamma _j)\mathrm{\Psi }_{\gamma ,ϵ}+i^N_{j=1}^{N1}\mathrm{\Psi }_{\gamma i\mathrm{}e_j,ϵ}_{mj}\frac{\lambda \gamma _m}{\gamma _j\gamma _m}\end{array}$$
(5.3b)
$$\begin{array}{c}D_N(\lambda )\mathrm{\Psi }_{\gamma ,ϵ}=i^N_{j=1}^{N1}\mathrm{\Psi }_{\gamma +i\mathrm{}e_j,ϵ}_{mj}\frac{\lambda \gamma _m}{\gamma _j\gamma _m}\end{array}$$
(5.3c)
In particular,
$$\begin{array}{c}\widehat{t}_N(\gamma _j)\mathrm{\Psi }_{\gamma ,ϵ}=i^N\mathrm{\Psi }_{\gamma ,ϵ}+i^N\mathrm{\Psi }_{\gamma ,ϵ}\end{array}$$
(5.4)
The problem is to find the corresponding solution for $`GL(N,)`$ Toda chain using the above information, i.e. in terms of the function $`\mathrm{\Psi }_{\gamma ,ϵ}(x)`$ construct the Weyl invariant function $`\psi _{\lambda _1,\mathrm{},\lambda _N}(x_1,\mathrm{},x_N)`$ satisfying to equations
$$\begin{array}{c}A_N(\lambda )\psi _{\lambda _1,\mathrm{},\lambda _N}=_{k=1}^N(\lambda \lambda _k)\psi _{\lambda _1,\mathrm{},\lambda _N}\end{array}$$
(5.5a)
$$\begin{array}{c}A_{N+1}(\lambda _n)\psi _{\lambda _1,\mathrm{},\lambda _N}=i^{N1}e^{x_{N+1}}\psi _{\lambda _1,\mathrm{},\lambda _ni\mathrm{},\mathrm{},\lambda _N}(n=1,\mathrm{},N)\end{array}$$
(5.5b)
and obeying the similar to (2.6)-(2.8) conditions
###### Lemma 5.1
Let $`\mathrm{\Psi }_{\gamma ,ϵ}(x,x_N)`$ be the auxiliary function (5.2). Let $`\lambda =(\lambda _1,\mathrm{},\lambda _N)^N`$ be the set of indeterminates. Let
$$\begin{array}{c}\mu (\gamma )=\frac{(2\pi \mathrm{})^{1N}}{(N1)!}_{j<k}\left\{\mathrm{\Gamma }\left(\frac{\gamma _j\gamma _k}{i\mathrm{}}\right)\mathrm{\Gamma }\left(\frac{\gamma _k\gamma _j}{i\mathrm{}}\right)\right\}^1\end{array}$$
(5.6)
$$\begin{array}{c}Q(\gamma _1,\mathrm{},\gamma _{N1}|\lambda _1,\mathrm{},\lambda _N)=_{j=1}^{N1}_{k=1}^Nh^{\frac{\gamma _j\lambda _k}{i\mathrm{}}}\mathrm{\Gamma }\left(\frac{\gamma _j\lambda _k}{i\mathrm{}}\right)\end{array}$$
(5.7)
Then the Weyl invariant solution to the spectral problem (5.5a)-(5.5b) with the properties discussed above is given by recurrent formula
$$\begin{array}{c}\psi _{\lambda _1,\mathrm{},\lambda _N}(x_1,\mathrm{},x_N)=\underset{𝒞}{}\mu (\gamma )Q(\gamma ;\lambda )\mathrm{\Psi }_{\gamma ;\lambda _1+\mathrm{}+\lambda _N}(x_1,\mathrm{},x_N)𝑑\gamma \end{array}$$
(5.8)
where the integration is performed along the horizontal lines with $`\text{Im}\gamma _j>\text{max}_k\{\text{Im}\lambda _k\}`$.
Proof. One needs to calculate the action of the operators $`A_N(\lambda )`$ and
$$\begin{array}{c}A_{N+1}(\lambda )=(\lambda p_{N+1})A_N(\lambda )+e^{x_{N+1}}C_N(\lambda )\end{array}$$
(5.9)
on the function (5.8) using the formulae (5.3b) and (5.5a), (5.3a). The shifted contours can be deformed to original ones using the facts that integrand in (5.8) is an entire function fast decreasing in any finite horizontal strip of complex plane as $`|\mathrm{Re}\gamma _j|\mathrm{}`$. The last step is to use the difference equations for the parts of the integrand with respect to the shifts $`\pm i\mathrm{}`$ of parameters $`\gamma _m`$ and $`\lambda _k`$.
Proof of Theorem 3.1. The proof of (3.1) is straightforward resolution of the recurrent relations (5.8) starting with trivial eigenfunction $`\psi _{\gamma _{_{11}}}(x_1)=\mathrm{exp}\{\frac{i}{\mathrm{}}\gamma _{11}x_1\}`$. Obviously, the function (3.1) is symmetric under the permutation of parameters $`\gamma `$. The asymptotics (2.6) can be proved using the steepest descent method. Using the Stirling formula for the $`\mathrm{\Gamma }`$-functions as $`\gamma _{N1,k}\gamma _k\pm \mathrm{}`$, it is easy to see that the asymptotics (2.8) hold. Hence, (3.1) is an appropriate solution to the spectral problem.
Further, the formula (3.2) can be proved as follows. The integrand in (5.8) decreases exponentially as $`|\gamma _j|\mathrm{},(j=1,\mathrm{},N1)`$ in the lower half-plane and, as consequence, the integrals over large semi-circles in the lower half-plane vanish. Using the Cauchy formula to calculate the integral (5.8) in the asymptotic region $`x_{k+1}x_k,(k=1,\mathrm{},N1)`$, it is easy to see that the asymptotics of the function $`\psi _\gamma `$ are determined precisely in terms of the corresponding Harish-Chandra function (3.3).
The scalar product (3.4) is the consequence of the Plancherel formula proved in for the $`SL(N,)`$ case. The formula (3.5) can be proved by induction.
###### Remark 5.1
In (eqs.(4.7),(4.18)) the eigenfunctions for an open Toda chain have been constructed in terms of Whittaker function which has a standard integral representation corresponding to the Iwasawa decomposition of semi-simple group (see for example ). Our expression (3.1), being obtained in the framework of Quantum Inverse Scattering Method, seems quite different. Nevertheless, both representations do define the same function (this can be shown by comparing the corresponding asymptotics and analytical properties). Hence, one can consider the representation (3.1) as a new one for Whittaker function -.
## 6 Periodic chain: $`\gamma `$-representation, eigenfunctions, <br>and Plancherel formula
Let $`\mathrm{\Psi }_E(x,x_N)`$ be the fast decreasing solution of the problem (4.11). We define the function $`C(\gamma ;E)`$ by the generalized Fourier transform:
$$\begin{array}{c}\delta (E_1ϵ)C(\gamma ;E)=\underset{\text{}^{N1}}{}\mathrm{\Psi }_E(x,x_N)\overline{\mathrm{\Psi }}_{\gamma ,ϵ}(x_0,x)𝑑x𝑑x_N\end{array}$$
(6.1)
###### Lemma 6.1
The function $`C(\gamma )`$ possesses the following properties:
* It is a symmetric function with respect to $`\gamma `$-variables.
* It is an entire function of $`\gamma ^{N1}`$.
* The function $`C(\gamma )`$ obeys the asymptotics
$$\begin{array}{c}C(\gamma ;E)|\gamma _k|^{N/2}\mathrm{exp}\left\{\frac{\pi N|\gamma _k|}{2\mathrm{}}\right\}\end{array}$$
(6.2)
as $`\mathrm{Re}\gamma _k\pm \mathrm{}`$ in the strip $`|\mathrm{Im}\gamma _k|\mathrm{}`$.
* The function $`C(\gamma )`$ satisfies the multi-dimensional Baxter equation
$$\begin{array}{c}t(\gamma _j;E)C(\gamma ;E)=i^NC(\gamma +i\mathrm{}e_j;E)+i^NC(\gamma i\mathrm{}e_j;E)\end{array}$$
(6.3)
where $`t(\gamma ;E)`$ is defined by (4.12).
Proof. The symmetry of the function $`C(\gamma )`$ is obvious.
We present here only a sketch of the proof of the statements (ii) and (iii).
The statement (ii) follows from the assertion that the auxiliary function $`\mathrm{\Psi }_{\gamma ,ϵ}`$ is an entire one while the solution of the periodic chain vanishes very rapidly as $`|x_kx_{k+1}|\mathrm{}`$. <sup>5</sup><sup>5</sup>5 Actually, the boundary conditions have the same importance here as the requirement of compact support in the theory of the analytic continuation for the usual Fourier transform. .
(iii) The asymptotics (6.2) is a combination of two factors. The first one comes from the asymptotics (2.8) while the additional factor $`|\gamma _k|^1\mathrm{exp}\{\pi |\gamma _k|/\mathrm{}\}`$ results from the stationary phase method while calculating the multiple integral including the function (3.1). The calculation is based heavily upon the exact asymptotics of the function $`\mathrm{\Psi }_E(x,x_N)`$ as $`|x_kx_{k+1}|\mathrm{}`$.
The proof of (iv) is simple. Using the definition (4.11) and integrating by parts (evidently, boundary terms vanish), one obtains
$$\begin{array}{c}\delta (E_1ϵ)t(\gamma _j;E)C(\gamma )\underset{\text{}^{N1}}{}\left\{\widehat{t}(\gamma _j)\mathrm{\Psi }_E(x,x_N)\right\}\overline{\mathrm{\Psi }}_{\gamma ,ϵ}(x,x_N)𝑑x𝑑x_N=\\ =\underset{\text{}^{N1}}{}\mathrm{\Psi }_E(x,x_N)\overline{\widehat{t}(\gamma _j)\mathrm{\Psi }_{\gamma ,ϵ}(x,x_N)}𝑑x𝑑x_N\end{array}$$
(6.4)
Taking into account the relation (5.4), the Baxter equation (6.3) follows from definition (6.1).
Now we prove Theorem 3.2. Using the completeness condition
$$\begin{array}{c}\underset{\text{}^N}{}\mu (\gamma )\mathrm{\Psi }_{\gamma ,ϵ}(x,x_N)\overline{\mathrm{\Psi }}_{\gamma ,ϵ}(y,y_N)𝑑\gamma 𝑑ϵ=\mathrm{\hspace{0.17em}2}\pi \mathrm{}\delta (xy)\delta (x_Ny_N)\end{array}$$
(6.5)
which is a corollary of (3.5), the inversion of the formula (6.1) results to expression
$$\begin{array}{c}\mathrm{\Psi }_E(x,x_N)=\frac{1}{2\pi }\underset{\text{}^{N1}}{}\mu (\gamma )C(\gamma ;E)\mathrm{\Psi }_{\gamma ,E_1}(x,x_N)𝑑\gamma \end{array}$$
(6.6)
The integral (6.6) is correctly defined. Indeed, the measure (3.6) is an entire function. Therefore, there are no poles in the integrand. Moreover,
$$\begin{array}{c}\mu (\gamma )|\gamma _k|^{N2}\mathrm{exp}\left\{\frac{\pi }{\mathrm{}}(N2)|\gamma _k|\right\}\end{array}$$
(6.7)
as $`|\gamma _k|\mathrm{}`$. Taking into account the asymptotics (2.8) and (6.2) one concludes that the integrand has the behavior $`|\gamma _k|^1\mathrm{exp}\{\pi |\gamma _k|/\mathrm{}\}`$ as $`|\gamma _k|\mathrm{}`$. Therefore, the integral (6.5) is convergent.
One can directly prove the spectral problem (4.11) calculating the action of the operator $`\widehat{t}_N(\lambda )=A_N(\lambda )+D_N(\lambda )`$ on the right hand side of (6.6) with the help of the formulae (5.3b), (5.3c). The calculation is performed similarly to those of Lemma 5.1, using the analytical properties of the integrand and the Baxter equation (6.3) (see for details).
The last step is to prove that the function (6.6) satisfies to integrability requirement (2.4). Using the scalar product
$$\begin{array}{c}\underset{\text{}^N}{}\overline{\mathrm{\Psi }}_{\gamma ^{},ϵ^{}}(x,x_N)\mathrm{\Psi }_{\gamma ,ϵ}(x,x_N)𝑑x𝑑x_N=(2\pi \mathrm{})\frac{\mu (\gamma )}{(N1)!}\delta (ϵϵ^{})_{sW}\delta (s\gamma \gamma ^{})\end{array}$$
(6.8)
one can write the Plancherel formula
$$\begin{array}{c}2\pi \mathrm{}\underset{\text{}^N}{}\overline{\mathrm{\Psi }}_E^{}(x,x_N)\mathrm{\Psi }_E(x,x_N)𝑑x𝑑x_N=\delta (E_1E_1^{})\underset{\text{}^{N1}}{}\mu (\gamma )\overline{C}(\gamma ;E^{})C(\gamma ;E)𝑑\gamma \end{array}$$
(6.9)
The integral in the r.h.s. of (6.9) is absolutely convergent due to asymptotics (6.2) and (6.7). Hence, the norm $`\mathrm{\Psi }_E`$ is finite modulo $`GL(1)`$ $`\delta `$-function $`\delta (E_1E_1^{})`$ (see the corresponding factor in (2.3) which leads to this function) and requirement (2.4) is fulfilled. Hence, Theorem 3.2 is proved.
## 7 Solution of Baxter equation
It is well known (see also for details) that the solution to the Baxter equation (3.10) with the asymptotics (3.11) can be written in the following separated form:
$$\begin{array}{c}C(\gamma ;E)=_{j=1}^{N1}\frac{c_+(\gamma _j;E)\xi c_{}(\gamma _j;E)}{\underset{k=1}{\overset{N}{}}\mathrm{sinh}\frac{\pi }{\mathrm{}}(\gamma _j\delta _k)}\end{array}$$
(7.1)
where $`\xi `$ and $`\delta _k`$ are an arbitrary constants and the entire functions $`c_\pm (\gamma )`$ are defined it terms of $`\times `$ determinants:
$$\begin{array}{c}c_+(\gamma )=\frac{1}{\underset{k=1}{\overset{N}{}}\mathrm{}^{{\scriptscriptstyle \frac{i\gamma }{\mathrm{}}}}\mathrm{\Gamma }\left(1{\scriptscriptstyle \frac{i}{\mathrm{}}}\left(\gamma \lambda _k\right)\right)}|\begin{array}{cccccc}1& \frac{1}{t\left(\gamma +i\mathrm{}\right)}& 0& \mathrm{}& \mathrm{}& \mathrm{}\\ \frac{1}{t\left(\gamma +2i\mathrm{}\right)}& 1& \frac{1}{t\left(\gamma +2i\mathrm{}\right)}& 0& \mathrm{}& \mathrm{}\\ 0& \frac{1}{t\left(\gamma +3i\mathrm{}\right)}& 1& \frac{1}{t\left(\gamma +3i\mathrm{}\right)}& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}|\end{array}$$
(7.2a)
$$\begin{array}{c}c_{}(\gamma )=\frac{1}{\underset{k=1}{\overset{N}{}}\mathrm{}^{{\scriptscriptstyle \frac{i\gamma }{\mathrm{}}}}\mathrm{\Gamma }\left(1+{\scriptscriptstyle \frac{i}{\mathrm{}}}\left(\gamma \lambda _k\right)\right)}|\begin{array}{cccccc}\mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& 0& \frac{1}{t\left(\gamma 3i\mathrm{}\right)}& 1& \frac{1}{t\left(\gamma 3i\mathrm{}\right)}& 0\\ \mathrm{}& \mathrm{}& 0& \frac{1}{t\left(\gamma 2i\mathrm{}\right)}& 1& \frac{1}{t\left(\gamma 2i\mathrm{}\right)}\\ \mathrm{}& \mathrm{}& \mathrm{}& 0& \frac{1}{t\left(\gamma i\mathrm{}\right)}& 1\end{array}|\end{array}$$
(7.2b)
and $`\lambda _k\lambda _k(E)`$ are the roots of the polynomial $`t(\gamma )t_N(\gamma ;E)`$.
On the other hand, the solution (7.1) is not an entire function in general since the denominator in (7.1) has an infinite number of poles at $`\gamma =\delta _k+i\mathrm{}n_k,n_k,k=1,\mathrm{},N`$. The poles are cancelled only if the following conditions hold:
$$\begin{array}{c}c_+(\delta _k+i\mathrm{}n_k)=\xi c_{}(\delta _k+i\mathrm{}n_k)\end{array}$$
(7.3)
In turn, this means that the Wronskian
$$\begin{array}{c}W(\gamma )=c_+(\gamma )c_{}(\gamma +i\mathrm{})c_+(\gamma +i\mathrm{})c_{}(\gamma )\end{array}$$
(7.4)
vanishes at $`\gamma =\delta _k+i\mathrm{}n_k`$. The Wronskian is $`i\mathrm{}`$-periodic function and possesses exactly $`N`$ real roots $`\delta _k(E)`$ . Therefore, the solution (7.1) has no poles if one takes $`\delta _k=\delta _k(E)`$ provided that the constant $`\xi `$ is chosen in such a way that
$$\begin{array}{c}\xi =\frac{c_+(\gamma )}{c_{}(\gamma )}|_{\gamma =\delta _k(E)}k=1,\mathrm{},N\end{array}$$
(7.5)
Hence, one arrives at the following
###### Lemma 7.1
The function
$$\begin{array}{c}C(\gamma ;E)=_{j=1}^{N1}\frac{c_+(\gamma _j;E)\xi (E)c_{}(\gamma _j;E)}{\underset{k=1}{\overset{N}{}}\mathrm{sinh}\frac{\pi }{\mathrm{}}\left(\gamma _j\delta _k(E)\right)}\end{array}$$
(7.6)
where $`\delta _k(E)`$ are real zeros of the Wronskian (7.4) and the constant $`\xi `$ is chosen according to (7.5), satisfies to conditions of Lemma 6.1.
The quantization conditions
$$\begin{array}{c}\frac{c_+(\delta _1)}{c_{}(\delta _1)}=\mathrm{}=\frac{c_+(\delta _N)}{c_{}(\delta _N)}\end{array}$$
(7.7)
determine the energy spectrum of the problem. They have been obtained for the first time by Gutzwiller using quite different method.
To prove Theorem 3.3, one should substitute the solution (7.6) into the integral formula (3.8) and calculate the residues coming from individual terms
$$\begin{array}{c}\frac{c_\pm (\gamma _j;E)}{\underset{k=1}{\overset{N}{}}\mathrm{sinh}\frac{\pi }{\mathrm{}}\left(\gamma _j\delta _k(E)\right)}\end{array}$$
(7.8)
The result is exactly the sum over all possible poles of expressions (7.8) and essentially coincides with (3.14) (see careful analysis in ).
## Acknowledgments
One of us (D.L.) is deeply indebted to the organizers of the MMRD 2000 Workshop in Leeds for stimulating and warm atmosphere. Our particular thanks go to M.Semenov-Tian-Shansky and E. Sklyanin for numerous and illuminating discussions.
The research was partly supported by grants INTAS 97-1312; RFBR 00-02-16477(S. Kharchev); RFBR 00-02-16530 (D.Lebedev) and by grant 00-15-96557 for Support of Scientific Schools.
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# Integrated random processes exhibiting long tails, finite moments and 1/f spectra
## I Introduction
The nature of the probability distribution of stock market prices has been discussed quantitatively for over a century . An early conjecture that the distribution was Gaussian was found not to be a good fit, largely because of the long tails found in financial data . A later suggestion that the Lévy distribution was a better fit seemed more promising, since this distribution does at least have long tails . On the other hand this distribution has no finite moments, which is a severe limitation. The solution of truncating the distribution in order to obtain finite moments is rather ad hoc and artificial . Moreover, the Lévy distribution has “too fat tails” –as opposed to the “too thin tails” of the Gaussian– to give a good fit to data .
This paper explores in more detail a model previously introduced in order to resolve these difficulties and give a more complete explanation of the appearance of non-Gaussian and self-similar fat tails in the probability distribution , while still keeping the important feature that all moments are finite. The model has the Lévy distribution as a limiting form and has the form of a type of “Edgeworth expansion” for the Lévy distribution. We will give the explicit form of this relationship later in the paper (the Edgeworth series is an expansion procedure which gives corrections to the Gaussian distribution in those cases where the Central Limit Theorem applies ).
While the motivation for the model was the explanation of stock market data, we want to stress in this paper the more general features of the model which we expect will have applications in other areas including physics. An important aspect of our integrated process is that it exhibits a $`1/f`$ spectrum. We say that noise is $`1/f`$ if any correlated random process has a power spectral density which is, in some range, inversely proportional to a power of the frequency, i.e. $`1/f^\nu `$ ($`\nu >0`$). Strictly speaking $`1/f`$ noise corresponds to the case when $`\nu 1`$ otherwise one talks of “flicker noise”. However, here we use the term $`1/f`$ noise to apply to any power law with $`\nu >0`$. From a historical point of view the first experimental observation of $`1/f`$ noise was made in 1925 by J. B. Johnson in studying non-stationary currents in vacuum tubes , and the first attempt of a theoretical explanation was given by Schottky in 1926. Since then $`1/f`$ noise has appeared, not only in semiconductor physics , but in many areas of physics and even natural and social sciences -. In spite of its ubiquity $`1/f`$ noise is not well understood from a theoretical point of view, and although several mechanisms to generate the noise have been proposed -, as is pointed out in they appear to be very specialized and do not address the universality of this noise. In this sense, we hope that the approach herein will shed some light in the understanding of $`1/f`$ noise as it can be considered a physical mechanism generating such kind of noises.
The paper is organized as follows. In Sect. II we present the main results of the paper for clarity, since many details are quite technical. In Sect III we describe in detail our integrated process. In Sect. IV we impose the self-scaling property on the probability distribution and obtain several relevant functions. Section V is devoted to moments and cumulants and the evaluation of the power spectrum which shows the $`1/f`$ character of our process. In Section VI we study the asymptotic behavior of the distribution while in Sect. VII we present, starting form our distribution, an Edgeworth-type series for the Lévy distribution. Conclusions are drawn in Sect. VIII and some more technical details are in Appendices.
## II Main Results
In this section we wish to state the main results of the paper without giving the derivations or being careful to state the range of validity for which they hold. We hope that this will give a good indication of the scope and nature of our results; the interested reader can then fill out this basic framework by proceeding to later sections.
As explained in the next section, our integrated process $`X(t)`$ is a continuos superposition of colored shot noises parametrized by $`u`$. It takes the form
$$X(t)=_{\mathrm{}}^{\mathrm{}}\left[\underset{k=1}{\overset{\mathrm{}}{}}A_k(u)\varphi (tT_k(u);u)\right]𝑑u.$$
There are several components: (i) The pulse shape function, $`\varphi (t,u)`$, (ii) the jump amplitude $`A_k(u)`$ of the $`k`$th pulse, (iii) the jump time $`T_k(u)`$ of the $`k`$th pulse. The jump amplitudes, $`A_k(u)`$, are identically distributed random variables distributed according the probability density function (pdf) $`h(x,u)`$, i.e. $`h(x,u)dx=\text{Prob}\{x<A_k(u)<x+dx\}`$. The jump times are assumed to follow a Poisson process with parameter $`\lambda (u)`$. Therefore the model is characterized by the functions $`\varphi (t,u)`$, $`h(x,u)`$ and $`\lambda (u)`$. There are very few restrictions on these functions (one is the causality condition $`\varphi (t,u)=0`$ for $`t<0`$) but we will argue that it is natural to assume the scaling forms
$$\varphi (t,u)=\varphi [\lambda (u)t]\mathrm{and}h(x,u)=\frac{1}{\sigma (u)}h\left[\frac{x}{\sigma (u)}\right],$$
where $`\sigma (u)`$ is the standard deviation of the jump amplitudes. The model is now specified by the four functions of a single variable $`\varphi `$, $`h`$, $`\sigma `$ and $`\lambda `$.
One of the main goals of this paper is to find the probability distribution of the process $`X(t)`$, this distribution will be obtained through the characteristic function (cf) $`\stackrel{~}{p}(\omega ,t)=\mathrm{exp}[i\omega X(t)],`$ which is the Fourier transform of the density $`p(x,t)`$. By assuming the scaling forms given above the one time distribution of $`X(t)`$ is explicitly given by the following cf
$$\stackrel{~}{p}(\omega ,t)=\mathrm{exp}\left\{bt\omega ^\alpha _0^{\mathrm{}}\frac{dz}{z^{1+\alpha }}\left[1_0^1\stackrel{~}{h}\left[z\varphi \left(\frac{bt\omega ^\alpha s}{z^\alpha }\right)\right]𝑑s\right]\right\},$$
where $`b>0`$. When $`\varphi (x)=\theta (x)`$ is the Heaviside step function then the input shot noises are white and $`X(t)`$ is the Lévy process:
$$\stackrel{~}{p}(\omega ,t)=e^{Mt\omega ^\alpha },$$
where $`M=b_0^{\mathrm{}}𝑑z[1\stackrel{~}{h}(z)]/z^{1+\alpha }`$. This provides us with an alternative interpretation of the Lévy process since, in our case, Lévy processes are continuos superpositions of families of white Poissonian shot noises.
Contrary to Lévy processes, our integrated process can have finite moments of any order, thus cumulants are given by
$$X^n(t)=\frac{i^n\stackrel{~}{h}^{(n)}(0)}{n}(bt)^{n/\alpha }_0^{\mathrm{}}\frac{\varphi ^n(x)}{x^{n/\alpha }}𝑑x,$$
$`(n=1,2,3,\mathrm{})`$. Note that the second cumulant, (i.e., the variance), is proportional to $`(bt)^{2/\alpha }`$ and $`X(t)`$ presents anomalous diffusion behavior (see Sect. V for a detailed discussion on limiting values and bounds for exponent $`\alpha `$ and other parameters related to the asymptotic behavior of the pulse shape function $`\varphi (x)`$).
We define the stationary correlation function by $`C(\tau )=lim_t\mathrm{}X(t+\tau )X(t)`$. For our process this reads
$$C(\tau )=b_0^{\mathrm{}}\frac{dz}{z^{1+\alpha }}_0^{\mathrm{}}\varphi \left(bz^\alpha t^{}\right)\varphi \left(bz^\alpha (t^{}+\tau )\right)𝑑t^{}.$$
The power spectral density of the process $`X(t)`$, given by the Fourier transform of the stationary correlation, is
$$\stackrel{~}{C}(\omega )=\frac{K}{\omega ^{1+2/\alpha }},$$
and $`X(t)`$ is $`1/f`$ noise with exponent $`\nu =1+2/\alpha `$.
We also can perform the asymptotic analysis of the probability distribution of $`X(t)`$ without having to specify any particular form for $`h(x)`$ and $`\varphi (x)`$ thus keeping the maximum level of generality. Specifically we show in Sect. VI that the center of the distribution is approached by a Lévy distribution
$$\stackrel{~}{p}(\omega ,t)e^{L(t)\omega ^\delta }\omega \mathrm{},$$
where $`0<\delta <2`$. We refer the reader to see Sect. VI for more details and for the behavior of the tails of the distribution which are mainly determined by the behavior of the jump pdf $`h(x)`$.
The relation to the Lévy distribution is explored in more detail in Sect. VII where we present and alternative (and exact) expression for the cf which decompose the distribution of the integrated process $`X(t)`$ into that of Lévy plus an additional term:
$`\mathrm{ln}\stackrel{~}{p}(\omega ,t)=\mathrm{ln}\stackrel{~}{p}_{\mathrm{Levy}}(\omega ,t){\displaystyle \frac{1}{\alpha }}{\displaystyle _0^{\mathrm{}}}[\stackrel{~}{h}\left([bt/x]^{1/\alpha }\omega \right)`$ $``$ $`\stackrel{~}{h}\left([bt/x]^{1/\alpha }\omega \varphi (x)\right)]dx`$
$`+`$ $`{\displaystyle _0^{\mathrm{}}}x{\displaystyle \frac{\varphi ^{}(x)}{\varphi (x)}}\left[1\stackrel{~}{h}\left([bt/x]^{1/\alpha }\omega \varphi (x)\right)\right]𝑑x.`$
Note that when $`\varphi (x)`$ is the rectangular step function this equation reduces to the Lévy distribution. Therefore, when $`\varphi (x)`$ is a step-like function close to the Heaviside function, this alternative expression can be used as the starting point of an Edgeworth-type expansion procedure giving corrections to the Lévy distribution.
## III The integrated process
Let $`X(t)`$ be a random process formed by a continuos superposition of independent shot-noise processes:
$$X(t)=_{\mathrm{}}^{\mathrm{}}Y(u,t)𝑑u,$$
(1)
where for any fixed time $`t`$, $`Y(u,t)`$ are independent random variables for different values of parameter $`u`$ (see Eq. (11) below) and for any fixed value of $`u`$, $`Y(u,t)`$ is a colored shot-noise process represented by a countable superposition of pulses of identical shape:
$$Y(u,t)=\underset{k=1}{\overset{\mathrm{}}{}}A_k(u)\varphi [tT_k(u);u],$$
(2)
where $`T_k(u)`$ marks the onset of the $`k`$th pulse, and $`A_k(u)`$ is its amplitude. Both $`T_k(u)`$ and $`A_k(u)`$ are independent and identically distributed random variables with probability density functions given by $`h(a,u)`$ and $`\psi (t,u)`$, respectively. The pulse shape $`\varphi (t,u)`$ has to fulfill the “causality condition”, i.e., $`\varphi (t,u)=0`$ for $`t<0`$ .
We assume that the occurrence of jumps is a Poisson process, in this case the shot-noise $`Y(t,u)`$ is Markovian, and the pdf for the time interval between jumps, $`\psi (t,u)dt=\text{P}\{t<T_k(u)T_{k1}(u)<t+dt\}`$, is exponential:
$$\psi (t,u)=\lambda (u)e^{\lambda (u)t}(t0),$$
(3)
where $`\lambda (u)`$ is the mean jump frequency, i.e., $`1/\lambda (u)`$ is the mean time between two consecutive jumps . We recall that jump amplitudes $`A_k(u)`$ are identically distributed (for all $`k=1,2,3,\mathrm{}`$) and independent random variables (for all $`k`$ and $`u`$). In what follows we will assume that they have zero mean and a pdf, $`h(x,u)dx=\text{P}\{x<A_k(u)<x+dx\}`$, depending on a single “dimensional” parameter which, without loss of generality, we assume to be the standard deviation of jumps $`\sigma (u)=\sqrt{A_k^2(u)}`$. That is,
$$h(x,u)=\frac{1}{\sigma (u)}h\left[\frac{x}{\sigma (u)}\right].$$
(4)
Before proceeding further with the probability distribution of the integrated process $`X(t)`$ given by Eq. (1), we note that following Rice’s method one can easily obtain all the probability distributions of the shot noise $`Y(t,u)`$ via their cf’s
$$\stackrel{~}{p}_Y(\omega _1,t_1;\mathrm{};\omega _n,t_n;u)=\mathrm{exp}\left[i\underset{k=1}{\overset{n}{}}\omega _kY(t_k,u)\right].$$
(5)
In Appendix A we show that
$$\mathrm{ln}\stackrel{~}{p}_Y(\omega ,t;u)=\lambda (u)\left[t_0^t\stackrel{~}{h}[\omega \sigma (u)\varphi (t^{},u)]𝑑t^{}\right],$$
(6)
and (supposing that $`t_2>t_1`$)
$`\mathrm{ln}\stackrel{~}{p}_Y(\omega _1,t_1;\omega _2,t_2;u)=\lambda (u)[t_2{\displaystyle _0^{t_1}}\stackrel{~}{h}[\omega _1\sigma (u)\varphi (t^{},u)`$ $`+`$ $`\omega _2\sigma (u)\varphi (t^{}+t_2t_1,u)]dt^{}`$ (7)
$``$ $`{\displaystyle _0^{t_2t_1}}\stackrel{~}{h}[\omega _2\sigma (u)\varphi (t^{},u)]dt^{}],`$ (8)
where $`\stackrel{~}{h}(\omega )`$ is the Fourier transform of the jump pdf $`h(x)`$.
Let us now evaluate the probability distribution of the integrated process $`X(t)`$. In terms of the cumulants, $`Y(t,u)`$ of the shot noise $`Y(t,u)`$ we see that the one time characteristic function of $`X(t)`$ can be written as
$$\stackrel{~}{p}_X(\omega ,t)=\mathrm{exp}\left\{\underset{k=1}{\overset{\mathrm{}}{}}\frac{(i\omega )^k}{k!}\left[_{\mathrm{}}^{\mathrm{}}Y(u,t)𝑑u\right]^k\right\}.$$
(9)
That is,
$$\mathrm{ln}\stackrel{~}{p}_X(\omega ,t)=\underset{k=1}{\overset{\mathrm{}}{}}\frac{(i\omega )^k}{k!}_{\mathrm{}}^{\mathrm{}}\mathrm{}_{\mathrm{}}^{\mathrm{}}Y(u_1,t)\mathrm{}Y(u_k,t)𝑑u_1\mathrm{}𝑑u_k$$
(10)
But, by our assumptions on the process $`Y(u,t)`$ we have
$$Y(u_1,t)\mathrm{}Y(u_k,t)=Y^k(u_1,t)\delta (u_1u_2)\mathrm{}\delta (u_{k1}u_k).$$
(11)
Therefore,
$$\mathrm{ln}\stackrel{~}{p}_X(\omega ,t)=\underset{k=1}{\overset{\mathrm{}}{}}\frac{(i\omega )^k}{k!}_{\mathrm{}}^{\mathrm{}}Y^k(u,t)𝑑u,$$
that is,
$$\mathrm{ln}\stackrel{~}{p}_X(\omega ,t)=_{\mathrm{}}^{\mathrm{}}\mathrm{ln}\stackrel{~}{p}_Y(\omega ,t;u)𝑑u.$$
(12)
Note that this line of reasoning can be easily extend to the $`n`$th time distribution, with the result:
$$\mathrm{ln}\stackrel{~}{p}_X(\omega _1,t_1;\mathrm{};\omega _n,t_n)=_{\mathrm{}}^{\mathrm{}}\mathrm{ln}\stackrel{~}{p}_Y(\omega _1,t_1;\mathrm{};\omega _n,t_n;u)𝑑u.$$
(13)
Going back to our integrated process we have from Eqs. (6)-(8) and (12)-(13) that the one time characteristic function of $`X(t)`$ reads
$$\mathrm{ln}\stackrel{~}{p}(\omega ,t)=_{\mathrm{}}^{\mathrm{}}𝑑u\lambda (u)\left[t_0^t\stackrel{~}{h}[\omega \sigma (u)\varphi (t^{},u)]𝑑t^{}\right],$$
(14)
while the two time cf is ($`t_2>t_1`$)
$`\mathrm{ln}\stackrel{~}{p}(\omega _1,t_1;\omega _2,t_2)={\displaystyle _{\mathrm{}}^{\mathrm{}}}du\lambda (u)\{{\displaystyle _0^{t_1}}dt^{}[\stackrel{~}{h}[\omega _1\sigma (u)\varphi (t^{},u)`$ $`+`$ $`\omega _2\sigma (u)\varphi (t^{}+t_2t_1,u)]1]`$ (15)
$`+`$ $`{\displaystyle _0^{t_2t_1}}dt^{}[\stackrel{~}{h}[\omega _2\sigma (u)\varphi (t^{},u)]1]\},`$ (16)
where we have dropped the subscript $`X`$. Obviously these are formal expressions, as long as we do not provide the functional dependence of $`\lambda (u)`$ and $`\sigma (u)`$ on the parameter $`u`$. We will do so in the next section using scaling arguments.
## IV Scaling
In order to proceed further we need to specify the functional form of $`\lambda (u)`$ and $`\sigma (u)`$. Of course that form will depend on the specific features of the problem at hand. At this point we choose what seems to us one of the most general ways of proceeding, we thus suppose that our integrated process $`X(t)`$ possesses self-scaling properties. Following this path we must first assume that the pulse function is of the form
$$\varphi (u,t)=\varphi [\lambda (u)t],$$
(17)
which turns $`\varphi (u,t)`$ into a function of the single dimensionless variable $`\lambda (u)t`$. Substituting this into Eq. (14), defining new integration variables $`s=t^{}/t`$ and $`z=\omega \sigma (u)`$, and supposing that $`\sigma (\mathrm{})=0`$ and $`\sigma (\mathrm{})=\mathrm{}`$, we obtain
$$\mathrm{ln}\stackrel{~}{p}(\omega ,t)=_0^{\mathrm{}}𝑑z\frac{\lambda t}{\omega \sigma ^{}}\left\{1_0^1\stackrel{~}{h}[z\varphi (\lambda st)]𝑑s\right\},$$
(18)
where the prime on $`\sigma `$ denotes derivative. We now impose the self-scaling property on the cf, that is, we assume that $`\stackrel{~}{p}(\omega ,t)`$ is a function of the single variable $`\omega t^{1/\alpha }`$:
$$\stackrel{~}{p}(\omega ,t)=f(\omega t^{1/\alpha }).$$
(19)
On the other hand we note that in Eq. (18), the quantities $`\sigma ^{}`$ and $`\lambda `$ are functions of $`z`$ and $`\omega `$. Then scaling (19) implies
$$\lambda =B(z)\omega ^\alpha ,\frac{\lambda }{\omega \sigma ^{}}=A(z)\omega ^\alpha ,$$
(20)
where $`A(z)`$ and $`B(z)`$ are arbitrary functions to be determined. From these two relations we get $`\sigma ^{}=C(z)/\omega `$, where $`C(z)=B(z)/A(z)`$. In the original variable $`u`$ we have
$$\sigma ^{}=\frac{C(\omega \sigma (u))}{\omega }.$$
(21)
But $`\sigma ^{}=\sigma ^{}(u)`$ is independent of $`\omega `$. Therefore, the unknown function $`C`$ has to be of the form $`C(\omega \sigma )=k\omega \sigma `$ where $`k`$ is a constant. Hence $`\sigma ^{}(u)=k\sigma (u)`$, whence $`\sigma (u)=\sigma _0e^{ku}`$. Finally absorbing the constant $`k`$ inside the variable $`u`$ we obtain the functional dependence of the jump variance $`\sigma `$ on parameter $`u`$:
$$\sigma (u)=\sigma _0e^u,$$
(22)
where $`\sigma _0`$ is a constant. Moreover we see from Eq. (20) that $`\lambda =B(\omega \sigma (u))\omega ^\alpha `$. But, again $`\lambda =\lambda (u)`$ is independent of $`\omega `$. In consequence $`B(\omega \sigma )=b(\omega \sigma )^\alpha `$ where $`b`$ is an arbitrary constant. Substituting this into the first relation of Eq. (20) yields the “dispersion relation” between the mean frequency $`\lambda `$ and the jump variance $`\sigma `$:
$$\lambda =b/\sigma ^\alpha .$$
(23)
The functional dependence of $`\lambda `$ on $`u`$ is obtained by combining Eqs. (22)-(23),
$$\lambda (u)=\lambda _0e^{\alpha u},$$
(24)
where $`\lambda _0=b/\sigma _0^\alpha `$, or equivalently
$$b=\lambda _0\sigma _0^\alpha .$$
(25)
Collecting results we see from Eq. (18) and Eqs. (22)-(25) that the one-time characteristic function reads
$$\stackrel{~}{p}(\omega ,t)=\mathrm{exp}\left\{bt\omega ^\alpha _0^{\mathrm{}}\frac{dz}{z^{1+\alpha }}\left[1_0^1\stackrel{~}{h}\left[z\varphi \left(bt\omega ^\alpha s/z^\alpha \right)\right]𝑑s\right]\right\}.$$
(26)
It is sometimes convenient to rewrite this equation and use the alternative form of $`\stackrel{~}{p}(\omega ,t)`$ given by
$$\stackrel{~}{p}(\omega ,t)=\mathrm{exp}\left\{bt_0^{\mathrm{}}\frac{dz}{z^{1+\alpha }}\left[1_0^1\stackrel{~}{h}\left[z\omega \varphi \left(bts/z^\alpha \right)\right]𝑑s\right]\right\},$$
(27)
or equivalently
$$\stackrel{~}{p}(\omega ,t)=\mathrm{exp}\left\{b_0^{\mathrm{}}\frac{dz}{z^{1+\alpha }}_0^t𝑑t^{}\left[\stackrel{~}{h}\left[z\omega \varphi \left(bt^{}/z^\alpha \right)\right]1\right]\right\}.$$
(28)
Starting from Eq. (16) and following an analogous reasoning based on the scaling assumption, we obtain the following expression for the two time characteristic function of the integrated process (with $`t_2>t_1`$)
$`\mathrm{ln}\stackrel{~}{p}(\omega _1,t_1;\omega _2,t_2)=b{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dz}{z^{1+\alpha }}}\{{\displaystyle _0^{t_1}}dt^{}[\stackrel{~}{h}[z\omega _1\varphi (bz^\alpha t^{})`$ $`+`$ $`z\omega _2\varphi (bz^\alpha (t^{}+t_2t_1))]1]`$ (29)
$`+`$ $`{\displaystyle _0^{t_2t_1}}dt^{}[\stackrel{~}{h}[z\omega _2\varphi (bz^\alpha t^{})]1]\}.`$ (30)
Equations (26)-(30) are some of the key results of the paper. Since, as we will see next, they constitute a generalization of the Lévy distribution with finite moments.
## V Moments, Cumulants and Power Spectrum
We first note from Eq. (26) that if the pulse shape function is the Heaviside step function:
$$\varphi (t)=\{\begin{array}{cc}1,\hfill & \text{if }t>0\hfill \\ 0,\hfill & \text{otherwise,}\hfill \end{array}$$
(31)
then the integrated process $`X(t)`$ is identically a Lévy process, regardless the jump pdf $`h(x)`$:
$$\stackrel{~}{p}(\omega ,t)=\mathrm{exp}\left[Mt\omega ^\alpha \right],$$
(32)
where
$$M=b_0^{\mathrm{}}\frac{dz}{z^{1+\alpha }}\left[1\stackrel{~}{h}(z)\right].$$
(33)
Therefore, following our model, Lévy processes can be viewed as a continuos superposition of families of rectangular pulses occurring at random Poisson times. The usual range of the exponent $`\alpha `$ in Lévy flights is $`0<\alpha <2`$. In such a case $`X(t)`$ has no finite moment but the first one . In actual situations, one is unlikely to meet with perfect rectangular pulses (showing sudden changes) in such a case all moments can be finite and are easily evaluated from Eq. (27) through
$$X^n(t)=i^n\frac{^n\stackrel{~}{p}(\omega ,t)}{\omega ^n}|_{\omega =0}.$$
Thus for instance the second moment is given by (note that due to Eq. (4) $`\stackrel{~}{h}^{\prime \prime }(0)=1`$)
$$X^2(t)=bt_0^{\mathrm{}}\frac{dz}{z^{\alpha 1}}_0^1\varphi ^2(bts/z^\alpha )𝑑s.$$
(34)
As an illustrative example suppose that our pulse function has the form
$$\varphi (t)=\{\begin{array}{cc}1e^{kt},\hfill & \text{if }t>0\hfill \\ 0,\hfill & \text{otherwise,}\hfill \end{array}$$
(35)
where $`k>0`$ is a constant (note that if $`k`$ is large then $`\varphi (t)`$ approaches to the rectangular pulse (31)). In the Appendix B we show that (see also Eq. (46) below)
$$X^2(t)=Dt^{2/\alpha },(2>\alpha >2/3),$$
(36)
where $`D=\alpha b^2k^{1+2/\alpha }(12^{2/\alpha 2})\mathrm{\Gamma }(22/\alpha )/(2\alpha )`$. Since $`1<(2/\alpha )<3`$, Eq. (36) clearly shows a superdiffusive behavior.
In fact we can easily obtain a closed expression not for the $`n`$th moment but for the $`n`$th cumulant
$$X^n(t)i^n\frac{^n\mathrm{ln}\stackrel{~}{p}(\omega ,t)}{\omega ^n}|_{\omega =0}.$$
¿From Eq. (27) we have
$$X^n(t)=i^nbt\stackrel{~}{h}^{(n)}(0)_0^{\mathrm{}}z^{n1\alpha }𝑑z_0^1\varphi ^n(bts/z^\alpha )𝑑s.$$
If in the double integral on the right hand side of this equation we define a new integration variable $`x`$ by $`s=(z^\alpha /bt)x`$ and exchange the order of integration we get
$$_0^{\mathrm{}}z^{n1\alpha }𝑑z_0^1\varphi ^n(bts/z^\alpha )𝑑s=\frac{1}{bt}_0^{\mathrm{}}\varphi ^n(x)𝑑x_0^{(bt/x)^{1/\alpha }}z^{n1}𝑑z,$$
but the last integral is trivially evaluated, and for the $`n`$th cumulant we have
$$X^n(t)=\frac{i^n\stackrel{~}{h}^{(n)}(0)}{n}(bt)^{n/\alpha }_0^{\mathrm{}}\frac{\varphi ^n(x)}{x^{n/\alpha }}𝑑x.$$
(37)
Taking into account that $`\stackrel{~}{h}^{(n)}(0)=0`$ for $`n`$ odd (we have assumed a symmetric jump distribution $`h(x)`$) we write
$$X^{2n1}(t)=0,$$
(38)
and
$$X^{2n}(t)=\frac{(1)^n\stackrel{~}{h}^{(2n)}(0)}{2n}(bt)^{2n/\alpha }_0^{\mathrm{}}\frac{\varphi ^{2n}(x)}{x^{2n/\alpha }}𝑑x,$$
(39)
($`n=1,2,3,\mathrm{}`$). In order to check the convergence of these expressions and the existence of moments, we first assume that
$$\varphi (x)x^\beta ,(x0),$$
(40)
($`\beta >0`$) then the convergence of the integral on the right hand side of Eq. (39) as $`x0`$ implies that the scaling exponent $`\alpha `$ has a lower bound:
$$\alpha >\frac{2n}{1+2n\beta }.$$
(41)
On the other hand if we assume that
$$\varphi (x)x^\gamma ,(x\mathrm{}),$$
(42)
then the convergence of (39) when $`x\mathrm{}`$ implies that the scaling exponent $`\alpha `$ also has an upper bound:
$$\frac{1}{\alpha }>\gamma +\frac{1}{2n}.$$
(43)
Moreover when $`\gamma 0`$ then if
$$\frac{1}{\gamma +1/2}>\alpha >\frac{1}{\beta }$$
(44)
all cumulants will exist. Note that Eq. (43) holds whenever $`\gamma 1/2n`$. Therefore, $`\gamma 0`$ is a sufficient condition for its validity. On the other hand if $`\gamma <1/2`$ there is no upper bound on the accepted values of $`\alpha `$. We also observe that for a step-like function, as that of Eq. (35), where $`\gamma =0`$ then all cumulants will exist if
$$2>\alpha >1/\beta .$$
Finally, for any integrable function $`\varphi (t)`$ over $`[0,\mathrm{})`$ there is no upper bound for $`\alpha `$ and the only condition on $`\alpha `$ for having all moments finite is that $`\alpha >1/\beta `$.
We close this discussion on moments and cumulants with an example. Suppose that the pulse function is given by the step-like function (35). In this case $`\gamma =0`$, $`\beta =1`$ and all moments (and cumulants) will exist if $`1<\alpha <2`$. Cumulants are given by Eqs. (38)-(39). In Appendix B we show that
$$_0^{\mathrm{}}\frac{(1e^{kx})^n}{x^n}𝑑x=A_nk^{1+n/\alpha }\mathrm{\Gamma }(1+nn/\alpha ),$$
(45)
where the numbers $`A_n`$ are given by Eq. (B3) of Appendix B. Finally,
$$X^{2n}(t)=\frac{(1)^n\stackrel{~}{h}^{(2n)}(0)}{2n}A_{2n}k^{1+2n/\alpha }\mathrm{\Gamma }(1+2n2n/\alpha )(bt)^{2n/\alpha }.$$
(46)
We finish this section evaluating the power spectrum of the integrated process $`X(t)`$. Let us first evaluate the correlation function
$$X(t+\tau )X(t)=\frac{^2}{\omega _1\omega _2}\stackrel{~}{p}(\omega _1,t;\omega _2,t+\tau )|_{\omega _1=\omega _2=0}.$$
¿From Eq. (30) we get
$$X(t+\tau )X(t)=b_0^{\mathrm{}}\frac{dz}{z^{1+\alpha }}_0^t\varphi (bz^\alpha t^{})\varphi (bz^\alpha (t^{}+\tau )).$$
(47)
Let $`C(\tau )`$ be the correlation function in the stationary limit $`t\mathrm{}`$, i.e.
$$C(\tau )=\underset{t\mathrm{}}{lim}X(t+\tau )X(t).$$
¿From Eq. (47) we have
$$C(\tau )=b_0^{\mathrm{}}\frac{dz}{z^{1+\alpha }}_0^{\mathrm{}}\varphi (bz^\alpha t^{})\varphi (bz^\alpha (t^{}+\tau )).$$
(48)
Note that the (stationary) variance $`C(0)=\mathrm{}`$ which agrees with the superdiffusive behavior of $`X(t)`$ given by Eq. (36).
The power spectral density of our process is thus given by the Fourier transform of the stationary correlation function
$$\stackrel{~}{C}(\omega )=_{\mathrm{}}^{\mathrm{}}e^{i\omega \tau }C(\tau )𝑑\tau .$$
Substituting Eq. (48) into this equation, performing simple changes of variables and taking into account the causality of the pulse function $`\varphi (t)`$ we finally obtain
$$\stackrel{~}{C}(\omega )=\frac{K}{\omega ^{1+2/\alpha }},$$
(49)
where
$$K=\frac{b^{2/\alpha }}{\alpha }_0^{\mathrm{}}\xi ^{2/\alpha }|\stackrel{~}{\varphi }(\xi )|^2𝑑\xi ,$$
(50)
and $`\stackrel{~}{\varphi }(\xi )`$ is the Fourier transform of $`\varphi (t)`$. The power spectral density $`\stackrel{~}{C}(\omega )`$ exists if the integral
$$J=_0^{\mathrm{}}\xi ^{2/\alpha }|\stackrel{~}{\varphi }(\xi )|^2𝑑\xi <\mathrm{}.$$
In order to prove the existence of $`J`$ we first need that the Fourier transform of the pulse function, $`\stackrel{~}{\varphi }(\xi )`$, exists. Note that any step or step-like function does not have a Fourier transform and consequently the power spectrum is infinite. For the existence of $`\stackrel{~}{\varphi }(\xi )`$ it suffices that $`\varphi (x)`$ be absolutely integrable, and from the asymptotic behavior given by Eq. (42):
$$\varphi (x)x^\gamma (x\mathrm{}),$$
we have to impose that $`\gamma <1`$. This in turn implies that $`\stackrel{~}{\varphi }(\xi )\xi ^{1\gamma }`$ as $`\xi 0`$ and, since $`\gamma <1`$, the integral $`J`$ at its lower limit is always finite for any $`\alpha >0`$. On the other hand if $`\varphi (x)`$ satisfies Eq. (40) as $`x0`$, then $`\stackrel{~}{\varphi }(\xi )\xi ^{1\beta }`$ as $`\xi \mathrm{}`$. Hence, $`J`$ will be finite if $`\alpha >1/(\beta +1/2)`$ (see also Eq. (41)). Therefore, the process $`X(t)`$ has a finite power spectrum if $`\varphi (x)`$ is absolutely integrable on the real line, and $`X(t)`$ has a finite second cumulant, Eq. (41), i.e.
$$\alpha >\frac{1}{\beta +1/2}.$$
(51)
In such a case, we see from Eq. (49) that $`X(t)`$ is $`1/f`$ noise with exponent $`\nu =1+2/\alpha `$. Moreover (recall that $`\alpha >0`$, and Eq. (51))
$$1<\nu <2(1+\beta ).$$
(52)
Finally, when $`\varphi (x)`$ is analytic at $`x=0`$ then $`\beta =1,2,3,\mathrm{}`$ is a positive integer and $`X(t)`$ is flicker noise $`1/f^\nu `$ with $`1<\nu <2(1+n)`$ ($`n=1,2,3,\mathrm{}`$).
## VI Asymptotic Behavior
We will now examine the asymptotic behavior of the one time probability density function (pdf) of the integrated process $`X(t)`$, $`p(x,t)dx=\text{P}\{x<X(t)<x+dx\}`$. For this analysis we distinguish two regions: the “center” ($`x0`$) and the “tails” ($`x\pm \mathrm{}`$) of the distribution. We cannot have a closed expression for the pdf $`p(x,t)`$ until the pulse function $`\varphi (x)`$ and the jump pdf $`h(x)`$ are both specified. Therefore, we will perform the asymptotics on the cf $`\stackrel{~}{p}(\omega ,t)`$. As a well known feature of the harmonic analysis the center of the distribution is determined by the large $`\omega `$ behavior of the cf, while the tails are determined by $`\stackrel{~}{p}(\omega ,t)`$ when $`\omega 0`$ .
We deal first with the center of the distribution where $`\omega \mathrm{}`$. If we assume that the pulse function $`\varphi (x)`$, as $`x\mathrm{}`$, satisfies Eq. (42):
$$\varphi (x)x^\gamma ,(x\mathrm{}),$$
then
$$\stackrel{~}{h}\left[z\varphi \left(bt\omega ^\alpha s/z^\alpha \right)\right]\stackrel{~}{h}\left[\left(bt\omega ^\alpha \right)^\gamma z^{1\alpha \gamma }s^\gamma \right],(\omega \mathrm{}).$$
Substituting this into Eq. (26) and performing the change of variables $`\xi =(bt\omega ^\alpha )^\gamma z^{1\alpha \gamma }`$ we obtain the following Lévy distribution:
$$\stackrel{~}{p}(\omega ,t)\mathrm{exp}\left\{L(bt)^{1/(1\alpha \gamma )}\omega ^{\alpha /(1\alpha \gamma )}\right\},(\omega \mathrm{}),$$
(53)
where
$$L=\frac{1}{1\alpha \gamma }_0^{\mathrm{}}\frac{d\xi }{\xi ^{(1+\alpha \alpha \gamma )/(1\alpha \gamma )}}\left[1_0^1\stackrel{~}{h}(\xi s^\gamma )𝑑s\right].$$
(54)
Note that due to the bounds discussed above (see Eq. (43)) we have $`1\alpha \gamma >\alpha /2`$ hence the Lévy exponent in Eq. (53) satisfies
$$0<\frac{\alpha }{1\alpha \gamma }<2,$$
(55)
and Eq. (53) is well defined. We also note that for a step-like pulse function $`\varphi (t)`$ where $`\gamma =0`$ we obtain the same Lévy distribution, Eq. (32), that satisfies the model for sudden pulses (31). Therefore, for any pulse shape function satisfying condition (42) the center of the pdf is given by a Lévy distribution.
Let us now obtain an asymptotic expression of the pdf $`p(x,t)`$ when $`x\pm \mathrm{}`$, which will be valid if $`\varphi (t)`$ obeys Eqs. (40) and (42), and the exponent $`\alpha `$ is bounded by
$$\beta ^1<\alpha <1/(\gamma +1/2).$$
In this case, we see from Eq. (44) that all cumulants exist. So taking the $`\omega 0`$ limit of Eq (27) we find
$$\stackrel{~}{p}(\omega ,t)1bt_0^{\mathrm{}}\frac{dz}{z^{1+\alpha }}\left[1_0^1\stackrel{~}{h}\left[z\omega \varphi \left(bts/z^\alpha \right)\right]𝑑z\right].$$
(56)
The Fourier inversion of Eq. (56) yields
$$p(x,t)bt_0^{\mathrm{}}\frac{dz}{z^{1+\alpha }}_0^1h(x,z,s)𝑑z(x\pm \mathrm{}),$$
(57)
where we have dropped delta function terms which have no contribution as $`x\mathrm{}`$. Moreover assuming symmetric jump distributions $`h(x)`$ we have
$$h(x,z,s)\frac{1}{\pi }_0^{\mathrm{}}\stackrel{~}{h}\left[z\omega \varphi \left(bts/z^\alpha \right)\right]\mathrm{cos}\omega xd\omega =\frac{1}{z\varphi (bts/z^\alpha )}h\left[\frac{x}{z\varphi (bts/z^\alpha )}\right].$$
Substituting this into Eq. (57) we see that
$$p(x,t)bt_0^{\mathrm{}}\frac{dz}{z^{2+\alpha }}_0^1\frac{ds}{\varphi (bts/z^\alpha )}h\left[\frac{x}{z\varphi (bts/z^\alpha )}\right](x\pm \mathrm{}).$$
(58)
Therefore the tails of the distribution are determined by the jump pdf $`h(x)`$ and the pulse shape function $`\varphi (t)`$.
Finally for the rectangular pulse (31) we have
$$p(x,t)\frac{bt}{|x|^{1+\alpha }}_0^{\mathrm{}}y^\alpha h(y)𝑑y,$$
(59)
which agrees with the expected tail behavior of the Lévy distribution .
## VII Relation to the Lévy distribution
In section V we obtained general expressions (38)-(37) for all of the cumulants, from which it follows that
$`\mathrm{ln}\stackrel{~}{p}(\omega ,t)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{i^n\omega ^n}{n!}}X^n(t)`$ (60)
$`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^n\omega ^{2n}\stackrel{~}{h}^{(2n)}(0)}{2n(2n)!}}(bt)^{2n/\alpha }{\displaystyle _0^1}{\displaystyle \frac{\varphi ^{2n}(x)}{x^{2n/\alpha }}}𝑑x.`$ (61)
In addition, in section VI, we showed that $`p(x,t)`$ was a Lévy distribution at the center of the distribution ($`x0`$) and took the form (58) in the tails ($`x\pm \mathrm{}`$) of the distribution. In this section we will show how the distribution can be separated into a Lévy distribution plus an additional term. This term takes the form of a single integral which can be evaluated once the functions $`\varphi `$ and $`h`$ have been specified.
We begin the analysis by changing variables from $`s`$ to $`x=bts/z^\alpha `$ ($`z`$ fixed) in Eq. (27). This gives
$`\mathrm{ln}\stackrel{~}{p}(\omega ,t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dz}{z}}{\displaystyle _0^{bt/z^\alpha }}𝑑x\left\{1\stackrel{~}{h}\left[z\omega \varphi (x)\right]\right\}`$ (62)
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑x{\displaystyle _0^{(bt/x)^{1/\alpha }}}{\displaystyle \frac{dz}{z}}\left\{1\stackrel{~}{h}\left[z\omega \varphi (x)\right]\right\}`$ (63)
changing the order of integration. At this point we factor out the contribution from the Lévy process by writing (63) as
$$_0^{\mathrm{}}𝑑x_0^{(bt/x)^{1/\alpha }}\frac{dz}{z}\left\{1\stackrel{~}{h}\left[z\omega \right]\right\}_0^{\mathrm{}}𝑑x_0^{(bt/x)^{1/\alpha }}\frac{dz}{z}\left\{\stackrel{~}{h}\left[z\omega \right]\stackrel{~}{h}\left[z\omega \varphi (x)\right]\right\},$$
or, after defining
$$\stackrel{~}{g}(\omega ,t)1\stackrel{~}{h}(\omega ,t),$$
(64)
as
$$_0^{\mathrm{}}𝑑x_0^{(bt/x)^{1/\alpha }}\frac{\stackrel{~}{g}(\omega ,t)}{z}𝑑z+_0^{\mathrm{}}𝑑x_0^{(bt/x)^{1/\alpha }}\frac{dz}{z}\left\{\stackrel{~}{g}\left[z\omega \right]\stackrel{~}{g}\left[z\omega \varphi (x)\right]\right\},$$
The first term is just $`\mathrm{ln}\stackrel{~}{p}(\omega ,t)`$ for the Lévy processes (see Eqs. (32)-(33)). The second term can be simplified by first writing it as
$$_0^{\mathrm{}}𝑑x_{(bt/x)^{1/\alpha }\varphi }^{(bt/x)^{1/\alpha }}\frac{\stackrel{~}{g}(z\omega )}{z}𝑑z$$
and then integrating by parts to give
$$\left|x_{(bt/x)^{1/\alpha }\varphi }^{(bt/x)^{1/\alpha }}\frac{\stackrel{~}{g}(z\omega )}{z}𝑑z\right|_{x=0}^{x=\mathrm{}}_0^{\mathrm{}}𝑑xx\frac{}{x}_{(bt/x)^{1/\alpha }\varphi }^{(bt/x)^{1/\alpha }}\frac{\stackrel{~}{g}(z\omega )}{z}𝑑z.$$
(65)
We assume that $`\stackrel{~}{h}(\omega )`$ is analytic at $`\omega =0`$ and integrable, then
$$\stackrel{~}{g}(\omega )\omega ^2(\omega 0)\mathrm{and}\stackrel{~}{g}(\omega )1(\omega \mathrm{}),$$
and since $`0<\alpha <2`$ the first term in Eq. (65) is zero. Finally
$`\mathrm{ln}\stackrel{~}{p}(\omega ,t)=\mathrm{ln}\stackrel{~}{p}_{\mathrm{Levy}}(\omega ,t)+{\displaystyle \frac{1}{\alpha }}{\displaystyle _0^{\mathrm{}}}[\stackrel{~}{g}\left([bt/x]^{1/\alpha }\omega \right)`$ $``$ $`\stackrel{~}{g}\left([bt/x]^{1/\alpha }\omega \varphi (x)\right)]dx`$ (66)
$`+`$ $`{\displaystyle _0^{\mathrm{}}}x{\displaystyle \frac{\varphi ^{}(x)}{\varphi (x)}}\stackrel{~}{g}\left([bt/x]^{1/\alpha }\omega \varphi (x)\right)𝑑x,`$ (67)
where $`\varphi ^{}(x)`$ is the derivative of the pulse shape function and
$$\mathrm{ln}\stackrel{~}{p}_{\mathrm{Levy}}(\omega ,t)=Mt\omega ^\alpha $$
where $`M`$ is given by Eq. (33). Note that when $`\varphi (x)`$ is the Heaviside step function the integrals on the right hand side of Eq. (67) vanish and Eq. (67) reduces to the Lévy distribution. Therefore, we can look at the second term on the right hand side of Eq. (67) as a correction to the Lévy distribution when $`\varphi (x)`$ is not exactly a Heaviside function but a step-like function very close to the Heaviside function. This may be evaluated, in principle, for any given $`\varphi `$ and $`\stackrel{~}{h}`$. For instance we could take $`\varphi `$ to be of the form (35) with $`k`$ large and the Lorentzian
$$\stackrel{~}{g}(\omega )=\frac{\omega ^2/2}{1+\omega ^2/2},$$
corresponding to $`h(x)=e^{\sqrt{2}|x|}/\sqrt{2}`$. The form of the correction terms depends on the choice of the functions to an extent, and so we will not discuss the explicit form it takes here.
## VIII Conclusions
In this paper we have presented and analyzed a dynamical model based on a process which is a superposition of colored Poisson noises. The model was shown to have several attractive features. The probability density function has long tails which emerged in a natural way and, unlike the Lévy distribution, all the moments of the distribution are finite. We believe that these properties make the distribution an ideal candidate for describing stock market prices .
A property what may have more relevance to physics and other natural sciences is the appearance of $`1/f`$ noise in the power spectrum. Once again we would like to stress that this result flowed naturally from the nature of the model and the scaling assumptions which reduce $`h`$ and $`\varphi `$ from functions of two variables to functions af a single variable.
In a more mathematical context, we believe that the decomposition of the cf of our model into that of the Lévy plus additional terms is interesting, both as an example of an Edgeworth-type expansion and for the nature of the corrections to the Lévy distributions when the parameters of our model are chosen so that our distribution is near to the Lévy one.
There are still some open questions. One of them is the extension of the model to the increments of the process $`Z(\tau ,tt_0)=X(tt_0+\tau )X(tt_0)`$ $`(t>t_0)`$, since in this case we believe that the process $`Z(\tau ,tt_0)`$ becomes stationary when it starts in the infinite past ($`t_0\mathrm{}`$). Another interesting and open question is the actual application of the model to financial time series where some non-white correlation is observed . Both points are presently being investigated.
## ACKNOWLEDGMENTS
This work has been supported in part by Dirección General de Investigación Científica y Técnica under contract No. PB96-0188 and Project No. HB119-0104, and by Generalitat de Catalunya under contract No. 1998 SGR-00015 (JM and MM) AM thanks the British Council for support under “Acciones Integradas” scheme.
## A Characteristic function for colored shot noise
By generalizing Rice’s method , we will now obtain the probability distribution of the shot noise $`Y(u,t)`$ defined by Eq. (2):
$$Y(u,t)=\underset{k=1}{\overset{\mathrm{}}{}}A_k(u)\varphi [tT_k(u);u],$$
(A1)
where we assume that the random variables $`A_k(u)`$ and $`T_k(u)`$ are identically distributed and statistically independent. The jump amplitudes are described by the pdf $`h(x,u)dx=\text{Prob}\{x<A_k(u)<x+dx\}`$ and the jump times $`T_k(u)`$ follow a Poisson distribution of parameter $`\lambda (u)`$. Define
$$p(x_1,t_1;x_2,t_2;u)dx_1dx_2=\text{Prob}\{x_1<Y(u,t_1)<x_1+dx_1;x_2<Y(u,t_2)<x_2+dx_2\}$$
to be the joint pdf of the process with $`t_2t_1`$. This pdf can be written as
$$p(x_1,t_1;x_2,t_2;u)=\underset{n_1=0}{\overset{\mathrm{}}{}}\underset{n_2=0}{\overset{\mathrm{}}{}}p(x_1,t_1;x_2,t_2;u|n_1,n_2)P(n_1,t_1;n_2,t_2;u),$$
(A2)
where $`p(x_1,t_1;x_2,t_2;u|n_1,n_2)`$ is the conditional pdf assuming that exactly $`n_1`$ pulses have occurred at time $`t_1`$ and $`n_2`$ pulses at time $`t_2`$. $`P(n_1,t_1;n_2,t_2;u)`$ is the joint probability for the occurrence of such pulses. Since $`t_2t_1`$ then $`n_2n_1`$ and
$$P(n_1,t_1;n_2,t_2;u)=\{\begin{array}{cc}P(n_2n_1;t_2t_1;u)P(n_1,t_1;u),\hfill & \text{if }n_1n_2\hfill \\ 0,\hfill & \text{otherwise,}\hfill \end{array}$$
(A3)
where
$$P(m,\tau ;u)=\frac{[\lambda (u)\tau ]^m}{m!}e^{\lambda (u)\tau }$$
(A4)
is the Poisson distribution. Substituting Eq. (A3) into Eq. (A2) and defining $`t_1=t`$, $`t_2=t+\mathrm{\Delta }t`$, $`n_1=n`$ and $`n_2n_1=m`$, we obtain
$$p(x_1,t;x_2,t+\mathrm{\Delta }t;u)=\underset{n=0}{\overset{\mathrm{}}{}}\underset{m=0}{\overset{\mathrm{}}{}}p(x_1,t;x_2,t+\mathrm{\Delta }t;u|n,n+m)P(m,\mathrm{\Delta }t;u)P(n,t;u),$$
and the characteristic function reads
$$\stackrel{~}{p}(\omega _1,t;\omega _2,t+\mathrm{\Delta }t;u)=\underset{n=0}{\overset{\mathrm{}}{}}\underset{m=0}{\overset{\mathrm{}}{}}\stackrel{~}{p}(\omega _1,t;\omega _2,t+\mathrm{\Delta }t;u|n,n+m)P(m,\mathrm{\Delta }t;u)P(n,t;u).$$
(A5)
Note that $`\stackrel{~}{p}(\omega _1,t;\omega _2,t+\mathrm{\Delta }t;u|n,n+m)`$ is the joint characteristic function of the truncated process
$$Y_n(u,t)=\underset{k=1}{\overset{n}{}}A_k(u)\varphi [tT_k(u);u].$$
Hence
$`\stackrel{~}{p}(\omega _1,t;\omega _2,t+\mathrm{\Delta }t;u|n,n+m)=\mathrm{exp}[i\omega _2{\displaystyle \underset{k=1}{\overset{n}{}}}A_k(u)\varphi (t+\mathrm{\Delta }t`$ $``$ $`T_k(u);u)+i\omega _1{\displaystyle }_{l=1}^nA_l(u)\varphi (tT_l(u);u)`$
$`+`$ $`i\omega _2{\displaystyle \underset{k=n+1}{\overset{n+m}{}}}A_k(u)\varphi (t+\mathrm{\Delta }tT_k(u);u)].`$
Taking into account that $`A_k(u)`$ and $`T_k(u)`$ are independent and identically distributed random variables, we have
$`\stackrel{~}{p}(\omega _1,t;\omega _2,t+\mathrm{\Delta }t;u|n,n+m)=`$ $`\left[\mathrm{exp}\{i\omega _2A_k(u)\varphi (t+\mathrm{\Delta }tT_k(u);u)+i\omega _1A_k(u)\varphi (tT_k(u);u)\}\right]^n`$
$`\times \left[\mathrm{exp}\{i\omega _2A_k(u)\varphi (t+\mathrm{\Delta }tT_k(u);u)\}\right]^m,`$
and, since the random times $`T_k(u)`$ are Poissonian,
$`\mathrm{exp}\{i\omega _2A_k(u)\varphi (t+\mathrm{\Delta }tT_k(u);u)`$ $`+`$ $`i\omega _1A_k(u)\varphi (tT_k(u);u)\}`$
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}h(a,u)𝑑a{\displaystyle _0^t}{\displaystyle \frac{dt^{}}{t}}\mathrm{exp}\{ia[\omega _2\varphi (t+\mathrm{\Delta }tt^{};u)+\omega _1\varphi (tt^{};u)]\}`$
$`=`$ $`{\displaystyle \frac{1}{t}}{\displaystyle _0^t}\stackrel{~}{h}[\omega _2\varphi (t+\mathrm{\Delta }tt^{};u)+\omega _1\varphi (tt^{};u);u]𝑑t^{},`$
where $`\stackrel{~}{h}(\omega ;u)`$ is the Fourier transform of the jump pdf $`h(x;u)`$. Analogously
$$\mathrm{exp}\{i\omega _2A_k(u)\varphi (t+\mathrm{\Delta }tT_k(u);u)\}=\frac{1}{\mathrm{\Delta }t}_t^{t+\mathrm{\Delta }t}\stackrel{~}{h}[\omega _2\varphi (t+\mathrm{\Delta }tt^{};u);u]𝑑t^{}.$$
Therefore,
$$\stackrel{~}{p}(\omega _1,t;\omega _2,t+\mathrm{\Delta }t;u|n,n+m)=\left[\frac{1}{t}F(\omega _2,t+\mathrm{\Delta }t;\omega _1,t;u)\right]^n\left[\frac{1}{\mathrm{\Delta }t}G(\omega _2,t+\mathrm{\Delta }t;t;u)\right]^m,$$
(A6)
where
$$F(\omega _2,t+\mathrm{\Delta }t;\omega _1,t;u)_0^t\stackrel{~}{h}[\omega _2\varphi (t+\mathrm{\Delta }tt^{};u)+\omega _1\varphi (tt^{};u);u]𝑑t^{},$$
(A7)
and
$$G(\omega _2,t+\mathrm{\Delta }t;t;u)_t^{t+\mathrm{\Delta }t}\stackrel{~}{h}[\omega _2\varphi (t+\mathrm{\Delta }tt^{};u);u]𝑑t^{}.$$
(A8)
Substituting Eq. (A6) into Eq. (A5) yields
$$\stackrel{~}{p}(\omega _1,t;\omega _2,t+\mathrm{\Delta }t;u)=\left\{\underset{n=0}{\overset{\mathrm{}}{}}\left[\frac{1}{t}F(\omega _1,t;\omega _2,t+\mathrm{\Delta }t;u)\right]^nP(n,t;u)\right\}\left\{\underset{m=0}{\overset{\mathrm{}}{}}\left[\frac{1}{\mathrm{\Delta }t}G(\omega _2,t+\mathrm{\Delta }t;t;u)\right]^mP(m,\mathrm{\Delta }t;u)\right\}.$$
Introducing Eq. (A4) into this and performing the resulting sums we get
$$\stackrel{~}{p}(\omega _1,t;\omega _2,t+\mathrm{\Delta }t;u)=\mathrm{exp}\left\{\lambda (u)\left[tF(\omega _2,t+\mathrm{\Delta }t;\omega _1,t;u)\right]\lambda (u)\left[\mathrm{\Delta }tG(\omega _2,t+\mathrm{\Delta }t;\omega _1,t;u)\right]\right\}.$$
Finally
$`\stackrel{~}{p}(\omega _1,t;\omega _2,t+\mathrm{\Delta }t;u)=\mathrm{exp}\{\lambda (u){\displaystyle _0^t}dt^{}[\stackrel{~}{h}[\omega _1\varphi (t^{},u)`$ $`+`$ $`\omega _2\varphi (t^{}+\mathrm{\Delta }t,u)]1]`$ (A9)
$`+`$ $`\lambda (u){\displaystyle _0^{\mathrm{\Delta }t}}dt^{}[\stackrel{~}{h}\left[\omega _2\varphi (t^{},u)\right]1]\},`$ (A10)
which agrees with Eq. (8). If in (A10) we set $`\omega _1=\omega `$ and $`\omega _2=0`$, we obtain the one time characteristic function (6):
$$\stackrel{~}{p}(\omega ,t;u)=\mathrm{exp}\left\{\lambda (u)_0^t𝑑t^{}\left[\stackrel{~}{h}\left[\omega \varphi (t^{},u)\right]1\right]\right\}.$$
(A11)
## B Cumulants for a step-like function
We will derive closed expressions for the cumulants (39):
$$X^{2n}(t)=\frac{(1)^n\stackrel{~}{h}^{(2n)}(0)}{2n}(bt)^{2n/\alpha }_0^{\mathrm{}}\frac{\varphi ^{2n}(x)}{x^{2n/\alpha }}𝑑x.$$
(B1)
when $`\varphi (x)`$ is the step-like function (35):
$$\varphi (t)=\{\begin{array}{cc}1e^{kt},\hfill & \text{if }t>0\hfill \\ 0,\hfill & \text{otherwise.}\hfill \end{array}$$
(B2)
The substitution of Eq. (B2) into Eq. (B1) leads us to evaluate the following integral
$`I_{2n}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{(1e^{kx})^{2n}}{x^{2n/\alpha }}}𝑑x`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dx}{x^{2n/\alpha }}}\left[kx{\displaystyle _0^1}e^{kxu}𝑑u\right]^{2n}`$
$`=`$ $`k^{2n}{\displaystyle _0^1}𝑑u_1\mathrm{}{\displaystyle _0^1}𝑑u_{2n}{\displaystyle _0^{\mathrm{}}}x^{2n(11/\alpha )}\mathrm{exp}\left\{kx\left({\displaystyle \underset{i=1}{\overset{2n}{}}}u_i\right)\right\}𝑑x.`$
Define the new integration variable
$$\xi =kx\left(\underset{i=1}{\overset{2n}{}}u_i\right),$$
then
$$I_{2n}=k^{1+2n/\alpha }_0^1𝑑u_1\mathrm{}_0^1𝑑u_{2n}\left(\underset{i=1}{\overset{2n}{}}u_i\right)^{12n(11/\alpha )}_0^{\mathrm{}}e^\xi \xi ^{2n(11/\alpha )}𝑑\xi .$$
But
$$_0^{\mathrm{}}e^\xi \xi ^{2n(11/\alpha )}𝑑\xi =\mathrm{\Gamma }(1+2n2n/\alpha ).$$
Defining the numbers
$$A_{2n}_0^1𝑑u_1\mathrm{}_0^1𝑑u_{2n}\left(\underset{i=1}{\overset{2n}{}}u_i\right)^{12n(11/\alpha )},$$
(B3)
we finally have
$$X^{2n}(t)=\frac{(1)^n\stackrel{~}{h}^{(2n)}(0)}{2n}A_{2n}k^{1+2n/\alpha }\mathrm{\Gamma }(1+2n2n/\alpha )(bt)^{2n/\alpha },$$
(B4)
which is Eq. (46).
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# Knowledge on Treelike Spaces
## 1 Introduction
The notion of possible world dominates the literature in modal logic, via Kripke models, as well as in any logic dealing with the epistemic state of a reasoner. The heart of this popularity lies in the identification of an intentional state through common properties of extensional objects. Apart from genuine problems such as logical omniscience this representation suffers from, it is limited in a static description of the reasoner’s epistemic state. The “logic of knowing” is not only embodied in the representation of knowledge but also in the way knowledge is acquired. We do not refer to temporal properties but rather to methodology (though both can be intertwined).
Recently a family of logics was introduced (\[MP92\],\[Geo94a\],\[Geo93\],\[DMP\]) with the intention to fill this void. It succeeds in doing so by attaching familiar mathematical structures such as spaces of subsets, topologies and complete lattices of subsets corresponding to a natural knowledge acquisition. This paper extends this work by introducing a bimodal logic belonging to the same family of logics and establishes a correspondence between a particular epistemic process of knowledge acquisition with a space of subsets forming a tree (treelike space).
In our framework the view of a reasoner will be represented by a set of possible worlds. Each of these worlds represents an alternative state compatible with the reasoner’s knowledge of actual state. This treatment of knowledge agrees with the traditional one (\[Hin62\], \[HM84\], \[PR85\], \[CM86\]) expressed in a variety of contexts (artificial intelligence, distributed processes, economics, etc).
We are interested in formulating a basic logical framework for reasoning about a resource-conscious acquiring of knowledge. Such a framework can be applied to many settings such as the ones involving time, computation, physical experiments or observations. In these settings an (discrete or continuous) increase of information available to us takes place and results in an increase of our knowledge. How could this simple idea be embodied in the formentioned semantical framework? An increase of knowledge can be represented with a restriction of the knower’s view, i.e. of the equivalence class of the alternative worlds. This restriction is nondeterministic (we do not know what kind of additional information will be available to us, if at all) but not arbitrary: it will always contain the actual state of the knower, i.e. it is a neighborhood restriction of the actual state. In this way, set-theoretic considerations come in.
A discrete version of our epistemic framework can arise in scientific experiments or tests. We acquire knowledge by “a step-by-step” process, each step being an experiment or test. The outcome of such an experiment or test is unknown to us beforehand, but after being known it restricts our attention to a smaller set of possibilities. A sequence of experiments, tests, or actions comprises a strategy of knowledge acquisition. This model is in many respects similar to Hintikka’s “oracle” (see \[Hin86\]). In Hintikka’s model the “inquirer” asks a series of questions $`Q_1,Q_2,\mathrm{},Q_n,\mathrm{}`$ to an external information source, called “oracle” (can be thought as a knowledge base). The oracle answers yes or no and the inquirer increases his or her knowledge by this piece of additional evidence. At any point of this process the inquirer follows a branch of a tree determined by the possible answers to his or her series of questions. Such an interrogative model is recognized by Gadamer (\[Gad75\]) as an important part of the epistemic process. Consider the following example:
Example: Suppose that our view, the set of possible worlds, is $`\{q_1,q_2,q_3,q_4\}`$ and our query consists of two questions $`Q_1`$, $`Q_2`$, in that order. The answer to $`Q_1`$ is yes in $`q_1`$, $`q_2`$ and no in $`q_3`$, $`q_4`$. The answer to $`Q_2`$ is yes in $`q_1`$, $`q_2`$, $`q_3`$ and no in $`q_4`$. Then the possible sequences of knowledge states comprise a tree of subsets as shown in Figure 1. The space of subsets labeling the nodes of the tree will be called a treelike space.
The above example shows a transition from the symbolic description of the epistemic process to a description in spatial terms. Instead of going down a proof tree, the one which entails the desired formulae, we intersect nodes of a tree labeled by subsets of a space. This transition is direct; it enables us to think in geometric terms.
Now consider the following example:
Example: Suppose that a machine emits a stream of binary digits representing the output of a recursive function $`f`$. After time $`t_1`$ the machine emitted the stream $`111`$. The only information we have about the function being computed at this time on the basis of this (finite) observation is that
$$f(1)=f(2)=f(3)=1.$$
As far as our knowledge concerns, $`f`$ is indistinguishable from the constant function $`\mathrm{𝟏}`$, where $`\mathrm{𝟏}(n)=1`$ for all $`n`$. After some additional time $`t_2`$, i.e. spending more time and resources, $`0`$ might appear and thus we could be able to distinguish $`f`$ from $`\mathrm{𝟏}`$. In any case, each binary stream will be an initial segment of $`f`$ and this initial segment is a neighborhood of $`f`$. In this way, we can acquire better knowledge of the function the machine computes. The space of finite binary streams is a structure which models computation. The sets of binary streams under the initial segment ordering is an example of a treelike space.
The above example shows how the same epistemic process appears during observations of programs. Here possible worlds correspond to (total) computations and our view to observations. We can apply the same spatial reasoning to programs through the following correspondence:
| Knowledge states | $`=`$ | Sets | $`=`$ | Observations |
| --- | --- | --- | --- | --- |
| Possible worlds | $`=`$ | Points | $`=`$ | Computations. |
Therefore a common idea lies behind the knowledge-theoretic, spatial and computational framework. The connection between the last two is not new. Here is how this epistemic framework ties with previous work on establishing links between spatial reasoning and reasoning about programs.
We use two modalities $`𝖪`$ for knowledge and $`\mathrm{}`$ for effort, i.e. spending of resources. Consider the formula
$$A\mathrm{}𝖪A,$$
where $`A`$ is an atomic predicate and $`\mathrm{}`$ is the dual of the $`\mathrm{}`$, i.e. $`\mathrm{}\neg \mathrm{}\neg `$. It will be clear after the presentation of semantics in Section 2.1 that if the above formula is valid, then the set which $`A`$ represents is an open set of the topology generated by the subsets of the treelike space as a basis. Under the reading of $`\mathrm{}`$ as “possible” and $`𝖪`$ as “is known”, the above formula says that
> “if $`A`$ is true then it is possible for $`A`$ to be known”,
i.e. $`A`$ is affirmative. Vickers defines similarly an affirmative assertion in \[Vic89\]
> “an assertion is affirmative iff it is true precisely in the circumstances when it can be affirmed.”
Affirmative and refutative assertions are closed under infinite disjunctions and conjunctions, respectively. Smyth in \[Smy83\] observed first these properties in semi-decidable properties. Semi-decidable properties are those properties whose truth set is r.e. and are a particular kind of affirmative assertions. In fact, changing our power of affirming or computing we get another class of properties with a similar knowledge-theoretic character. For example, using polynomial algorithms affirmative assertions become polynomially semi-decidable, i.e. NP properties. If an object has this property then it is possible to know it with a polynomial algorithm even though it is not true we know it now.
Our approach has an independent theoretical interest. A new family of Kripke frames, called subset frames, arises. These are Kripke frames which are equivalent to sets of subsets. In particular, we have identified those which are equivalent to (complete) lattices of subsets and topologies (see \[Geo93\]). In this paper, we shall identify those which correspond to the above interrogative model, called treelike spaces. Treelike spaces have a particular interest; they correspond to an indeterminist’s theory of time called Ockhamism (see \[Pri67\]), which gives rise to branching time. We refer the reader to section 2.1 for a detailed discussion.
A family of logics for knowledge and time is studied in \[HV89\] and various complexity results are established. However, the framework of the above logics is restricted to distributed systems and their interpretation differs significantly from ours.
Interpreting the knowledge modal operator as a universal quantifier we present a novel way of understanding the meaning of quantifiers in varying (ordered) domains (see section 2.2 for a relevant discussion). This is one of the main difficulties in formulating a meaningful first-order system for modal logic (see \[Fit93\] for a discussion).
The language and semantics of our logical framework is presented in Section 2. In the same section, we present two systems which belong to the same family of logics, studied in \[MP92\], \[Geo93\] and \[Geo94a\]. In Section 3, we present an axiomatization, called $`\mathrm{𝐌𝐏𝐓}`$, for our semantics and we prove completeness, small model property, and decidability.
A preliminary version of this paper has appeared in \[Geo94b\].
## 2 Two Systems: $`\mathrm{𝐌𝐏}`$ and $`\mathrm{𝐌𝐏}^{}`$
### 2.1 Language and Semantics
We follow the notation of \[MP92\].
We construct a bimodal propositional modal logic. Formally, we start with a countable set $`𝖠`$ of atomic formulae, then the language $``$ is the least set such that $`𝖠`$ and closed under the following rules:
$$\frac{\varphi ,\psi }{\varphi \psi }\frac{\varphi }{\neg \varphi ,\mathrm{}\varphi ,𝖪\varphi }$$
We abbreviate, as usual, $`\varphi \neg \varphi `$ with $``$ and $`\neg `$ with $``$. The language $``$ can be interpreted inside any spatial context as follows.
Definition 1 Let $`X`$ be a set and $`𝒪`$ a subset of the powerset of $`X`$, i.e. $`𝒪𝒫(X)`$ such that $`X𝒪`$. We call the pair $`X,𝒪`$ a subset space. A model is a triple $`X,𝒪,i`$, where $`X,𝒪`$ is a subset space and $`i`$ a map from $`𝖠`$ to $`𝒫(X)`$ with $`i()=X`$ and $`i()=\mathrm{}`$ called initial interpretation.
We denote the set $`\{(x,U):U𝒪,\text{ and }xU\}X\times 𝒪`$ with $`X\dot{\times }𝒪`$. For each $`U𝒪`$ let $`U`$ be the lower closed set generated by $`U`$ in the partial order $`(𝒪,)`$, i.e. the set $`\{V:V𝒪\text{ and }VU\}`$.
Definition 2 The satisfaction relation $`_{}`$, where $``$ is the model $`X,𝒪,i`$, is a subset of $`(X\dot{\times }𝒪)\times `$ defined recursively by (we write $`x,U_{}\varphi `$ instead of $`((x,U),\varphi )_{}`$)
$$\begin{array}{cc}x,U_{}A\hfill & \text{iff}xi(A),\text{ where }A𝖠\hfill \\ x,U_{}\varphi \psi \hfill & \text{if}x,U_{}\varphi \text{ and }x,U_{}\psi \hfill \\ x,U_{}\neg \varphi \hfill & \text{if}x,U\vDash ̸_{}\varphi \hfill \\ x,U_{}𝖪\varphi \hfill & \text{if}\text{for all }yU,y,U_{}\varphi \hfill \\ x,U_{}\mathrm{}\varphi \hfill & \text{if}\text{for all }VU\text{ such that }xV,x,V_{}\varphi .\hfill \end{array}$$
If $`x,U_{}\varphi `$, for all $`(x,U)`$ belonging to $`X\dot{\times }𝒪`$, then $`\varphi `$ is valid in $``$, denoted by $`\varphi `$.
The case for atomic formulae shows that we deal with analytic sentences, i.e. sentences which do not change their truth value. If a formula $`\mathrm{}\varphi `$ does not contain $`𝖪`$ then it has the same interpretation as $`\varphi `$. This has also the consequence that the universal substitution rule does not hold. Thus, time does not affect the semantic value of sentences but rather the knowledge we have of them. This difference makes the $`\mathrm{}`$ modality not collapsing to a temporal modality but being closer to necessity.
We abbreviate $`\neg \mathrm{}\neg \varphi `$ and $`\neg 𝖪\neg \varphi `$ with $`\mathrm{}\varphi `$ and $`𝖫\varphi `$ respectively. We have that
$$\begin{array}{cc}x,U_{}𝖫\varphi \hfill & \text{if there exists }yU\text{ such that }y,U_{}\varphi \hfill \\ x,U_{}\mathrm{}\varphi \hfill & \text{if there exists }V𝒪\text{ such that }VU,xV,\text{ and }x,V_{}\varphi .\hfill \end{array}$$
Definition 3 A treelike space is a subset space $`X,𝒪`$ where for all $`U,V𝒪`$, either $`UV`$, or $`VU`$, or $`UV=\mathrm{}`$. A model induced by a tree space will be called a treelike model.
It is clear that in the countable case the set of subsets of a treelike space forms a tree under the subset ordering.
Example: Let
$$X=\{ff\text{ recursive }\}.$$
Now, let
$$[a_1,a_2,\mathrm{},a_n]=\{ff(k)=a_k,\text{ for }k=1,2,\mathrm{},n\}X,$$
where $`a_1,a_2,\mathrm{},a_n`$ are natural numbers, and
$$𝒪=\{[a_1,a_2,\mathrm{},a_n]n=1,2,\mathrm{}\}\{X\}.$$
Then it is easily verified, using definition 2.1, that $`X,𝒪`$ is a treelike space.
Now let $`\mathrm{𝟏}`$ be a predicate with
$$i(\mathrm{𝟏})=\{f\text{ there exists }n\text{ such that for all }m>n,f(n)=1\}.$$
Then the formula
$$\mathrm{}𝖫\mathrm{𝟏}$$
which translates to “it will never be known that $`0`$ appears infinitely often”, is valid in the treelike model $`X,𝒪,i`$. This comes with no surprise, since the knowledge of “infinitely often” requires an infinite amount of resources. This formula is an example of a refutative assertion (see introduction).
Treelike spaces get their name from treelike frames (see \[Pri67\]). A treelike frame is a pair $`T,<`$, where $`T`$ is a nonempty set and $`<`$ is a transitive ordering on $`T`$ such that if $`t_1<t`$ and $`t_2<t`$ then either $`t_1=t`$ or $`t_1<t_2`$ or $`t_2<t_1`$. Treelike frames have appeared as semantics for the Ockhamist’s concept of non-deterministic time and been used for treating historical necessity and conditionals (see \[Tho84\] and \[VF81\]). The validity on these frames is called Ockhamist validity. A treelike space is a special form of a treelike frame where the temporal instants of the frame are labeled by subsets of a space and whenever instants are incomparable the respective subsets are disjoint. It can be easily seen that the ordering among subsets is a treelike frame. The similarities do not end here. Let $`T,<`$ be a treelike frame and, for each $`tT`$, $`B_t`$ the set of maximal linear ordered subsets of $`T`$ containing $`t`$, i.e. the branches intersecting $`t`$. Then $`\{B_t\}_{tT}`$ is a treelike space. The difference lies on the interpretation of atomic formulae. We interpret atomic formulae on branches while an Ockhamist assignment interprets atomic formulae on temporal instances. This bring up another dimension of our logic. Our logic is not conservative over a logic which interprets $`\mathrm{}`$ as $`𝖥`$ (the “future” modality) for if $`\varphi `$ contains no occurrences of $`𝖪`$ then $`\mathrm{}\varphi `$ is valid in a treelike space exactly when $`\varphi `$ is. We adopt the indeterminist’s view of necessity (knowledge). Although $`\varphi `$ may be true in our world, $`𝖪\varphi `$ may be false. This is because there is no special world in our view which deserves to be called actual. Setting apart Ockhamist validity, treelike spaces are more general than treelike frames (and their derivative $`T\times W`$ frames) due to the fact that we do not assume an overall temporal ordering. In this sense treelike spaces are closer to a more general structure, first introduced by Kamp and subsequently called Kamp frames, where worlds do not participate in the same temporal structure (for definition and discussion see \[Tho84\]). In fact, it is easily seen that treelike spaces are equivalent to Ockhamist frames introduced by Zanardo in \[Zan85\] for the completeness of strong Ockhamist validity. At any rate, our work seems to have more than superficial links with work in historical necessity and questions such as what the connections between the two notions of validity are should be the subject of a more systematic investigation.
### 2.2 $`\mathrm{𝐌𝐏}`$ and $`\mathrm{𝐌𝐏}^{}`$
We saw that the semantics of the bimodal language is interpreted in any pair $`X,𝒪`$. What happens when we allow $`𝒪`$ to be any class of sets of subsets? If $`𝒪`$ is an arbitrary set of subsets then the system $`\mathrm{𝐌𝐏}`$ is complete for such subset spaces. The axiom system $`\mathrm{𝐌𝐏}`$ consists of axiom schemes 1 through 10 and rules of Table 1 (see page 1) and appeared first in \[MP92\].
The following was proved in \[MP92\].
###### Theorem 4
The axioms and rules of $`\mathrm{𝐌𝐏}`$ are sound and complete with respect to subset spaces.
If $`𝒪`$ is a complete lattice under set-theoretic union and intersection then the system $`\mathrm{𝐌𝐏}^{}`$ is canonically complete for this class of subset spaces. The axiom system $`\mathrm{𝐌𝐏}^{}`$ consists of the axiom schemes and rules of $`\mathrm{𝐌𝐏}`$ plus the following two additional axiom schemes:
$$\mathrm{}\mathrm{}\varphi \mathrm{}\mathrm{}\varphi $$
and
$$\mathrm{}(𝖪\varphi \psi )𝖫\mathrm{}(𝖪\varphi \chi )\mathrm{}(𝖪\mathrm{}\varphi \mathrm{}\psi 𝖫\mathrm{}\chi ).$$
The first axiom is a well-known formula which characterizes incestual frames, i.e. if two points $`\beta `$ and $`\gamma `$ in a frame can be accessed by a common point $`\alpha `$ then there is a point $`\delta `$ which can be accessed by both $`\beta `$ and $`\gamma `$. The second characterizes union.
The following was proved in \[Geo93\].
###### Theorem 5
The axioms and rules of $`\mathrm{𝐌𝐏}^{}`$ are sound and canonically complete with respect to subset spaces, which are complete lattices.
The proof of the above theorem was later shortened and improved through an elegant embedding of $`\mathrm{𝐒𝟒}`$ (and therefore intuitionistic logic via the Gödel translation) by Dabrowski, Moss and Parikh in \[DMP\]. This translation reveals that truth in intuitionistic logic coincides with “possibility of knowing” in our system. It also reveals a connection with another line of work, that of Fischer Servi. In \[FS80\] and \[FS84\] the semantics and syntax of the family $``$-IC of intuitionistic modal logics is studied. This family is is naturally embedded via the Gödel translation to the family ($`\mathrm{𝐒𝟒}`$-$``$) of bimodal logics, where $`\mathrm{𝐒𝟒}`$ is always one of the coordinates (like in our case). However, the semantics called double model structures (birelational modal frames) deviate from our space theoretic framework; a fact that declares itself on the presence of different connecting axioms, i.e. axioms involving both modalities.
## 3 The system $`\mathrm{𝐌𝐏𝐓}`$
We add the axioms 11 and 12 to form the system $`\mathrm{𝐌𝐏𝐓}`$ for the purpose of axiomatizing treelike spaces.
A word about the axioms (most of the following facts can be found in any introductory book about modal logic, e.g. \[Che80\] or \[Gol87\].) Axiom 2 expresses the fact that the truth of atomic formulae is independent of the choice of subset and depends only on the choice of point. Axioms 3 through 5 and Axioms 6 through 9 are used to axiomatize the normal modal logics S4 and S5 respectively. The former group of axioms expresses the fact that the passage from one subset to its restriction is done in a constructive way, as actually happens in an increase of information or a spending of resources (the classical interpretation of necessity in intuitionistic logic is axiomatized in the same way). The latter group is generally used for axiomatizing logics of knowledge.
Axiom 10 expresses the fact that if a formula holds in arbitrary subsets is going to hold as well in the ones which are neighborhoods of a point. The converse of this axiom is not sound.
Axiom 11 is a well-known axiom which characterizes reflexive, transitive and connected frames, i.e. if two points $`\beta `$ and $`\gamma `$ in a frame can be accessed by a common point $`\alpha `$ then either $`\beta `$ accesses $`\gamma `$ or $`\gamma `$ accesses $`\beta `$ (or both).
Soundness of Axioms 1 through 10 has already been established for arbitrary subset spaces (see \[MP92\]). The soundness of Axiom 11 is easy to see, since the subset frame (see \[Geo93\]), i.e. the birelational modal frame, of a tree model is connected.
###### Proposition 6
The axiom 12 is sound.
Proof. We shall show soundness for the equivalent formula
$$\mathrm{}𝖪\varphi \mathrm{}𝖫(\psi \mathrm{}\varphi )𝖫(\mathrm{}\psi \mathrm{}\varphi ).$$
Let $`x,U\mathrm{}𝖪\varphi \mathrm{}𝖫(\psi \mathrm{}\varphi )`$. Then there exists $`VU`$ such that $`x,V𝖫(\psi \mathrm{}\varphi )`$. This implies that there exists $`yV`$ such that $`y,V\psi \mathrm{}\varphi `$. Now, observe that $`y,U\mathrm{}\varphi `$. For, if $`WU`$ and $`yW`$ then there are two cases. Either $`WV`$ and $`y,W\varphi `$, since $`y,U\mathrm{}\varphi `$, and we are done, or $`WU`$ and $`WV`$ so we have $`VWU`$, since the subsets containing $`y`$ are linearly ordered. In this case, we have $`xW`$, since $`xV`$. By our assumption $`x,U\mathrm{}𝖪\varphi `$, we have $`x,W𝖪\varphi `$. So $`y,W\varphi `$. Now, $`yU`$ and $`y,U\mathrm{}\varphi `$ imply together $`y,U\mathrm{}\psi \mathrm{}\varphi `$.
Note that Axiom 10 follows from Axiom 12 (substitute $`\varphi `$ with $``$). Axiom 10 has a particular interest; if we replace $`𝖪`$ with the universal quantifier it becomes the well-known Barcan formula
$$x\mathrm{}\varphi (x)\mathrm{}x\varphi (x).$$
Our system (and therefore $`\mathrm{𝐌𝐏}`$ and $`\mathrm{𝐌𝐏}^{}`$, since this formula belongs to their axiomatization) can be thought as a propositional analogue of a first order modal system interpreted over varying restricting domains (see \[Fit93\]).
### 3.1 Completeness
Our proof of completeness is based on a construction of a treelike model which is (strongly) equivalent to each generated canonical submodel of the canonical model of $`\mathrm{𝐌𝐏𝐓}`$.
The canonical model of $`\mathrm{𝐌𝐏𝐓}`$ is the structure
$$𝒞=(S,\{\stackrel{\mathrm{}}{},\stackrel{\text{L}}{}\},v),$$
where
$$\begin{array}{ccc}& S=\{s|s\text{ is }\mathrm{𝐌𝐏𝐓}\text{-maximal consistent}\},\hfill & \\ & s\stackrel{\mathrm{}}{}t\text{ iff }\{\varphi |\mathrm{}\varphi s\}t,\hfill & \\ & s\stackrel{\text{L}}{}t\text{ iff }\{\varphi |𝖪\varphi s\}t,\hfill & \\ & v(A)=\{sS|As\},\hfill & \end{array}$$
along with the usual satisfaction relation (defined inductively):
$$\begin{array}{ccc}s_𝒞A\hfill & \text{iff}\hfill & sv(A)\hfill \\ s_𝒞\neg \varphi \hfill & \text{iff}\hfill & s\vDash ̸_𝒞\varphi \hfill \\ s_𝒞\varphi \psi \hfill & \text{iff}\hfill & s_𝒞\varphi \text{ and }s_𝒞\psi \hfill \\ s_𝒞\mathrm{}\varphi \hfill & \text{iff}\hfill & \text{for all }tS,s\stackrel{\mathrm{}}{}t\text{ implies }t_𝒞\varphi \hfill \\ s_𝒞𝖪\varphi \hfill & \text{iff}\hfill & \text{for all }tS,s\stackrel{\text{L}}{}t\text{ implies }t_𝒞\varphi .\hfill \end{array}$$
We write $`𝒞\varphi `$, if $`s_𝒞\varphi `$ for all $`sS`$.
A canonical model exists for all consistent bimodal systems with the normal axiom scheme for each modality (as $`\mathrm{𝐌𝐏𝐓}`$). We have the following well known theorems (see \[Che80\], or \[Gol87\]).
###### Theorem 7 (Truth Theorem)
For all $`sS`$ and $`\varphi `$,
$$s_𝒞\varphi \text{iff}\varphi s.$$
###### Theorem 8 (Completeness Theorem)
For all $`\varphi `$,
$$𝒞\varphi \text{iff}_{\mathrm{𝐌𝐏𝐓}}\varphi .$$
We shall now prove some properties of $`𝒞`$.
###### Proposition 9
1. The canonical frame is reflexive, transitive and connected with respect to the relation $`\stackrel{\mathrm{}}{}`$.
2. The relation $`\stackrel{\text{L}}{}`$ is an equivalence relation.
3. For all $`s,s^{},tS`$, if $`s\stackrel{\mathrm{}}{}s^{}\stackrel{\text{L}}{}t`$ then there exists $`t^{}S`$ such that $`s\stackrel{\text{L}}{}t^{}\stackrel{\mathrm{}}{}t`$.
4. For all $`s,s^{}S`$, if $`s\stackrel{\text{L}}{}s^{}`$ and $`s\stackrel{\mathrm{}}{}s^{}`$ then $`s=s^{}`$.
5. The relation $`\stackrel{\mathrm{}}{}`$ is antisymmetric.
Proof. For Part a, Axioms 3 through 5 and Axiom 11 characterize reflexive, transitive and connected frames (these axioms comprise the system $`\mathrm{𝐒𝟒}\mathbf{.3}`$).
For Part b, $`𝖪`$ is axiomatized with the $`\mathrm{𝐒𝟓}`$ axioms.
Part c is an immediate consequence of Axiom 10.
To show Part d, let
$$I_{s,s^{}}=\{ts\stackrel{\mathrm{}}{}t\stackrel{\mathrm{}}{}s^{}\},$$
for all pairs $`(s,s^{})`$ such that $`s\stackrel{\text{L}}{}s^{}`$ and $`s\stackrel{\mathrm{}}{}s^{}`$.
We shall prove by induction on the complexity of $`\varphi `$ that, for all pairs $`(s,s^{})`$ such that $`s\stackrel{\text{L}}{}s^{}`$ and $`s\stackrel{\mathrm{}}{}s^{}`$, $`\varphi `$ belongs to some $`tI_{s,s^{}}`$ if and only if $`\varphi `$ belongs to s. This shows that $`I_{s,s^{}}s`$. Further, we have $`I_{s,s^{}}=\{s\}`$, since $`sI_{s,s^{}}`$. Therefore $`s=s^{}`$.
If $`\varphi `$ is an atomic formula $`A`$ and $`At`$, for some $`tI_{s,s^{}}`$, then $`\mathrm{}As`$. Therefore, by axiom 2, $`\mathrm{}As`$. Hence, $`As`$.
The cases of negation and conjunction are straightforward.
If $`\varphi =\mathrm{}\psi `$, let $`\mathrm{}\psi t`$, for some $`tI_{s,s^{}}`$. In particular, $`\psi t`$ and by induction hypothesis, $`\psi s`$. Suppose, towards a contradiction, that $`\mathrm{}\neg \psi s`$. Then there exists $`rS`$ such that $`s\stackrel{\mathrm{}}{}r`$ and $`\neg \psi r`$. Since the frame is connected, $`s\stackrel{\mathrm{}}{}s^{}`$ and $`s\stackrel{\mathrm{}}{}r`$ imply that either $`s^{}\stackrel{\mathrm{}}{}r`$ or $`r\stackrel{\mathrm{}}{}s^{}`$. If $`s^{}\stackrel{\mathrm{}}{}r`$ then $`t\stackrel{\mathrm{}}{}r`$ which is contradiction, since $`\mathrm{}\psi t`$ and $`\neg \psi r`$. If $`r\stackrel{\mathrm{}}{}s^{}`$ then, by induction hypothesis, $`\neg \psi s`$ which is a contradiction, since $`\psi s`$ and $`s`$ is consistent. Hence $`\mathrm{}\psi s`$.
If $`\varphi =𝖪\psi `$, let $`𝖪\psi t`$ for some $`tI_{s,s^{}}`$. Suppose, towards a contradiction, that $`𝖫\neg \psi s`$. Then there exists $`rS`$ such that $`s\stackrel{\text{L}}{}r`$ and $`\neg \psi r`$. We have $`s^{}\stackrel{\text{L}}{}r`$, since $`s^{}\stackrel{\text{L}}{}s`$. Since $`t\stackrel{\mathrm{}}{}s^{}`$, there exists, by Part c, $`r^{}S`$ such that $`t\stackrel{\text{L}}{}r^{}\stackrel{\mathrm{}}{}r`$. We have $`\psi r^{}`$, since $`𝖪\psi t`$. Since $`s\stackrel{\mathrm{}}{}t\stackrel{\text{L}}{}r^{}`$, there exists, by Part c, $`r^{\prime \prime }S`$ such that $`s\stackrel{\text{L}}{}r^{\prime \prime }\stackrel{\mathrm{}}{}r^{}`$. Notice that $`r^{\prime \prime }\stackrel{\mathrm{}}{}r^{}\stackrel{\mathrm{}}{}r`$ and $`r^{\prime \prime }\stackrel{\text{L}}{}r`$, and so $`r^{\prime \prime },r^{},rI_{r^{\prime \prime },r}`$. By our previous assumption, we have $`\neg \psi r`$ and $`\psi r^{}`$. By induction hypothesis on $`I_{r^{\prime \prime },r}`$, both $`\neg \psi `$ and $`\psi `$ should belong to $`r^{\prime \prime }`$ which is a contradiction to its consistency.
For Part e, we shall prove by induction on the structure of $`\varphi `$ that, for all $`s,tS`$ such that $`s\stackrel{\mathrm{}}{}t\stackrel{\mathrm{}}{}s`$, $`\varphi s`$ if and only $`\varphi t`$.
The cases of atomic formula, negation, conjunction and $`\mathrm{}`$ are straightforward. We shall show the $`\varphi =𝖪\psi `$ step. Let $`𝖪\psi s`$, and suppose $`𝖫\neg \psi t`$ towards a contradiction. Then there exists $`rS`$ such that $`t\stackrel{\text{L}}{}r`$ and $`\neg \psi r`$. Since $`s\stackrel{\mathrm{}}{}t\stackrel{\text{L}}{}r`$, there exists $`pS`$ such that $`s\stackrel{\text{L}}{}p\stackrel{\mathrm{}}{}r`$. Also, $`\psi p`$, since $`𝖪\psi S`$. Now, since $`t\stackrel{\mathrm{}}{}s\stackrel{\text{L}}{}p`$ there exists $`r^{}S`$ such that $`t\stackrel{\text{L}}{}r^{}\stackrel{\mathrm{}}{}p`$. This implies $`r^{}\stackrel{\mathrm{}}{}p\stackrel{\mathrm{}}{}r`$ and $`r\stackrel{\text{L}}{}r^{}`$. Therefore, by Part d, $`r=r^{}`$. Thus we have $`r\stackrel{\mathrm{}}{}p\stackrel{\mathrm{}}{}r`$ with $`\neg \psi r`$ and $`\psi p`$ which is a contradiction to the induction hypothesis.
The canonical model is not a (model corresponding to a) treelike model. A counterexample will appear later on (see Figure 2). However, by defining a number of equivalence relations, we shall be able to construct a treelike model equivalent to each generated part of the canonical model.
For all $`tS`$, let $`[t]=\{sSs\stackrel{\text{L}}{}t\}`$, i.e. the equivalence class under $`\stackrel{\text{L}}{}`$ where $`t`$ belongs. Let $`𝒞_\text{K}=\{[t]tS\}`$. We define the following relation on $`𝒞_\text{K}`$.
$$[t_1][t_2]\text{iff}\text{there exist}s_1,s_2S\text{such that}s_1[t_1],s_2[t_2]\text{and}s_2\stackrel{\mathrm{}}{}s_1.$$
###### Proposition 10
The relation $``$ is a partial order.
Proof. Since $`t\stackrel{\mathrm{}}{}t`$, we have $`[t][t]`$ and reflexivity follows.
For antisymmetry, let $`[t_1][t_2]`$ and $`[t_2][t_1]`$ for some $`t_1,t_2S`$. Then there exist $`s_1,s_2,s_1^{},s_2^{}S`$ such that $`s_1,s_1^{}[t_1]`$, $`s_2,s_2^{}[t_2]`$, $`s_2\stackrel{\mathrm{}}{}s_1`$ and $`s_1^{}\stackrel{\mathrm{}}{}s_2^{}`$. Since $`s_2\stackrel{\mathrm{}}{}s_1\stackrel{\text{L}}{}s_1^{}`$, there exists $`s_2^{\prime \prime }S`$ such that $`s_2\stackrel{\text{L}}{}s_2^{\prime \prime }\stackrel{\mathrm{}}{}s_1^{}`$. So we have $`s_2^{\prime \prime }\stackrel{\mathrm{}}{}s_1^{}\stackrel{\mathrm{}}{}s_2^{}`$ and $`s_2^{\prime \prime }\stackrel{\text{L}}{}s_2^{}`$ which implies, by Proposition 9(d), $`s_2^{\prime \prime }=s_2^{}`$. Therefore $`s_1^{}=s_2^{}`$, by $`\stackrel{\mathrm{}}{}`$’s antisymmetry. Hence $`[t_1]=[s_1^{}]=[s_2^{}]=[t_2]`$.
For transitivity, let $`[t_3][t_2][t_1]`$ for some $`t_1,t_2,t_3S`$. Then there exist $`s_1[t_1]`$, $`s_2,s_2^{}[t_2]`$, and $`s_3[t_3]`$ such that $`s_1\stackrel{\mathrm{}}{}s_2`$ and $`s_2^{}\stackrel{\mathrm{}}{}s_3`$. Since $`s_1\stackrel{\mathrm{}}{}s_2\stackrel{\text{L}}{}s_2^{}`$, there exists $`s_1^{}S`$ such that $`s_1\stackrel{\text{L}}{}s_1^{}\stackrel{\mathrm{}}{}s_2^{}`$. So $`s_1^{}\stackrel{\mathrm{}}{}s_3`$, and therefore $`[t_3]=[s_3][s_1^{}]=[t_1]`$.
A subset $`X`$ of $`S`$, the domain of the canonical model $`𝒞`$, is called $`𝖪\mathrm{}`$-closed whenever
$$\text{ if }sX,\text{ and }s\stackrel{\mathrm{}}{}t\text{ or }s\stackrel{\text{L}}{}t,\text{then}tX.$$
The intersection of $`𝖪\mathrm{}`$-closed sets is still $`𝖪\mathrm{}`$-closed, therefore we can define the smallest $`𝖪\mathrm{}`$-closed containing $`t`$, for all $`tS`$. We shall denote this set by $`S^t`$. Fix $`t_0S`$. We define the model
$$𝒞^{t_0}=(S^{t_0},\stackrel{\mathrm{}}{}|_{S^{t_0}\times S^{t_0}},\stackrel{\text{L}}{}|_{S^{t_0}\times S^{t_0}},v^{t_0}),$$
where $`\stackrel{\mathrm{}}{}|_{S^{t_0}\times S^{t_0}}`$, $`\stackrel{\text{L}}{}|_{S^{t_0}\times S^{t_0}}`$ and $`v^{t_0}`$ are the restrictions of $`\stackrel{\mathrm{}}{}`$, $`\stackrel{\text{L}}{}`$ and $`v`$ to $`S^{t_0}\times S^{t_0}`$ and $`S^{t_0}`$ respectively. We shall call this model the submodel of $`𝒞`$ generated by $`t_0`$.
Observe that if we restrict the partial order $``$ to $`𝒞^{t_0}`$ then $`[t_0]`$ is the greatest element under $``$.
For each generated submodel of the canonical model, we shall construct a treelike model which is equivalent to it.
For each $`sS^{t_0}`$, let
$$[[s]]=\{t[t_0]\text{there exists}t^{}[s]\text{such that}t\stackrel{\mathrm{}}{}t^{}\}.$$
Notice that $`[[s]][t_0]`$.
For each $`sS^{t_0}`$, we define the following relation $`_s`$ on $`[[s]]`$
$$t_1_st_2\text{iff}\text{for all}[s][s^{}],t_1[[s^{}]]\text{iff}t_2[[s^{}]].$$
###### Proposition 11
For all $`sS^{t_0}`$, the relation $`_s`$ is an equivalence relation.
Proof. This is because $`_s`$ inherits the properties of $`\stackrel{\text{L}}{}`$.
We denote the equivalence class of $`t`$ under $`_s`$ with $`[t]_s`$. We have $`[t]_s[[s]][t_0]`$.
Let $`X,𝒪^{t_0}`$ be the subset space where
$$X=\{tt[t_0]\}$$
and
$$𝒪^{t_0}=\{[t]_st[[s]]\text{and}sS^{t_0}\}.$$
It is clear that $`𝒪^{t_0}𝒫(X)`$.
###### Lemma 12
If $`[s_1][s_2]`$ and $`t[[s_1]][[s_2]]`$ then $`[t]_{s_1}[t]_{s_2}`$.
Proof. Immediate from the definition of $`_s`$ .
To elaborate the above process, we present the following simple example.
Example: A part of the canonical model appears in Figure 2.
(Horizontal and downward arrows correspond to $`\stackrel{\text{L}}{}`$ and $`\stackrel{\mathrm{}}{}`$, respectively.) We would like to make subsets of a treelike space correspond to equivalence classes under $`\stackrel{\text{L}}{}`$. Canonical model worlds related with $`\stackrel{\mathrm{}}{}`$ will be represented by a single point. However, this model is not a treelike model: $`\{r_1,t_1\}`$ and $`\{r_3,s_1,t_1\}`$ should make two distinct points. To remedy that, we “trace back” each equivalence class under $`\stackrel{\text{L}}{}`$ to the uppermost one. For instance, $`[t_1]=\{t_1,t_2\}`$ is traced back to $`[r_1]=\{r_1,r_2,r_3,r_4\}`$. The latter forms $`[[t_1]]`$. Next, we split $`[[t_1]]`$ into equivalence classes under $`_{t_1}`$, i.e. $`[r_1]_{t_1}=\{r_1,r_2\}`$ and $`[r_3]_{t_1}=\{r_3,r_4\}`$, since $`r_1_{t_1}r_2`$ and $`r_3_{t_1}r_4`$. Finally, we replace $`[t_1]`$ with as many copies as these equivalence classes (see Figure 3).
The infinite case is taken care of by Lemma 14. The resulting space (of Figure 3) is a treelike space. Note that we could have replaced this procedure by one that employs maximal branches but we find the present one simpler.
###### Proposition 13
The subset space $`X,𝒪^{t_0}`$ is a treelike space.
Proof. Suppose $`[t_1]_{s_1}[t_2]_{s_2}\mathrm{}`$. Let $`t[t_1]_{s_1}[t_2]_{s_2}`$. We have either $`s_1\stackrel{\mathrm{}}{}s_2`$ or $`s_2\stackrel{\mathrm{}}{}s_1`$, since $`t\stackrel{\mathrm{}}{}s_1^{}`$, $`t\stackrel{\mathrm{}}{}s_2^{}`$, for some $`s_1^{}[s_1]`$ and $`s_2^{}[s_2]`$, and the canonical frame is connected. The former implies $`[s_1][s_2]`$. Thus, by Lemma 12, $`[t_1]_{s_1}=[t]_{s_1}[t]_{s_2}=[t_2]_{s_2}`$. Similarly, the latter implies $`[t]_{s_2}[t]_{s_1}`$.
Let $`X,𝒪^{t_0},i`$ be the treelike model where $`X`$ and $`𝒪^{t_0}`$ are as above, and $`i(A)=v^{t_0}(A)`$ where $`v^{t_0}`$ is the initial interpretation restricted on $`𝒞^{t_0}`$.
An element of $`X\dot{\times }𝒪^{t_0}`$ can have more than one representation. In order to prove the semantical equivalence we are opting for, we shall choose a canonical representation. So, given a pair $`(t,[t^{}]_s^{})X\dot{\times }𝒪^{t_0}`$, its canonical representation is $`(t,[t]_s)`$ where $`s`$ is such that $`t\stackrel{\mathrm{}}{}s\stackrel{\text{L}}{}s^{}`$. Its existence is assured by the definition of $`[t^{}]_s^{}`$ and uniqueness by Proposition 9(d). From now on, we shall use the canonical representation wherever is possible.
###### Lemma 14
Let $`t[t_0]`$ and $`sS^{t_0}`$ such that $`t\stackrel{\mathrm{}}{}s`$. Then for all $`s^{}[s]`$ there exists $`t^{}[t_0]`$ such that $`t^{}\stackrel{\mathrm{}}{}s^{}`$ and $`t_st^{}`$, i.e. $`t^{}[t]_s`$.
Proof. Let
$$\{t_i\}_{iI}$$
be the linear order of all members of $`S^{t_0}`$ under $`\stackrel{\mathrm{}}{}`$ such that $`t\stackrel{\mathrm{}}{}t_i\stackrel{\mathrm{}}{}s`$.
Now, let
$$\begin{array}{cccc}T^{}\hfill & =\hfill & & \{\mathrm{}\psi \psi s^{}\}\hfill \\ & & \hfill & \{\chi 𝖪\chi t\}\hfill \\ & & \hfill & \{\mathrm{}\omega 𝖪\omega t_i,\text{for some}iI\}\hfill \\ & & \hfill & \{\mathrm{}\varphi \mathrm{}𝖪\varphi t\text{and}\mathrm{}\varphi s^{}\}.\hfill \end{array}$$
$`T^{}`$ is consistent. For if not, then there would be $`\psi ,\omega _1,\omega _2,\mathrm{},\omega _n,\chi ,\varphi `$ as above with $`i_1,i_2,\mathrm{},i_nI`$ and $`i_1i_2\mathrm{}i_n`$ such that
$$_{\mathrm{𝐌𝐏𝐓}}\mathrm{}\psi \underset{k=1}{\overset{n}{}}\mathrm{}\omega _k\chi \mathrm{}\neg \varphi .$$
Thus
$$_{\mathrm{𝐌𝐏𝐓}}𝖪(\mathrm{}\psi \underset{k=1}{\overset{n}{}}\mathrm{}\omega _k\chi \mathrm{}\neg \varphi ).$$
We shall prove that the negation of the above formula belongs to $`t`$ and reach a contradiction. Since $`\psi \mathrm{}\varphi s^{}`$, we have $`𝖫(\psi \mathrm{}\varphi )s`$. Hence
$$\mathrm{}𝖫(\psi \mathrm{}\varphi )t_{i_1}.$$
Observe that $`\mathrm{}𝖪\varphi t_{i_1}`$ so, by applying axiom 12, we have
$$𝖫(\mathrm{}\psi \mathrm{}\varphi )t_{i_1}.$$
Since $`𝖪\omega _1t_{i_1}`$, we have
$$𝖫\left(\mathrm{}\psi \omega _1\mathrm{}\varphi \right)t_{i_1}.$$
Also, $`\mathrm{}𝖪\varphi t_{i_2}`$ and
$$\mathrm{}𝖫\left(\mathrm{}\psi \omega _1\mathrm{}\varphi \right)t_{i_2}.$$
So, by axiom 12,
$$𝖫\left(\mathrm{}(\mathrm{}\psi \omega _1)\mathrm{}\varphi \right)t_{i_2}.$$
Since $`𝖪\omega _2t_{i_2}`$, we have
$$𝖫\left(\mathrm{}\psi \omega _2\mathrm{}\omega _1\mathrm{}\varphi \right)t_{i_2}.$$
Also, $`\mathrm{}𝖪\varphi t_{i_3}`$ and
$$\mathrm{}𝖫\left(\mathrm{}\psi \omega _2\mathrm{}\omega _1\mathrm{}\varphi \right)t_{i_3}.$$
So, by axiom 12,
$$𝖫\left(\mathrm{}(\mathrm{}\psi \omega _2\mathrm{}\omega _1)\mathrm{}\varphi \right)t_{i_3},$$
i.e.
$$𝖫\left(\mathrm{}\psi \mathrm{}\omega _2\mathrm{}\omega _1\mathrm{}\varphi \right)t_{i_3}.$$
Arguing this way and by repeated applications of axiom 12 we have
$$𝖫\left(\mathrm{}\psi \underset{k=1}{\overset{n}{}}\mathrm{}\omega _k\mathrm{}\varphi \right)t.$$
Since $`𝖪\chi t`$, we have
$$𝖫\left(\mathrm{}\psi \underset{k=1}{\overset{n}{}}\mathrm{}\omega _k\chi \mathrm{}\varphi \right)t$$
which is the negation of the formula that $`\mathrm{𝐌𝐏𝐓}`$ proves. Therefore $`T^{}`$ is consistent. Let $`t^{}`$ be a maximal extension of $`T^{}`$.
We shall show that $`t^{}`$ is the required theory of the lemma. We begin by showing that if $`t^{}\stackrel{\mathrm{}}{}r^{}\stackrel{\mathrm{}}{}s^{}`$ then $`r^{}\stackrel{\text{L}}{}t_i`$, for some $`iI`$, i.e. $`t[[r^{}]]`$. So suppose that $`t^{}\stackrel{\mathrm{}}{}r^{}\stackrel{\mathrm{}}{}s^{}`$. If $`r^{}=s^{}`$ we are done. If not, let
$$R=\{\psi \mathrm{}\psi t\}\{𝖫\chi \chi r^{}\}.$$
$`R`$ is consistent. For if not, then there would be $`\psi `$ and $`\chi `$ as above such that
$$_{\mathrm{𝐌𝐏𝐓}}\psi \neg 𝖫\chi .$$
Since $`r^{}\stackrel{\mathrm{}}{}s^{}`$ and $`r^{}s^{}`$, there exists $`\chi ^{}r^{}`$ such that $`\mathrm{}\neg \chi ^{}s^{}`$. Let $`\varphi =\chi \chi ^{}`$. Observe that $`\mathrm{}\neg (\chi \chi ^{})s^{}`$, i.e. $`\mathrm{}\neg \varphi s^{}`$, and $`\varphi r^{}`$. Further,
$$_{\mathrm{𝐌𝐏𝐓}}\psi \neg 𝖫(\chi \chi ^{}),$$
i.e.
$$_{\mathrm{𝐌𝐏𝐓}}\psi \neg 𝖫\varphi ,$$
and therefore,
$$_{\mathrm{𝐌𝐏𝐓}}\mathrm{}\psi \mathrm{}𝖪\neg \varphi .$$
Now, we have $`\mathrm{}𝖪\neg \varphi t`$ and $`\mathrm{}\neg \varphi s^{}`$, since $`\mathrm{}\psi t`$. By definition of $`T^{}`$ above, we have $`\mathrm{}\neg \varphi T^{}`$, and therefore $`\mathrm{}\neg \varphi t^{}`$ ($`t^{}`$ is an extension of $`T^{}`$). In this case, $`\neg \varphi r^{}`$ which is a contradiction. Therefore $`R`$ is consistent. So a maximal extension $`r`$ of $`R`$ has the property $`t\stackrel{\mathrm{}}{}r\stackrel{\mathrm{}}{}s`$. Hence $`r=t_i`$, for some $`iI`$.
We must now prove that $`t^{}[[t_i]]`$, for all $`iI`$. Let
$$T_i^{}=\{\psi \mathrm{}\psi t^{}\}\{\omega 𝖪\omega t_i\}.$$
$`T_i^{}`$ is consistent. If not, then
$$_{\mathrm{𝐌𝐏𝐓}}\psi \neg \omega ,$$
for some $`\varphi `$ and $`\omega `$ as above, which implies
$$_{\mathrm{𝐌𝐏𝐓}}\mathrm{}\psi \mathrm{}\neg \omega ,$$
i.e.
$$_{\mathrm{𝐌𝐏𝐓}}\mathrm{}\psi \neg \mathrm{}\omega .$$
So $`\neg \mathrm{}\omega t^{}`$, since $`\mathrm{}\psi t^{}`$. But, by definition, $`\mathrm{}\omega T^{}t^{}`$ which is a contradiction. Therefore a maximal extension $`t_i^{}`$ of $`T_i^{}`$ is such that $`t^{}\stackrel{\mathrm{}}{}t_i^{}\stackrel{\text{L}}{}t_i`$. Hence $`t^{}[[t_i]]`$.
Combining the above proofs we have $`t_st^{}`$.
We now have the following theorem.
###### Theorem 15
For all $`sS^{t_0}`$ and $`tX`$ such that $`t\stackrel{\mathrm{}}{}s`$,
$$\varphi s\text{iff}t,[t]_s\varphi .$$
Proof. By induction on the structure of $`\varphi `$. For an atomic formula $`A`$, we have that $`ti(A)`$ if and only if $`si(A)=v^{t_0}(A)`$, i.e. $`As`$, because of Axiom 2 and $`t\stackrel{\mathrm{}}{}s`$.
Negation and conjunction are straightforward.
Suppose $`\varphi =\mathrm{}\psi `$. Let $`\mathrm{}\psi s`$ and $`t,[t]_s\mathrm{}\neg \psi `$, for some $`s`$ and $`t`$ as in the theorem’s statement. This implies that there exists $`s^{}S^{t_0}`$ such that $`[t]_s^{}[t]_s`$, $`t\stackrel{\mathrm{}}{}s^{}`$ and $`t,[t]_s^{}\neg \psi `$. By induction hypothesis, $`\neg \psi s^{}`$. We have now that $`t\stackrel{\mathrm{}}{}s`$ and $`t\stackrel{\mathrm{}}{}s^{}`$ which, by connectivity, implies either $`s^{}\stackrel{\mathrm{}}{}s`$ or $`s\stackrel{\mathrm{}}{}s^{}`$. In the former case, we have $`[s][s^{}]`$ and hence, by Lemma 12, $`[t]_s[t]_s^{}`$. So $`[t]_s=[t]_s^{}`$. Therefore $`s=s^{}`$, by Proposition 9(d), which is a contradiction to our hypothesis ($`\varphi s`$). In the latter case, we have $`s\stackrel{\mathrm{}}{}s^{}`$ which again contradicts our hypothesis ($`\varphi s`$).
For the other direction, suppose that $`t,[t]_s\mathrm{}\psi `$ and $`\mathrm{}\neg \psi s`$ for some $`s`$ and $`t`$ as above. Then there exists $`s^{}S^{t_0}`$ such that $`s\stackrel{\mathrm{}}{}s^{}`$ and $`\neg \psi s^{}`$. Thus, $`t,[t]_s^{}\neg \psi `$ by induction hypothesis. Moreover $`[t]_s^{}[t]_s`$ by Lemma 12, which is a contradiction.
If $`\varphi =𝖪\psi `$, let $`𝖪\psi s`$ and suppose $`t,[t]_s𝖫\neg \psi `$, for some $`s`$ and $`t`$ as in the theorem’s statement, towards a contradiction. Then there exists $`t^{}[t]_s`$ such that $`t^{},[t]_s\neg \psi `$, i.e. $`t^{},[t^{}]_s^{}\neg \psi `$, for some $`s^{}S^{t_0}`$ such that $`t^{}\stackrel{\mathrm{}}{}s^{}`$ and $`s^{}\stackrel{\text{L}}{}s`$, which is a contradiction.
For the other direction, suppose that $`t,[t]_s𝖪\psi `$ and $`𝖫\neg \psi s`$, for some $`s`$ and $`t`$ as above. Then there exist $`s^{}S^{t_0}`$ such that $`s\stackrel{\text{L}}{}s^{}`$ and $`\neg \psi s^{}`$. By Lemma 14, there exists $`t^{}[t]_s`$ such that $`t^{}\stackrel{\mathrm{}}{}s^{}`$. Then we have $`t^{},[t]_s\neg \psi `$ by induction hypothesis. Therefore $`t,[t]_s\neg 𝖪\psi `$ which is a contradiction.
Combining now Proposition 13 and Theorem 15 we have the following
###### Corollary 16
The system $`\mathrm{𝐌𝐏𝐓}`$ is complete with respect to treelike spaces.
### 3.2 Decidability
For each treelike model and formula $`\varphi `$, we shall construct an equivalent finite subset space of bounded size with respect to the complexity of $`\varphi `$. This is a kind of “semantic” filtration, based on geometric properties of treelike models, using a technique first introduced in \[Geo94a\].
In the following we assume that $`X,𝒪`$ is a treelike space. Our aim is to find a partition of $`𝒪`$, where a given formula $`\varphi `$ “retains its truth value” for each point throughout a member of this partition. It turns out that there exists a finite partition of this kind.
First we need some definitions. (Note that the following hold, although we refer to a treelike space $`𝒪`$, for an arbitrary family of subsets of $`X`$.)
Definition 17 Given a finite family $`=\{U_1,\mathrm{},U_n\}𝒫(X)`$, i.e. of subsets of $`X`$, we define the remainder of (the principal ideal in $`(𝒪,)`$ generated by) $`U_k`$ by
$$\mathrm{𝖱𝖾𝗆}^{}U_k=U_k\underset{U_kU_i}{}U_i,$$
where $`U_k=\{V𝒪VU_k\}`$. Note that $`\mathrm{𝖱𝖾𝗆}^{}U_k𝒪`$ (but not necessarily $`U_k𝒪`$).
###### Proposition 18
In a finite family $`=\{U_1,\mathrm{},U_n\}𝒫(X)`$ closed under intersection, we have
$$\mathrm{𝖱𝖾𝗆}^{}U_i=U_i\underset{U_jU_i}{}U_j,$$
for $`i=1,\mathrm{},n`$.
Proof.
$$\begin{array}{ccc}\hfill \mathrm{𝖱𝖾𝗆}^{}U_i& =& U_i_{U_iU_h}U_h\hfill \\ & =& U_i_{U_iU_h}(U_hU_i)\hfill \\ & =& U_i_{U_jU_i}U_i.\hfill \end{array}$$
We denote $`_{U_i}U_i`$ with $``$.
###### Proposition 19
If $`=\{U_1,\mathrm{},U_n\}`$ is a finite family of subsets of $`X`$ closed under intersection then
1. $`\mathrm{𝖱𝖾𝗆}^{}U_i\mathrm{𝖱𝖾𝗆}^{}U_j=\mathrm{}`$, for $`ij`$,
2. $`_{i=1}^n\mathrm{𝖱𝖾𝗆}^{}U_i=`$, i.e. $`\{\mathrm{𝖱𝖾𝗆}^{}U_i\}_{i=1}^n`$ is a partition of $``$. From now on we shall call a finite family of subsets $``$ closed under intersection a finite partition (of $``$),
3. if $`V_1,V_2𝒪`$, $`V_1\mathrm{𝖱𝖾𝗆}^{}U_i`$ and $`V_1V_2U_i`$ then $`V_2\mathrm{𝖱𝖾𝗆}^{}U_i`$, i.e. $`\mathrm{𝖱𝖾𝗆}^{}U_i`$ is convex,
4. if $`\{V_j\}_{jJ}\mathrm{𝖱𝖾𝗆}^{}U_i`$ then $`_{jJ}U_jU_i`$.
Proof. Parts a, c and d are immediate from the definition.
For Part b, suppose that $`V`$ then $`V\mathrm{𝖱𝖾𝗆}^{}_{VU_i}U_i`$.
Every partition of a set induces an equivalence relation on this set. The members of the partition comprise the equivalence classes. We denote the equivalence relation induced by $``$ by $`_{}`$.
Definition 20 Given a set of subsets $`𝒢`$, we define the relation $`_𝒢^{}`$ on $`𝒪`$ with $`V_1_𝒢^{}V_2`$ if and only if $`V_1UV_2U`$ for all $`U𝒢`$.
We have the following
###### Proposition 21
The relation $`_𝒢^{}`$ is an equivalence.
###### Proposition 22
Given a finite partition $``$, we have $`_{}^{}=_{}`$ i.e. the remainders of $``$ are the equivalence classes of $`_{}^{}`$.
Proof. Suppose $`V_1_{}^{}V_2`$ then $`V_1`$ and $`V_2`$ belong to $`\mathrm{𝖱𝖾𝗆}^{}U`$ where
$$U=\{U^{}|V_1,V_2U,U^{}\}.$$
For the opposite direction, suppose $`V_1,V_2\mathrm{𝖱𝖾𝗆}^{}U`$ and there exists $`U^{}`$ such that $`V_1U^{}`$ while $`V_2U^{}`$. Then we have $`V_1U^{}U`$, $`U^{}U`$ and $`U^{}UU`$ i.e. $`V_1\mathrm{𝖱𝖾𝗆}^{}U`$.
###### Proposition 23
If $`𝒢`$ is a finite set of subsets of $`X`$ then $`\mathrm{𝖢𝗅}(𝒢)`$, its closure under intersection, is a finite partition for $`𝒢`$.
The last proposition enables us to give yet another characterization of remainders: every family of points in a complete lattice closed under arbitrary joins comprises a closure system, i.e. a set of fixed points of a closure operator of the lattice (cf. \[GHK<sup>+</sup>80\].) Here the lattice is the powerset of $`X`$. If we restrict ourselves to a finite number of fixed points then we just ask for a finite set of subsets closed under intersection i.e. Proposition 23. Thus a closure operator in the lattice of the powerset of $`X`$ induces an equivalence relation to any family of subsets of $`X`$. Two subsets are equivalent if they have the same closure, and the equivalence classes of this relation are just the remainders of the subsets which are fixed points of the closure operator.
We now introduce the notion of stability corresponding to what we mean by “a formula retains its truth value on a set of subsets”.
Definition 24 Let $`𝒢𝒪`$ then $`𝒢`$ is stable for $`\varphi `$, if for all $`x`$, either $`x,V\varphi `$ for all $`V𝒢`$, or $`x,V\neg \varphi `$ for all $`V𝒢`$.
###### Proposition 25
Let $`𝒢_1`$,$`𝒢_2𝒪`$ then
1. if $`𝒢_1𝒢_2`$ and $`𝒢_2`$ is stable for $`\varphi `$ then $`𝒢_1`$ is stable for $`\varphi `$, and
2. if $`𝒢_1`$ is stable for $`\varphi `$ and $`𝒢`$ is stable for $`\chi `$ then $`𝒢_1𝒢_2`$ is stable for $`\varphi \chi `$.
Proof. Part a is easy to see while Part b is a corollary of Part a.
Definition 26 A finite partition $`=\{U_1,\mathrm{},U_n\}`$ is called a stable partition for $`\varphi `$, if $`\mathrm{𝖱𝖾𝗆}^{}U_i`$ is stable for $`\varphi `$, for all $`U_i`$.
###### Proposition 27
If $`=\{U_1,\mathrm{},U_n\}`$ is a stable partition for $`\varphi `$, so is
$$^{}=\mathrm{𝖢𝗅}(\{U_0,U_1,\mathrm{},U_n\}),$$
where $`U_0`$.
Proof. Let $`V^{}`$, then there exists $`U_l`$ such that $`\mathrm{𝖱𝖾𝗆}^{^{}}V\mathrm{𝖱𝖾𝗆}^{}U_l`$ (e.g. $`U_l=\{U_i|U_i,VU_i\}`$), i.e. $`^{}`$ is a refinement of $``$. But $`\mathrm{𝖱𝖾𝗆}^{}U_l`$ is stable for $`\varphi `$ and so is $`\mathrm{𝖱𝖾𝗆}^{^{}}V`$ by Proposition 25(a).
The above proposition says that a finite stable partition for a treelike space $`𝒪`$ remains stable if we “refine” it.
The following is the main theorem of this section. It says that for each formula $`\varphi `$ we can find a stable partition for $`\varphi `$ which is essentially a refinement of the stable partition corresponding to the subformulae of $`\varphi `$.
###### Theorem 28 (Partition Theorem)
Let $`=X,𝒪,i`$ be a treelike model. Then there exists a family $`\{^\psi \}_\psi `$ of finite stable partitions such that if $`\varphi `$ is a subformula of $`\psi `$ then $`^\varphi ^\psi `$ and $`^\psi `$ is a finite stable partition for $`\psi `$.
Proof. By induction on the structure of the formula $`\psi `$. In each step we refine the partition of the induction hypothesis. For each $`U^\psi `$, let $`U^\psi =\{xU:x,U\psi \}`$. This set determines completely the satisfaction of $`\psi `$ on $`\mathrm{𝖱𝖾𝗆}^^\psi U`$ whenever $`^\psi `$ is stable.
* If $`\psi =A`$ is an atomic formula then $`^A=\{X\}=\{i()\}`$, since $`𝒪`$ is stable for all atomic formulae. We have $`X^A=i(A)`$.
* If $`\psi =\neg \varphi `$ then let $`^\psi =^\varphi `$, since the statement of the theorem is symmetric with respect to negation. We also have $`U^\psi =(XU^\varphi )U`$, for all $`U^\psi `$.
* If $`\psi =\chi \varphi `$, let
$$^\psi =\mathrm{𝖢𝗅}(^\chi ^\varphi ).$$
Observe that $`^\chi ^\varphi ^{\chi \varphi }`$. Now, $`^\psi `$ is a stable partition for $`\chi \varphi `$ containing $`X`$, since it is a refinement of both $`^\chi `$ and $`^\varphi `$. Thus, $`^\psi `$ is a finite stable partition for $`\psi `$ containing $`X`$.
* Suppose $`\psi =𝖪\varphi `$. Then, by induction hypothesis, there exists a finite stable partition $`^\varphi =\{U_1,\mathrm{},U_n\}`$ for $`\varphi `$ containing $`X`$.
Now, if $`V\mathrm{𝖱𝖾𝗆}^^\varphi U_iU_i^\varphi `$, for some $`i\{1.\mathrm{},n\}`$, then $`x,V\varphi `$, for all $`xV`$, by definition of $`U_i^\varphi `$. Hence $`x,V𝖪\varphi `$, for all $`xV`$.
On the other hand, if $`V\mathrm{𝖱𝖾𝗆}^^\varphi U_iU_i^\varphi `$ then there exists $`xV`$ such that $`x,V\neg \varphi `$ (otherwise $`VU_i^\varphi `$). Thus we have $`x,V\neg 𝖪\varphi `$, for all $`xV`$. Hence $`\mathrm{𝖱𝖾𝗆}^^\varphi U_iU_i^\varphi `$ and $`\mathrm{𝖱𝖾𝗆}^^\varphi U_iU_i^\varphi `$ are stable for $`𝖪\varphi `$. Thus the set
$$\begin{array}{ccc}F\hfill & =\hfill & \{\mathrm{𝖱𝖾𝗆}^{}U_i|U_i^\varphi \mathrm{𝖱𝖾𝗆}^{}U_i\}\hfill \\ & & \{\mathrm{𝖱𝖾𝗆}^{}U_jU_j^\varphi ,\mathrm{𝖱𝖾𝗆}^{}U_jU_j^\varphi |U_j^\varphi U_j\}\hfill \end{array}$$
is a partition of $`𝒪`$ and its members are stable for $`𝖪\varphi `$. Let
$$^{\text{K}\varphi }=\mathrm{𝖢𝗅}(^\varphi U_i^\varphi ).$$
We have that $`^{\text{K}\varphi }`$ is a finite set of opens and $`^\varphi ^{\text{K}\varphi }`$. Thus $`^{\text{K}\varphi }`$ is finite and contains $`X`$. We have only to prove that $`^{\text{K}\varphi }`$ is a stable partition for $`𝖪\varphi `$, i.e. every remainder of an open in $`^{\text{K}\varphi }`$ is stable for $`𝖪\varphi `$. But for that, observe that $`^{\text{K}\varphi }`$ is a refinement of $`F`$. Therefore $`^{\text{K}\varphi }`$ is a finite stable partition for $`𝖪\varphi `$, using Proposition 25(a).
Now, if $`U^\psi `$ then either $`U^{\text{K}\varphi }=U`$ or $`U^{\text{K}\varphi }=\mathrm{}`$.
* Suppose $`\psi =\mathrm{}\varphi `$. Then, let
$$^\mathrm{}\varphi =^\varphi ,$$
where $`^\varphi `$ is a finite stable partition for $`\varphi `$ by induction hypothesis.
We shall show that $`^\varphi `$ is also a finite stable spitting for $`\mathrm{}\varphi `$. Pick $`U^\varphi `$ and $`xU`$. If $`x,V\neg \varphi `$, for all $`VU`$ such that $`xV`$, we are done, since $`x,V\neg \mathrm{}\varphi `$. If $`x,V\varphi `$, for some $`V\mathrm{𝖱𝖾𝗆}^^\varphi U`$, then $`x,W\varphi `$, for all $`W\mathrm{𝖱𝖾𝗆}^^\varphi U`$, since $`^\varphi `$ is stable for $`\varphi `$. Therefore $`x,W\mathrm{}\varphi `$ for all $`W\mathrm{𝖱𝖾𝗆}^^\varphi U`$. If $`x,V\varphi `$, for some $`VU`$ with $`V\mathrm{𝖱𝖾𝗆}^^\varphi U`$, then we have $`VW`$, for all $`W\mathrm{𝖱𝖾𝗆}^^\varphi U`$, since the set of subsets containing $`x`$ is linearly ordered and $`\mathrm{𝖱𝖾𝗆}^^\varphi U`$ is stable and convex. Hence $`x,W\mathrm{}\varphi `$, for all $`W\mathrm{𝖱𝖾𝗆}^^\varphi U`$.
The following corollary is “folklore”.
###### Corollary 29
The formula $`\mathrm{}\mathrm{}\varphi \mathrm{}\mathrm{}\varphi `$ is sound in treelike spaces.
Proof. Let $`x,U\mathrm{}\mathrm{}\varphi `$ in some model $`X,𝒪,i`$.
By the Partition theorem, there exists a finite stable partition $``$ for $`\varphi `$. Further, there is a $`V`$ which is “the least” in the following sense: if $`W,W^{}𝒪`$ contain $`x`$, $`W\mathrm{𝖱𝖾𝗆}^{}V`$, and $`W^{}W`$ then we will also have $`W^{}\mathrm{𝖱𝖾𝗆}^{}V`$. The existence of such a set $`V`$ is assured by the fact that $``$ is finite, the members of the partition which $``$ induces are convex, and the set of subsets in $`𝒪`$ which contain $`x`$ is linearly ordered. Moreover, $`\mathrm{𝖱𝖾𝗆}^{}V`$ contains at least one subset which contains $`x`$, say $`W`$.
Now, we have either $`UV`$ or $`VU`$. In the former case, we have $`U\mathrm{𝖱𝖾𝗆}^{}V`$. Hence $`x,U\mathrm{}\varphi `$ as $`\mathrm{𝖱𝖾𝗆}^{}V`$ is stable for $`\varphi `$. In the latter case, $`x,W\mathrm{}\varphi `$, since $`WVU`$. Thus we have $`x,W\mathrm{}\varphi `$ for the same reasons as above. Hence $`x,U\mathrm{}\mathrm{}\varphi `$.
A finite partition does not have a treelike form. Therefore we cannot perform a filtration in a direct manner. First, we shall consider no partition member (remainder) that contains no subset belonging to the initial treelike space. Next, we shall impose a relation $``$ among the remaining members (Definition 30). Two remainders will be related just in case they contain subsets with common elements. This relation is not a partial order. However, it respects the initial treelike ordering (Lemma 32 through 35). Finally, using a number of equivalence relations based on $``$, one for each member of the partition, we shall construct a treelike model equivalent to the initial one (Propositions 38 and 39). Moreover, the underlying space of this model will contain a finite number of subsets.
By the Partition theorem, given a treelike model $`X,𝒪,i`$ and a formula $`\varphi `$, there exists a finite partition $`^\varphi `$ on $`𝒪`$ stable for $`\varphi `$. For each $`U^\varphi `$, let
$$\overline{U}=\mathrm{𝖱𝖾𝗆}^^\varphi U$$
and
$$\overline{^\varphi }=\{UU^\varphi \text{and}\overline{U}\mathrm{}\}.$$
We have the following
###### Lemma 30
If $`U_1,U_2\overline{^\varphi }`$ with $`\overline{U_1}\overline{U_2}`$, and $`V_1,V_2𝒪`$ with $`V_1\mathrm{𝖱𝖾𝗆}^^\varphi U_1`$, $`V_2\mathrm{𝖱𝖾𝗆}^^\varphi U_2`$ and $`V_1V_2\mathrm{}`$, then $`V_1V_2`$.
Proof. Since $`V_1V_2\mathrm{}`$, then, by connectedness, we have either $`V_1V_2`$ or $`V_2V_1`$. If $`V_1V_2`$ then $`V_1V_2`$ since they belong to distinct equivalence classes. If $`V_2V_1`$ then we have $`V_2V_1\overline{U_1}\overline{U_2}`$. Hence $`V_1\mathrm{𝖱𝖾𝗆}^^\varphi U_2`$, by Proposition 19(c).
Definition 31 Let $`<`$ be the following relation on $`\overline{^\varphi }`$
$$\begin{array}{ccc}U_1<U_2\hfill & \text{iff}\hfill & \overline{U_1}\overline{U_2}\mathrm{},\text{and}\hfill \\ & & \text{for all}x,V_1,V_2\text{such that}x\overline{U_1}\overline{U_2},V_1\mathrm{𝖱𝖾𝗆}^^\varphi U_1\text{with}xV_1,\hfill \\ & & \text{and}V_2\mathrm{𝖱𝖾𝗆}^^\varphi U_2\text{with}xV_2,V_1V_2.\hfill \end{array}$$
Clearly, we cannot have $`U_1<U_2`$ and $`U_2<U_1`$. Let $`U_1U_2`$, if either $`U_1=U_2`$ or $`U_1<U_2`$.
The following lemma allows us to weaken the conditions of the definition of $`<`$.
###### Lemma 32
Let $`U_1,U_2\overline{^\varphi }`$ with $`U_1U_2`$. If there exist $`x\overline{U_1}\overline{U_2}`$ and $`V_1\mathrm{𝖱𝖾𝗆}^^\varphi U_1`$, $`V_2\mathrm{𝖱𝖾𝗆}^^\varphi U_2`$ with $`xV_1V_2`$ such that $`V_1V_2`$, then $`U_1<U_2`$.
Proof. Suppose, towards a contradiction, that for $`y\overline{U_1}\overline{U_2}`$ there exist $`W_1\mathrm{𝖱𝖾𝗆}^^\varphi U_1`$ and $`W_2\mathrm{𝖱𝖾𝗆}^^\varphi U_2`$ such that $`yW_1W_2`$ and $`W_2W_1`$. By our hypothesis, $`U_1U_2`$ and $`V_1V_2`$, and so we have $`U_2U_1`$. This implies that $`\mathrm{𝖱𝖾𝗆}^^\varphi U_2U_1=\mathrm{}`$. Therefore $`W_2\mathrm{𝖱𝖾𝗆}^^\varphi U_2`$ which is a contradiction. Thus $`W_1W_2`$. Hence $`U_1<U_2`$.
###### Lemma 33
Let $`U_1,U_2\overline{^\varphi }`$. If $`\overline{U_1}\overline{U_2}\mathrm{}`$ then $`U_1U_2`$ or $`U_2U_1`$.
Proof. Suppose that $`U_1U_2`$ and let $`x\overline{U_1}\overline{U_2}`$. Let $`V_1\mathrm{𝖱𝖾𝗆}^^\varphi U_1`$ and $`V_2\mathrm{𝖱𝖾𝗆}^^\varphi U_2`$ such that $`xV_1V_2`$. Since $`𝒪`$ is a treelike space, we have either $`V_1V_2`$ or $`V_2V_1`$. Suppose that the former holds. Since $`U_1U_2`$, we have $`V_1V_2`$. By Lemma 32, $`U_1<U_2`$. Similarly, if $`V_2V_1`$ then $`U_2<U_1`$.
###### Lemma 34
$``$ is reflexive and antisymmetric.
Proof. Reflexivity is straightforward. For antisymmetry, suppose that $`\overline{U_1}\overline{U_2}\mathrm{}`$, $`U_1U_2`$ and $`U_2U_1`$. If $`U_1U_2`$ then we have $`U_1<U_2`$ and $`U_2<U_1`$ which is a contradiction.
Instead of transitivity, we have the following property of $``$:
###### Lemma 35
Let $`U_1,U_2,U_3\overline{^\varphi }`$. If $`U_1U_2`$, $`U_2U_3`$ and $`\overline{U_1}\overline{U_2}\overline{U_3}\mathrm{}`$ then $`U_1U_3`$.
Proof. If either $`U_1=U_2`$ or $`U_2=U_3`$ we are done, so suppose that $`U_1<U_2`$ and $`U_2<U_3`$. Let $`x\overline{U_1}\overline{U_2}\overline{U_3}`$, $`V_1\mathrm{𝖱𝖾𝗆}^^\varphi U_1`$ and $`V_3\mathrm{𝖱𝖾𝗆}^^\varphi U_3`$ such that $`xV_1`$ and $`xV_3`$. Since $`x\overline{U_2}`$, there exists $`V_2\mathrm{𝖱𝖾𝗆}^^\varphi U_2`$ such that $`xV_2`$. Also, we have $`V_1V_2V_3`$, since $`U_1<U_2`$ and $`U_2<U_3`$. So, by Lemma 32, $`U_1<U_3`$.
Since $`\overline{^\varphi }^\varphi `$, $`\overline{^\varphi }`$ is finite. Let $`\overline{^\varphi }=\{U_1,U_2,\mathrm{},U_n\}`$, for some $`n`$. Now, let $`_i`$ be the following equivalence relation on $`\overline{U_i}`$
$$\begin{array}{ccc}x_iy\hfill & \text{iff}\hfill & \text{for all}\overline{U_j},j\{1,2,\mathrm{},n\}\text{such that}U_iU_j,\hfill \\ & & x\overline{U_j}\text{iff}y\overline{U_j}.\hfill \end{array}$$
We denote the equivalence of $`x`$ under $`_i`$ with $`[x]_i`$. Observe that the number of equivalence classes is finite, since it depends only on the number of members of the partition.
###### Lemma 36
Let $`U_k,U_l\overline{^\varphi }`$, $`k,l\{1,2,\mathrm{},n\}`$, with $`\overline{U_k}\overline{U_l}\mathrm{}`$. Then
1. if $`U_kU_l`$ then $`[x]_k[x]_l`$ , for all $`x\overline{U_k}\overline{U_l}`$, and
2. if $`[x]_k[x]_l`$, for some $`x\overline{U_k}\overline{U_l}`$, then $`U_k<U_l`$.
Proof. For Part a, if $`U_k=U_l`$ then we are done. Suppose $`U_k<U_l`$ and let $`z[x]_k`$. Let $`U_m\overline{^\varphi }`$, $`m\{1,2,\mathrm{},n\}`$, be such that $`U_lU_m`$. If $`x\overline{U_m}`$ then $`x\overline{U_k}\overline{U_l}\overline{U_m}`$. So, by Lemma 35, $`U_kU_m`$. So $`z\overline{U_m}`$, since $`x_iz`$. For the other direction, suppose $`z\overline{U_m}`$. Then we have $`z\overline{U_l}`$, since $`U_kU_l`$ and $`x_1z`$. So $`z\overline{U_k}\overline{U_l}\overline{U_m}`$. Hence, by Lemma 35, $`U_kU_m`$. Also, $`x\overline{U_m}`$, since $`x_iz`$. Therefore $`z[x]_l`$.
For Part b, we have either $`U_k<U_l`$ or $`U_lU_k`$, since $`\overline{U_k}\overline{U_l}\mathrm{}`$. Suppose the latter towards a contradiction. Then, by Part a and Lemma 33, we have $`[x]_l[x]_k`$ which is a contradiction to our hypothesis.
###### Lemma 37
Let $`U_k,U_l\overline{^\varphi }`$, $`k,l\{1,2,\mathrm{},n\}`$, with $`\overline{U_k}\overline{U_l}\mathrm{}`$. If $`U_k<U_l`$ then $`[x]_k[x]_l`$, for all $`x\overline{U_k}\overline{U_l}`$.
Proof. By Lemma 36(a), we have $`[x]_k[x]_l`$. Suppose $`[x]_k=[x]_l`$. Let $`V\mathrm{𝖱𝖾𝗆}^^\varphi U_i`$ such that $`xV`$. We have $`V[x]_i`$. Thus $`V[x]_j\overline{U_j}`$. So, for each $`yV`$, there exists $`V_y\mathrm{𝖱𝖾𝗆}^^\varphi U_j`$ such that $`V_yV`$. But then $`V=_{yV}V_y\mathrm{𝖱𝖾𝗆}^^\varphi U_j`$ which is a contradiction, since $`U_iU_j`$.
Now, let
$$[\overline{^\varphi }]=\{[x]_ix\overline{U_i},i\{1,2,\mathrm{},n\}\}.$$
###### Proposition 38
The subset space $`X,[\overline{^\varphi }]`$ is a treelike space.
Proof. First notice that $`X\overline{^\varphi }`$, since $`X^\varphi `$. Thus $`X=U_{i_0}`$, for some $`i_0\{1,2,\mathrm{},n\}`$. Moreover, $`x_{i_0}y`$, for all $`x,yX`$. Hence $`X=[x]_{i_0}[\overline{^\varphi }]`$.
Now, let $`[x]_i[y]_j\mathrm{}`$, for some $`x\overline{U_i}`$ and $`y\overline{U_j}`$. Let $`z[x]_i[y]_j`$. We have $`[x]_i=[z]_i`$ and $`[y]_j=[z]_j`$. Further, $`z\overline{U_i}`$ and $`z\overline{U_j}`$, i.e. $`\overline{U_i}\overline{U_j}\mathrm{}`$. So, by Lemma 33, we have either $`U_iU_j`$ or $`U_jU_i`$. By Lemma 36(a), we have either $`[z]_i[z]_j`$ or $`[z]_j[z]_i`$, respectively. Therefore either $`[x]_i[y]_j`$ or $`[y]_i[x]_j`$.
Let $`\overline{}=X,[\overline{^\varphi }],\overline{i}`$ be the treelike model where $`\overline{i}(A)=\{[x]_ixi(A)\}`$.
###### Proposition 39
For all $`xX`$, $`V𝒪`$ and $`\psi `$ such that $`\psi `$ is a subformula of $`\varphi `$, if $`V\mathrm{𝖱𝖾𝗆}^^\varphi U_i`$, for some $`i\{1,2,\mathrm{},n\}`$, then
$$x,V_{}\psi \text{iff}x,[x]_i_\overline{}\psi .$$
Proof. By induction on the complexity of $`\varphi `$. The only interesting case is that of $`\varphi =\mathrm{}\psi `$. Suppose $`x,[x]_i\mathrm{}\neg \psi `$ but $`x,V_{}\mathrm{}\psi `$, for some $`V\mathrm{𝖱𝖾𝗆}^^\psi U_i`$. The latter implies that there is $`j\{1,2,\mathrm{},n\}`$ such that $`x,[x]_j\neg \psi `$ and $`[x]_j[x]_i`$. We have $`x,[x]_i\psi `$, by $`x,V_{}\psi `$ and induction hypothesis. Hence $`[x]_j[x]_i`$. By Lemma 36(b), we have $`U_j<U_i`$. By induction hypothesis, we have $`x,V^{}_{}\neg \psi `$, for all $`V^{}\mathrm{𝖱𝖾𝗆}^^\psi U_j`$ such that $`xV^{}`$. Also, we have $`V^{}V`$, since $`\overline{U_j}<\overline{U_i}`$. Hence $`x,V_{}\neg \psi `$, a contradiction.
Now, suppose $`x,[x]_i\mathrm{}\psi `$ but $`x,V_{}\mathrm{}\neg \psi `$, for some $`V\mathrm{𝖱𝖾𝗆}^^\psi U_i`$. So there exists $`V^{}\mathrm{𝖱𝖾𝗆}^^\psi U_j`$ such that $`xV^{}`$, $`V^{}V`$, and $`x,V^{}_{}\neg \psi `$. Also, $`x,[x]_i\psi `$ so, by induction hypothesis, $`x,V_{}\psi `$. The latter implies $`U_iU_j`$, since $`\mathrm{𝖱𝖾𝗆}^^\psi U_i`$ is stable for $`\psi `$. Therefore we have $`U_j<U_i`$. Hence, by Lemma 37, $`[x]_j[x]_i`$. Thus $`x,[x]_j\neg \psi `$, by induction hypothesis. Hence $`x,[x]_i\mathrm{}\neg \psi `$, a contradiction to our hypothesis.
Constructing the above model is not adequate for generating a finite model, since there may still be an infinite number of points. It turns out that we only need a finite number of them.
Let $`=X,𝒪,i`$ be a treelike model, and define an equivalence relation $``$ on $`X`$ by $`xy`$ iff
1. for all $`U𝒪`$, $`xU`$ iff $`yU`$, and
2. for all atomic $`A`$, $`xi(A)`$ iff $`yi(A)`$.
Further, denote by $`x^{}`$ the equivalence class of $`x`$, and let $`X^{}=\{x^{}:xX\}`$. For every $`U𝒪`$, let $`U^{}=\{x^{}:xU\}`$, then $`𝒪^{}=\{U^{}:U𝒪\}`$ is a treelike space on $`X^{}`$. Define a map $`i^{}`$ from the atomic formulae to the powerset of $`X^{}`$ by $`i^{}(A)=\{x^{}:xi(A)\}`$. The entire model $``$ lifts to the model $`^{}=X^{},𝒪^{},i^{}`$ in a well-defined way.
###### Lemma 40
For all $`x`$, $`U`$, and $`\varphi `$,
$$x,U_{}\varphi \text{iff}x^{},U^{}_{^{}}\varphi .$$
Proof. By induction on $`\varphi `$.
###### Theorem 41
If $`\varphi `$ is satisfied in any treelike space then $`\varphi `$ is satisfied in a finite treelike space.
Proof. Let $`=X,𝒪,i`$ be such that, for some $`xU𝒪`$, $`x,U_{}\varphi `$. Let $`^\varphi `$ be a finite stable partition (by Theorem 28) for $`\varphi `$ and its subformulae with respect to $``$. By Proposition 39, $`x,U_𝒩\varphi `$, where $`𝒩=X,,i`$. We may assume that $``$ is a treelike space, and we may also assume that the overall language has only the (finitely many) atomic symbols which occur in $`\varphi `$. Then the relation $``$ has only finitely many classes. So the model $`𝒩^{}`$ is finite. Finally, by Lemma 40, $`x^{},U^{}_𝒩^{}\varphi `$.
Observe that the finite treelike space is a quotient of the initial one under two equivalences. The one equivalence is on the elements of the treelike space and the number of equivalence classes is a function of the complexity of $`\varphi `$. The other equivalence is on the points of the treelike space and the number of equivalence classes is a function of the atomic formulae appearing in $`\varphi `$. So the overall size of the (finite) treelike space is bounded by a function of the complexity of $`\varphi `$. Thus if we want to test if a given formula is invalid we have a finite number of finite treelike spaces where we have to test its validity. Thus we have the following
###### Theorem 42
The theory of treelike spaces is decidable.
Acknowledgments: The author is indebted to Rohit Parikh for bringing this problem to his attention and wishes to thank Bernhard Heinemann, Larry Moss, and Timothy Williamson, as well as, the anonymous referees for helpful comments.
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# Contents
## 1 Introduction
The idea that spacetime has more than four dimensions is actually quite old. Already in the 1920’s, Kaluza suggested that gravity and electromagnetism can be interpreted as the degrees of freedom of the metric of a five-dimensional spacetime . Later, Klein gave an explanation for the fact that the extra dimension is not observed by suggesting that the extra dimension is compact and very small. Since then the idea has been studied from many different perspectives e.g. in Kaluza-Klein supergravity theories and also in string theories – where more than three spatial dimensions naturally arise but the extra dimensions are usually assumed to be of Planck size for not been directly observable. In another direction it has been suggested that spacetime can have more than three noncompact spatial dimensions if we live on a four-dimensional domain wall which is embedded in the higher dimensions. More recently, there has been a renewed interest in the topic since progresses in string theories have modified the old scenario (where the extra dimensions cannot exceed the tiny scale $`1`$ TeV$`{}_{}{}^{1}10^{19}`$ m) suggesting that Standard Model gauge interactions could be confined to a four-dimensional subspace – or brane–world – whereas gravity can still propagate in the whole bulk spacetime. Actually, the possibility that part of the standard model particles live in large (TeV) extra dimensions was first put forward in connection to the problem of supersymmetry breaking in string theory . These scenarios presents us with the enticing possibility to explain some long-standing particle physics problems by geometrical means .
In the canonical example of , spacetime is a direct product of ordinary four-dimensional spacetime and a (flat) spatial $`d`$-torus of common linear size $`r_c`$. Within this simple model, the large hierarchy between the weak scale and the fundamental scale of gravity can be eliminated. However, the hierarchy only arises in the presence of a large volume for the compactified dimension which is very difficult to justify. A more compelling scenario was introduced by Randall and Sundrum (herein RS). Reviving an old idea , RS proposed a set–up with the shape of a gravitational condenser in which two branes of opposite tension (which gravitationally repel each other) are stabilized by a slab of anti-de Sitter (AdS) space . In this model the extra dimension is strongly curved, and the distance scales on the brane with negative tension are exponentially smaller than those on the positive tension brane. Such exponential suppression can then naturally explain why the observed physical scales are so much smaller than the Plank scale. In further work RS found that gravitons can be localized on a brane which separates two patches of $`AdS_5`$ spacetime , suggesting that it is possible to have an infinite extra dimension . The question whether this scenario reproduces the usual four-dimensional gravity beyond the Newton’s law has been analyzed and cosmological considerations of models with large extra dimensions confirms that they are at least consistent candidates for describing our world . These ideas have raised a lot of interest in the subject and several groups have begun to work on possible experimental signatures of the extra dimension(s) .
In this paper we shall discuss the creation of brane–worlds in $`AdS`$ bulk. The approximation scheme to be used is the minisuperspace restriction of the canonical Wheeler–DeWitt formalism. The basic idea of this approach, commonly adopted in quantum cosmology calculations , is to separate the space-like metric into “modes”, and then insist that all the “translational” modes are “frozen out” by using the classical field equations, leaving only the scale factor to be quantized. The outline of the paper is as follows. We begin in section 2 by deriving a brane-big-bang in $`AdS_3`$. This lower–dimensional model provides a simple setting in which certain basic physical phenomena can be easily demonstrated while avoiding the mathematical complexities associated with the higher–dimensional counterparts. In section 3 we consider multi-dimensional brane-worlds, discussing the possible cosmologies within the WKB approximation. In section 4 we analyze the implications of the $`AdS/CFT`$ correspondence to quantum cosmology.
## 2 Brane–world in $`AdS_3`$
### 2.1 Wheeler–DeWitt Equation
In this section we consider the creation of a onebrane in $`AdS_3`$ within the framework of quantum cosmology . Thus the universe will initially be described by three–dimensional Anti-de Sitter space in which onebrane bubbles can nucleate spontaneously. As we shall see below, these bubbles appear (classically) at a critical size and then expand.
We thus begin by considering the action for a onebrane coupled to gravity,
$$S_{\mathrm{tot}}=\frac{L_p}{16\pi }_\mathrm{\Omega }d^3x\sqrt{g}\left(R+\frac{2}{\mathrm{}^2}\right)+\frac{L_p}{8\pi }_\mathrm{\Omega }d^2x\sqrt{\gamma }𝔎+T_\mathrm{\Omega }d^2x\sqrt{\gamma },$$
(1)
where $`𝔎`$ stands for the trace of the extrinsic curvature of the boundary, $`\gamma `$ is the induced metric on the brane, and $`T`$ is the brane tension.<sup>4</sup><sup>4</sup>4In our convention, the extrinsic curvature is defined as $`𝔎_{\mu \nu }=1/2(_\mu \widehat{n}_\nu +_\nu \widehat{n}_\mu )`$, where $`\widehat{n}^\nu `$ is the outward pointing normal vector to the boundary $`\mathrm{\Omega }`$. Lower Greek subscripts run from 0 to 2, capital Greek subscripts from 0 to $`d`$, and capital Latin subscripts from 0 to $`(d1)`$. Throughout the paper we adopt geometrodynamic units so that $`G1`$, $`c1`$ and $`\mathrm{}L_p^2M_p^2`$, where $`L_p`$ and $`M_p`$ are the Planck length and Planck mass, respectively. The first term is the usual Einstein-Hilbert (EH) action with a negative cosmological constant ($`\mathrm{\Lambda }=1/\mathrm{}^2`$). The second term is the Gibbons-Hawking (GH) boundary term, necessary for a well defined variational problem . The third term corresponds to a constant “vacuum energy”, i.e. a cosmological term on the boundary.
We wish to consider a brane which bounds two regions of $`AdS_3`$. If we further specialize to the case of spherical symmetry where
$$ds_3^2=\left(1+\frac{y^2}{\mathrm{}^2}\right)dt^2+\left(1+\frac{y^2}{\mathrm{}^2}\right)^1dy^2+y^2d\varphi ^2,$$
(2)
the geometry is uniquely specified by a single degree of freedom, the “radius” of the brane $`A(\tau )`$. The $`\tau `$ coordinate denotes proper time as measured along the brane-world. The computation of the GH boundary term has now reduced to that of computing the two non-trivial components of the second fundamental form. From Eq. (2) we find (see Appendix for details):
$$𝔎_\varphi ^\varphi =\frac{1}{A}\left[1+\frac{A^2}{\mathrm{}^2}+\dot{A}^2\right]^{1/2},$$
(3)
and
$$𝔎_\tau ^\tau =\left[\ddot{A}+\frac{A}{\mathrm{}^2}\right]\left[1+\frac{A^2}{\mathrm{}^2}+\dot{A}^2\right]^{1/2},$$
(4)
(where the dot denote a derivative with respect to proper time). After integration by parts the gravitational Lagrangian restricted to this minisuperspace may be identified as,
$$=\frac{L_p}{2}\left\{\dot{A}\mathrm{arcsinh}\left[\frac{\dot{A}}{\sqrt{1+A^2/\mathrm{}^2}}\right]+\sqrt{1+\frac{A^2}{\mathrm{}^2}+\dot{A}^2}\right\}.$$
(5)
The classical Wheeler–DeWitt Hamiltonian is now easily extracted. In order to do this we compute the conjugate momentum to $`A`$,
$$p=\frac{}{\dot{A}}=\frac{L_p}{2}\mathrm{arcsinh}\left[\frac{\dot{A}}{\sqrt{1+A^2/\mathrm{}^2}}\right].$$
(6)
This relation may be inverted to yield, $`\dot{A}=(1+A^2/\mathrm{}^2)^{1/2}\mathrm{sinh}(2p/L_p)`$, so that the Wheeler–DeWitt Hamiltonian is
$$_{\mathrm{tot}}p\dot{A}_{\mathrm{tot}}=2\pi AT\frac{L_p}{2}\sqrt{1+\frac{A^2}{\mathrm{}^2}}\mathrm{cosh}(2p/L_p).$$
(7)
Eq. (7) can be rewritten as,
$$_{\mathrm{tot}}=2\pi AT\frac{L_p}{2}\sqrt{1+\frac{A^2}{\mathrm{}^2}+\dot{A}^2}.$$
(8)
The Hamiltonian constraint – which follows from the requirement of diffeomorphism invariance – is $`_{\mathrm{tot}}=0`$, or equivalently,
$$\dot{A}^2=1A^2\left(\frac{1}{\mathrm{}^2}16\pi ^2\frac{T^2}{L_p^2}\right).$$
(9)
Observe that the constraint equation is consistent with the covariant conservation of the stress-energy tensor and reproduces the classical Einstein field equations of motion. It is easy to see from Eq. (9) that in order to obtain a real solution we need $`T0`$. Furthermore, the brane–world is (classically) bounded by a minimum radius
$$A_0^2=\left(\frac{1}{\mathrm{}^2}+\frac{16\pi ^2T^2}{L_p^2}\right)^1,$$
(10)
with $`16\pi ^2\mathrm{}^2T^2/L_p^2>1`$. In other words, the brane bubbles appear classically at a critical size and then their expansion is governed by (9). Note that as the world approaches the minimum size the expansion tends to zero. Once the world is dynamically stable it experiences an everlasting expansion. However, we shall soon see that quantum effects permit well–behaved wave functions for vanishing $`T`$. With the classical dynamics of the model understood and the Wheeler–DeWitt Hamiltonian at hand, quantization is straightforward. Canonical quantization proceeds via the usual replacement $`pi\mathrm{}/A`$. Naturally, the resulting quantum Hamiltonian has a factor order ambiguity. This factor-ordering ambiguity may be removed in a natural (though not unique) way by demanding that the quantum Hamiltonian be Hermitian,
$$\widehat{}_{\mathrm{tot}}=\frac{L_p}{2}\left(1+\frac{A^2}{\mathrm{}^2}\right)^{1/4}\mathrm{cos}\left[2L_p\frac{}{A}\right]\left(1+\frac{A^2}{\mathrm{}^2}\right)^{1/4}+2\pi AT.$$
(11)
That this Hamiltonian is Hermitian may formally be seen by Taylor-series expansion of the cosine. A more precise statement is that this Hamiltonian acts on the Hilbert space of square-integrable functions defined on the half–interval $`[0,\mathrm{})`$ subject to the constraint $`\psi (0)=0`$. This is most easily seen by noting that the Hamiltonian is Hermitian on $`L^2([0,\mathrm{}))`$ only if $`\psi (0)=0`$.<sup>5</sup><sup>5</sup>5Here, the Hilbert space scalar-product is given by the sum over all possible configurations (i.e. sizes) of the brane-world. Henceforth, for any operator $`\widehat{𝒪}`$ acting over any brane-wave-function $`\psi _j`$, $`<\psi _k|\widehat{𝒪}|\psi _j>=_0^{\mathrm{}}\psi _k^{}\widehat{𝒪}\psi _j𝑑A`$. In other words, $`A`$ is a parameter which governs the evolution of the brane along the null geodesic congruence of the extra dimension.
The wave function of the brane-world is determined in the usual fashion by the Wheeler–DeWitt equation $`\widehat{}_{\mathrm{tot}}\psi (A)=0`$. For the special case of $`T=0`$ we find the following solution:
$$\psi _{mn}(A)=C_{mn}(\phi _m\phi _n),$$
(12)
with
$$\phi _j=\left(1+\frac{A^2}{\mathrm{}^2}\right)^{1/4}\mathrm{exp}\left[\left(j+\frac{1}{2}\right)\frac{\pi }{2}\frac{A}{L_p}\right].$$
(13)
(See Fig.1. for a plot of some of these wavefunctions). Here $`m`$ and $`n`$ are integer valued quantum numbers describing the internal state of the brane. Negatives values of $`m`$, $`n`$ are not normalizable and so need to be discarded, as is the case when $`m=n`$.
Note that the appropriate normalization is $`|\psi |^2𝑑A=1`$, and that $`\psi (0)=0`$, as required. In fact the two terms in $`\psi _{mn}`$ individually satisfy the differential equation $`\widehat{}_{\mathrm{gravity}}\psi =0`$, but do not individually satisfy the boundary condition. By appropriate choice of $`C_{mn}`$ these states may be normalized, though they are not orthogonal to one another. The normalization constant takes the rather complicated form:
$`C_{mn}=\{{\displaystyle \frac{\mathrm{}\pi }{2}}[H_0(\mathrm{}(m+1/2)\pi /L_p)+H_0(\mathrm{}(n+1/2)\pi /L_p)2H_0(\mathrm{}(m+n+1)\pi /2L_p)`$
$`N_0(\mathrm{}(m+1/2)\pi /L_p)N_0(\mathrm{}(n+1/2)\pi /L_p)+2N_0(\mathrm{}(m+n+1)\pi /2L_p)]\}^{1/2},`$ (14)
where $`H_0(z)`$ is the Struve function and $`N_0(z)`$ is Neumann’s function<sup>6</sup><sup>6</sup>6These functions have the following respective integral representations: $`H_0(z)=\frac{2}{\pi }_0^1\frac{\mathrm{sin}(zt)dt}{\sqrt{1t^2}}`$ and $`N_0(z)=\frac{2}{\pi }_1^{\mathrm{}}\frac{\mathrm{cos}(zt)dt}{\sqrt{t^21}}`$.. With the wavefunctions at hand, one can calculate the mean value of the “radius” of the brane, i.e.
$$A=\frac{A|\psi _{mn}|^2𝑑A}{|\psi _{mn}|^2𝑑A}.$$
(15)
The integral in the numerator can be evaluated exactly but involves Meijer’s $`G`$-function $`G_{13}^{31}`$ and so the expression for $`A`$ is not of much practical use (anyhow, by dimensional analysis one would expect that this number would be of order $`L_p`$). Using numerical integration we have found $`A_{0,1}=0.54`$, $`A_{1,2}=0.01`$ and $`A_{2,3}=10^3`$ in units of $`L_p`$ and for $`\mathrm{}=1`$.
### 2.2 Qualitative Behaviour of Wavefunctions
In this subsection we take a first look at the problem of finding solutions to the Wheeler–DeWitt equation with $`T0`$ (in the next subsection we discuss the solutions in the WKB approximation). The relevant equation is:
$$\widehat{}_{\mathrm{tot}}\psi (A)=0$$
(16)
or, more explicitly
$$\frac{L_p}{2}\left(1+\frac{A^2}{\mathrm{}^2}\right)^{\frac{1}{4}}\mathrm{cos}\left[2L_p\frac{}{A}\right]\left(1+\frac{A^2}{\mathrm{}^2}\right)^{1/4}\psi +2\pi TA\psi =0.$$
(17)
After defining $`\phi (1+A^2/\mathrm{}^2)^{1/4}\psi `$, and the operator
$$\mathrm{\Delta }\underset{n=0}{\overset{\mathrm{}}{}}(1)^n\frac{(2L_p)^{2n}}{2n!}\frac{^{2n}}{A^{2n}},$$
(18)
we see that we have to solve
$$\mathrm{\Delta }\phi +\frac{4\pi TA}{L_p\sqrt{1+A^2/\mathrm{}^2}}\phi =0.$$
(19)
In order to gain some intuition for the behaviour of the solutions of this equation we will look for solutions in the two limits $`A/\mathrm{}1`$ and $`A/\mathrm{}1`$. In the case $`A/\mathrm{}1`$, using a trial solution of the form $`\phi =e^{\lambda A}`$, we find that Eq. (19) reduces to the following condition:
$$\mathrm{cos}(2L_p\lambda )=4\pi \frac{T\mathrm{}}{L_p}.$$
(20)
This is essentially the same condition as found before. Indeed, if $`16\pi ^2\mathrm{}^2T^2/L_p^2>1`$ then $`\lambda `$ will have pure imaginary values, leading to an oscillatory solution at infinity, that is not acceptable since it is not normalizable (it would in any case imply a delta function normalization, that we are not considering here) and should be something like a “classical” solution. On the other hand, if $`16\pi ^2\mathrm{}^2T^2/L_p^21`$, $`\lambda `$ will have two real solutions $`\lambda _\pm `$ of which one has to choose the negative one, with the same criteria of normalizability as before.
In order to analyze more carefully the behaviour of the wave function for $`A/\mathrm{}1`$, let us consider performing the following change of variables
$$A2L_pa.$$
(21)
With this change of variables and after scaling $`\psi \phi `$ as above, the Wheeler–DeWitt equation reads:
$$\stackrel{~}{\mathrm{\Delta }}\phi +\frac{8\pi Ta}{\sqrt{1+4L_p^2a^2/\mathrm{}^2}}\phi =0,$$
(22)
where
$$\stackrel{~}{\mathrm{\Delta }}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(1)^n}{(2n)!}\frac{^{2n}}{a^{2n}}.$$
(23)
It makes sense, since the factor $`L_p/\mathrm{}`$ is small, to analyze the behaviour of this equation for small values of the variable $`a`$ and expand the square root in series.
The plot in Fig. 2 shows the behaviour, in the interval $`[0,1]`$, of the square of the wave function $`|\psi |^2`$ that solves numerically Eq.(22), where we have considered eighteen orders of derivatives in $`\stackrel{~}{\mathrm{\Delta }}`$. It is observed that $`|\psi |^2`$ has a maximum at $`AL_p`$ as expected. It is worth to point out that the solutions with less derivatives have similar behaviour in the interval considered. It should be interesting to find a method to analyze the complete series.
Let us now consider a qualitative analysis of the possible wavefunctions. For this, let us set the $`AdS`$ radius to one, and analyze the behaviour for different relations between $`T`$ and $`L_p`$. In Fig. 3 we plot a schematic representation of the potential energy $`V_{\mathrm{eff}}=A^2(1/\mathrm{}^2+16\pi ^2T^2/L_p^2)`$. Classically, motion is confined to the region below the solid line (on which $`V_{\mathrm{eff}}=1`$). Strictly speaking, when $`T=2L_p`$ classical motion is only allowed for $`A>\alpha `$, while for $`T=L_p`$ the condition is $`A>\beta `$. In this region the wave function $`\psi `$ presents an oscillatory behaviour modulated by $`(1+A^2/\mathrm{}^2)^{1/4}`$, whereas from the turning point to zero, $`\psi `$ is exponentially decreasing. The complete shape of $`\psi `$ can be seen in Fig. 4. On the other hand, if $`16\pi ^2\mathrm{}^2T^2/L_p^2<1`$, $`V_{\mathrm{eff}}`$ remains greater than $`1`$ in the whole parameter space, and the classical motion is always forbidden. In this case, $`\psi `$ can be expressed in terms of exponentials with real arguments, yielding just vacuum fluctuations. This is of course consistent with the behaviour we saw for $`T=0`$.
### 2.3 WKB Approximation
In this subsection we discuss solutions which are valid in the near–classical domain. Since the potential is slowly varying (see Fig. 3), one expects the wave function to closely approximate the free particle state wavefunction $`\psi (A)=f(A)e^{ipA/\mathrm{}}`$. Thus, we will look for solutions of the form $`\psi (A)=f(A)e^{iS(A)/\mathrm{}}`$. Following , the semi-classical quantization condition may be written in the generalized form,
$$p(T,A)𝑑A=(n_A+\delta )\mathrm{},$$
(24)
where $`n_A`$ stands for the “radial” quantum number, and $`\delta `$ is related to the Maslov index . For a Hamiltonian quadratic in momenta, the usual WKB method shows that $`\delta `$ is typically a simple fraction. In other cases, $`\delta `$ depends on both the Hamiltonian and the boundary conditions and is often transcendental. In the present discussion, a precise calculation of $`\delta `$ would add little to our understanding, thus, it will not be evaluated but shall merely be carried along as an arbitrary constant.
The precise form of the WKB wavefunction is determined by the following constraint. In the semi–classical limit ($`\mathrm{}0`$), the classical average in time of any quantity $`Q(x)`$,
$$\overline{Q}(x)=\frac{1}{\tau }_0^\tau Q(x(t))𝑑t=\frac{1}{\tau }_0^\tau \frac{Q(x(t))}{v(x)}𝑑x,$$
(25)
has to be equal to the quantum average,
$$<\psi |Q|\psi >=\frac{|\psi (x)|^2Q(x)𝑑x}{|\psi (x)|^2𝑑x};$$
(26)
where $`v=/p`$, and the classical time average $`\tau =_0^\tau v^1𝑑x`$. Thus, the semi–classical approximation in the classically allowed region is given by,
$$\psi _{\mathrm{WKB}}(A)=\left|\frac{(p(T,A),A)}{p}\right|^{1/2}\mathrm{exp}\left[\pm \frac{i}{\mathrm{}}^Ap(T,x)𝑑x\right],$$
(27)
while in the classical forbidden region it reads,
$$\psi _{\mathrm{WKB}}(A)=\left|\frac{(p(T,A),A)}{p}\right|^{1/2}\mathrm{exp}\left[\pm \frac{1}{\mathrm{}}^Ap(T,x)𝑑x\right].$$
(28)
It is easily seen that for the typical Hamiltonian quadratic in momentum this generalized prescription reduces to the usual WKB approximation.
The conjugate momenta results in a multi-valued function:
$$p(T,A)=\pm \frac{L_p}{2}\left\{\mathrm{arccosh}\left[\frac{4\pi TA}{L_p\sqrt{1+A^2/\mathrm{}^2}}\right]+2\pi in\right\}.$$
(29)
Here $`\mathrm{arcosh}(x)`$ is taken to map $`[1,\mathrm{})[0,\mathrm{})`$, and $`\pm `$ refers to outgoing/ingoing directions. In this scheme the imaginary contribution to $`p(T,A)`$ does not contribute to the quantization condition. The quantum number $`n`$, however, does contribute when estimating the WKB wave function. In the classical allowed region we get,
$$\psi _{\mathrm{WKB}}(A)=\frac{\mathrm{exp}[n\pi A/L_p]}{|1A^2(16\pi ^2T^2/L_p^21/\mathrm{}^2)|^{1/4}}e^{\pm i\mathrm{\Theta }(A)},$$
(30)
where
$$\mathrm{\Theta }=\frac{1}{\mathrm{}}^A\frac{L_p}{2}\mathrm{arccosh}\left[\frac{4\pi Tx}{L_p\sqrt{1+x^2/\mathrm{}^2}}\right]𝑑x.$$
(31)
Note that in the limit $`A/\mathrm{}<<1`$,
$$\mathrm{\Theta }=\frac{A}{2L_p}\mathrm{arccosh}\left[\frac{4\pi TA}{L_p^2}\right].$$
(32)
Thus, we recover the behaviour found in the previous subsection, $`\psi `$ exponentially increases from zero to the turning point.
If we now flip $`TT`$, and use $`\mathrm{arcosh}(x)=\mathrm{arccosh}(x)+i\pi `$, we find that
$$\psi _{\mathrm{WKB}}(A)=\frac{\mathrm{exp}[(n+1)\pi A/L_p]}{|1A^2(16\pi ^2T^2/L_p^21/\mathrm{}^2)|^{1/4}}e^{\pm i\mathrm{\Theta }(A)}$$
(33)
are WKB eigenmodes corresponding to an eigenvacuumenergy $`T`$.
This semiclassical solution blows up at the turning points, where $`\dot{A}`$ goes to zero. This in itself may be tolerated if the wavefunction is normalizable. The matching of the wavefunction at the turning points may still be done by examining the wave equation more closely in the vicinity of the turning point.
## 3 Brane–world in $`AdS_{d+1}`$
### 3.1 Cosmology on the Brane
We turn now to a more general analysis independent of the dimension, i.e., for $`AdS_{d+1}`$ with $`d>1`$. The expression for the total action is given by,
$$S_{\mathrm{tot}}=\frac{L_p^{(3d)}}{16\pi }_\mathrm{\Omega }d^{d+1}x\sqrt{g}\left(R+\frac{d(d1)}{\mathrm{}^2}\right)+\frac{L_p^{(3d)}}{8\pi }_\mathrm{\Omega }d^dx\sqrt{\gamma }𝔎+T_\mathrm{\Omega }d^dx\sqrt{\gamma }.$$
(34)
Let us also generalize the possible symmetries on the bulk which yield different Robertson–Walker like cosmologies. The most general $`AdS_{d+1}`$ metric can be written as,
$$ds^2=\left(k+\frac{y^2}{\mathrm{}^2}\right)dt^2+\left(k+\frac{y^2}{\mathrm{}^2}\right)^1dy^2+y^2d\mathrm{\Sigma }_k^2,$$
(35)
where $`k`$ takes the values $`0,1,1`$ for flat, hyperbolic, or spherical geometries respectively and where $`d\mathrm{\Sigma }_k^2`$ is the corresponding metric on the unit $`(d1)`$-dimensional plane, hyperboloid, or sphere. It should be stressed that if $`k=1`$, an event horizon appears at $`y=\mathrm{}`$. With this in mind, one can trivially generalize the discussion in the appendix to get,
$$𝔎_{\varphi _i}^{\varphi _i}=\frac{1}{A}\left[k+\frac{A^2}{\mathrm{}^2}+\dot{A}^2\right]^{1/2},$$
(36)
and
$$𝔎_\tau ^\tau =\left[\ddot{A}+\frac{A}{\mathrm{}^2}\right]\left[k+\frac{A^2}{\mathrm{}^2}+\dot{A}^2\right]^{1/2},$$
(37)
where $`i`$ runs from 1 to $`(d1)`$. In terms of these quantities, the Einstein equation reads ,
$$Tg_{_{\mathrm{\Xi }\mathrm{{\rm Y}}}}\delta _A^^\mathrm{\Xi }\delta _B^^\mathrm{{\rm Y}}=\frac{L_p^{3d}}{4\pi }[𝔎_{AB}\mathrm{tr}(𝔎)g_{_{\mathrm{\Xi }\mathrm{{\rm Y}}}}\delta _A^^\mathrm{\Xi }\delta _B^^\mathrm{{\rm Y}}].$$
(38)
Its non–trivial components are,
$$T=\frac{L_p^{(3d)}}{4\pi }\frac{(d1)}{A}\left(k+\frac{A^2}{\mathrm{}^2}+\dot{A}^2\right)^{1/2},$$
(39)
and
$$T=\frac{L_p^{(3d)}}{4\pi }\left\{\frac{(d2)}{A}\left(k+\frac{A^2}{\mathrm{}^2}+\dot{A}^2\right)^{1/2}+\frac{\ddot{A}+A/\mathrm{}^2}{\sqrt{k+\dot{A}^2+A^2/\mathrm{}^2}}\right\}.$$
(40)
It is easily seen that Eqs. (39) and (40) imply the conservation of the stress energy. The evolution of the system is thus governed by,
$$\dot{A}^2=kA^2\left(\frac{1}{\mathrm{}^2}\frac{16\pi ^2T^2}{(d1)^2L_p^{2(3d)}}\right).$$
(41)
A somewhat unusual feature of brane physics can be analyzed from Eq. (41) (the five–dimensional case was already discussed by Kraus, Ref. ). Recall that in the spherical case, the classical behaviour of the brane is bounded by a minimum radius
$$A_0^2=\left(\frac{1}{\mathrm{}^2}+\frac{16\pi ^2T^2}{(d1)^2L_p^{2(3d)}}\right)^1,$$
(42)
but once the brane reaches that “size” it expands forever. Thus, contrary to the standard Robertson Walker cosmology, the spherically symmetric brane – corresponding to $`k=1`$ – represents an open world. Furthermore, depending on the value of $`T`$ we can also obtain a closed world with hyperbolic symmetry, i.e. with $`k=1`$. On the one hand, if
$$\frac{16\pi ^2T^2\mathrm{}^2}{(d1)^2L_p^{2(3d)}}1,$$
(43)
the classical solution does not have turning points yielding an open world. It should be remarked, however, that for $`k=1`$ the spacetime has no event horizons, whereas if $`k=1`$, the brane crosses an event horizon (at $`A=\mathrm{}`$) in a finite proper time.
On the other hand, if
$$\frac{16\pi ^2T^2\mathrm{}^2}{(d1)^2L_p^{2(3d)}}<1,$$
(44)
the classical solution has two turning points representing a big–bang and a big–crunch. Again, the spacetime has an event horizon at finite proper distance from the brane. If $`k=0`$, one obtains a solution only if the inequality (43) is satisfied. In the critical, case the solution represents the RS<sub>d+1</sub> brane–world. At this stage, it is noteworthy that a comprehensive analysis of a domain wall that inflates, either moving through the bulk or with the bulk inflating too, was first discussed by Chamblin–Reall .
### 3.2 Semiclassical Corrections
With the field equations for an expanding $`(d1)`$-brane in hand, the generalization of the WKB approximation to $`AdS_{d+1}`$ is straightforward. Of particular interest is $`AdS_5`$.<sup>7</sup><sup>7</sup>7Note that if $`k=0`$ and $`T=3/(4\pi L_p\mathrm{})`$, one recovers the RS–world. Let us specialize again to the case of a spherically symmetric brane. In such a case, Eq. (34) can be re–written as
$`S_{\mathrm{tot}}`$ $`=`$ $`{\displaystyle \frac{1}{L_p}}{\displaystyle }d\tau \{{\displaystyle \frac{A^3}{3\mathrm{}^2}}{\displaystyle \frac{\sqrt{\dot{A}^2+A^2/\mathrm{}^2+1}}{1+A^2/\mathrm{}^2}}+3A\sqrt{1+A^2/\mathrm{}^2+\dot{A}^2}`$ (45)
$``$ $`2A\dot{A}\mathrm{arcsinh}\left[{\displaystyle \frac{\dot{A}}{\sqrt{1+A^2/\mathrm{}^2}}}\right]\}+T{\displaystyle }_\mathrm{\Omega }d^4x\sqrt{\gamma }.`$
For positive eigenvalues of $`T`$, the solution in the classical allowed region is then given by,
$$\psi _{\mathrm{WKB}}(A)=\frac{\mathrm{exp}[2\pi n(A/L_p)^2]}{|1A^2/\mathrm{}^2+G^2)|^{1/4}}e^{\pm i^Ap𝑑x},$$
(46)
with $`p_5/\dot{A}`$, and $`G(A)=4\pi A^2TL_p/3`$. The oscillating part will be a real exponential term in the classically forbidden region.
## 4 Relation to $`AdS/CFT`$ Correspondence
### 4.1 Generalities
Another, seemingly different, but in fact closely related subject we will discuss in this section is the $`AdS/CFT`$ correspondence . This map provides a “holographic” projection of the $`AdS`$ gravitational system into the physics of the gauge theory. In the standard noncompact $`AdS/CFT`$ set up, gravity is decoupled from the dual boundary theory. The prime example here being the duality between Type IIB on $`AdS_5\times S^5`$ and $`𝒩=4`$ supersymmetric $`U(N)`$ Yang-Mills in $`d=4`$ with coupling $`g_{YM}`$ (the t’Hooft coupling is defined as $`\lambda =g_{YM}^2N`$). In this case it is known that the parameters of the $`CFT`$ are related to those of the supergravity theory by
$`\mathrm{}`$ $`=`$ $`\lambda ^{1/4}l_s`$ (47)
$`{\displaystyle \frac{\mathrm{}^3}{L_p^3}}`$ $`=`$ $`{\displaystyle \frac{2N^2}{\pi }},`$ (48)
where $`l_s`$ is the string length. The supergravity description is valid when $`\lambda `$ and $`N`$ are large (so that stringy effects are small). However, it is natural to suppose (in the spirit of $`AdS/CFT`$) that any RS-like model should properly be viewed as a coupling of gravity to whatever strongly coupled conformal theory the $`AdS`$ geometry is dual to. In the following discussion, inspired in , we unfold on this hypothesis: The most general action for a RS–like model in $`AdS_{d+1}`$ is given by
$$S_{RS}=S_{EH}+S_{GH}+2S_1+S_m,$$
(49)
where $`S_1`$ is the counterterm $`(T/2)d^dx\sqrt{\gamma }`$. The last term $`S_m`$ is the action for matter on the brane which was not included in Eq. (34), but it is included here for completeness. Now, to apply the $`AdS/CFT`$-correspondence, there is the question of the definition of the gravitational action in $`AdS_{d+1}`$. The standard action – corresponding to the two first terms of Eq. (34) – is divergent for generic geometries and one must add certain “counterterms” to obtain a finite action . Then we have schematically,
$$S_{\mathrm{grav}}=S_{EH}+S_{GH}+S_1+S_2+S_3+\mathrm{},$$
(50)
where $`S_k`$ is of order $`2(k1)`$ in derivatives of the boundary metric. Specifically, $`S_2`$ and $`S_3`$ are the counterterms discussed in . They are expressed in terms of the boundary metric:
$$S_2d^dx\sqrt{\gamma }\stackrel{~}{R}$$
(51)
and
$$S_3d^dx\sqrt{\gamma }\left(\stackrel{~}{R}_{ij}\stackrel{~}{R}^{ij}\frac{d}{4(d1)}\stackrel{~}{R}^2\right).$$
(52)
Some of the higher–order counterterms were computed in . For a given dimension $`d`$, however, one only needs to add a finite number of counterterms, specifically terms of order $`2n<d`$ in derivatives of the boundary metric.
In the counterterms were found for $`AdS_3`$, $`AdS_4`$ and $`AdS_5`$ by requiring a finite mass density of the spacetime. In the first case it was found, that only $`S_1`$ is needed, while in the latter cases both $`S_1`$ and $`S_2`$ are needed. Kraus et al. later derived a method for generating the required counterterms for any dimension $`d`$. Furthermore, in it was also noted that for the case of $`AdS_5`$ one could add terms of higher order in derivatives of the metric, as for example the counterterm $`S_3`$ but without changing the mass of the spacetime. Confronted with this ambiguity we face the question of which counterterms should be added in for example $`AdS_5`$. For that, we note that in order to apply the $`AdS/CFT`$–correspondence we should require that the symmetries on both sides of the correspondence match. The Weyl anomaly was computed in for gravity theories in $`AdS_{d+1}`$ and we can then apply this result to fix the possible counterterms. For $`d`$ odd there is no such anomaly and the divergent part of the (super)gravity action is canceled by the addition of the above mentioned counterterms. This implies, for example, that for $`AdS_4`$ we should only add $`S_1`$ and $`S_2`$. For $`d`$ even there is a nonvanishing anomaly . For $`AdS_3`$ this means that both $`S_1`$ and $`S_2`$ should be added and for $`AdS_5`$ we should add the terms $`S_1`$, $`S_2`$ and $`S_3`$. So, the requirement of finiteness of the action together with the matching of Weyl anomalies fixes the precise form of the supergravity action in $`AdS_{d+1}`$.
### 4.2 Dual Boundary Theory
Now, using the $`AdS/CFT`$-correspondence, one can easily show that the RS–model in dimension $`d+1`$ is dual to a $`d`$–dimensional $`CFT`$ (which we call the RS $`CFT`$) with a coupling to matter fields and the domain wall given by the action $`2S_2+2S_3+\mathrm{}+S_m`$, where we should remember that for $`AdS_3`$ and $`AdS_4`$, the $`S_3`$–term is absent but appears in all higher–dimensional cases. To illustrate this point, let us now analyze the $`AdS/CFT`$ for the simplest three–dimensional example. We will work in Euclidean space in order to avoid definition problems in the path integral. In this case the RS action (without matter) is given by
$$S_{RS}=\frac{L_p}{16\pi }_\mathrm{\Omega }d^3x\sqrt{g}\left(R+\frac{2}{\mathrm{}^2}\right)\frac{L_p}{8\pi }_\mathrm{\Omega }d^2x\sqrt{\gamma }𝔎\frac{L_p}{4\pi }_\mathrm{\Omega }d^2x\sqrt{\gamma },$$
(53)
which is essentially the same as in Eq. (1) but now with the tension $`T`$ fixed to be $`L_p/(4\pi )`$. (More on this below). Our set–up is as illustrated in Fig. 5: we have two regions $`R_1`$ and $`R_2`$ bounded by a two–dimensional domain wall and on each of these regions the metric is the $`AdS_3`$ metric $`g_{ij}`$ which induces the metric $`\gamma _{ij}`$ on the wall.<sup>8</sup><sup>8</sup>8For details of Penrose diagrams the reader is referred to . Following , let us compute the partition function obtained by integrating over the bulk metric with boundary value $`\gamma _{ij}`$ on the wall:
$$Z_{RS}[\gamma ]=e^{2S_1}\left(_{R_1R_2}𝒟ge^{S_{EH}[g]S_{GH}[g]}\right),$$
(54)
where the integral is over the two patches $`R_1`$ and $`R_2`$ of $`AdS`$. (Note that even though $`S_{GH}`$ is a two–dimensional term it depends on the bulk metric through the extrinsic curvature of the domain wall and can therefore not be taken out of the path integral). Since the integral over the two regions of $`AdS`$–space are independent, we can write it as an integral over a single patch of $`AdS`$–space<sup>9</sup><sup>9</sup>9 Note that the result of the integral over the regions $`R_1R_2`$ is not the addition of the integrals, but the product. Indeed, since we are dealing with independent processes, we have the product of the probabilities amplitudes instead of the sum, that would produce ‘interference effects’ not present in the RS set-up.:
$$Z_{RS}[\gamma ]=e^{2S_1}\left(_{R_1}𝒟ge^{S_{EH}[g]S_{GH}[g]}\right)^2.$$
(55)
Now according to the discussion above, the partition function for a consistent gravity theory in $`AdS_3`$, with finite mass of spacetime and appropriate central charge, is
$`Z_{\mathrm{grav}}[\gamma ]`$ $`=`$ $`{\displaystyle _{[\gamma ]}}𝒟ge^{S_{EH}[g]S_{GH}[g]S_1[\gamma ]S_2[\gamma ]}`$ (56)
$`=`$ $`e^{S_1[\gamma ]S_2[\gamma ]}{\displaystyle _{[\gamma ]}}𝒟ge^{S_{EH}[g]S_{GH}[g]}`$
$`=`$ $`e^{W_{CFT}[\gamma ]},`$
and according to the $`AdS/CFT`$ it should be identified with the generating functional for connected Green’s functions of the RS $`CFT`$ as above. By combining Eq. (55) and (56) we finally obtain:
$$Z_{RS}[\gamma ]=e^{2W_{CFT}[\gamma ]+2S_2[\gamma ]}.$$
(57)
This shows that the RS–like model in $`AdS_3`$ is equivalent to a $`CFT`$ coupled to gravity with action $`2S_2`$. This dual gravity theory is actually two–dimensional since $`2S_2`$ is the Einstein–Hilbert action for two–dimensional gravity. Similar correspondences can be derived in higher–dimensional cases. For example we have:
$$S_{RS}^{(4)}W_{RS}^{(4)}2S_2+S_m,$$
(58)
while
$$S_{RS}^{(5)}W_{RS}^{(5)}2S_22S_3+S_m.$$
(59)
Here $`W_{RS}`$ stands for the generating functional of connected Green’s functions of the boundary (RS) $`CFT`$, that is twice the CFT induced on the brane. Note that, as in the case of $`AdS_3`$, $`2S_2`$ is the Einstein-Hilbert action for $`d`$–dimensional gravity and so the RS model is equivalent to $`d`$-dimensional gravity coupled to a $`CFT`$ with corrections to gravity coming from the third counterterm $`S_3`$ (at least for $`d>3`$). This alone, however, does not tell us what the RS $`CFT`$ actually is<sup>10</sup><sup>10</sup>10The boundary $`CFT`$ can be found for the case of $`AdS_3`$., but rather that the RS model in $`d+1`$ dimensions can be viewed as a $`d`$-dimensional gravity (including corrections) coupled to a $`CFT`$ with matter.<sup>11</sup><sup>11</sup>11Related ideas were discussed in . And so, for example, in the case of $`AdS_5`$ this is another way to see why gravity is trapped on the four–dimensional domain wall and why there are corrections to Einstein gravity. (However, there are no such corrections in the case of $`AdS_3`$ and $`AdS_4`$ as we argued above).
### 4.3 Physical Implications
Up to this point we have kept the tension of the domain wall, $`T`$, arbitrary. Because of the various bounds described in sections 2 and 3 for different behaviours of the braneworld, it is important to see what one might expect. Let us again first restrict to $`AdS_3`$ for simplicity. It is well known that gravity in asymptotically $`AdS_3`$ spacetime has a holographic description as a 1+1 dimensional conformal field theory with central charge $`c=3\mathrm{}M_p/2`$ . In order to recover the geometry discussed in section 2, one must glue two copies of such bounded $`AdS_3`$ spacetimes, and then integrate over boundary metrics. Consequently, one has two copies of the matter action on the boundary, with total central charge $`c=3\mathrm{}M_p`$. In addition, if $`\stackrel{~}{R}>0`$ the conformal anomaly of the $`CFT`$ increases the effective tension on the domain wall, $`T>L_p/4\pi \mathrm{}`$, yielding a de Sitter universe with an effective cosmological constant driving inflation.<sup>12</sup><sup>12</sup>12A few words of caution; it is quite possible that the truncation from an infinite number of degrees of freedom down to only one degree of freedom, $`A(\tau )`$, has also drastically truncated the real physics. Unfortunately, a treatment using Wheeler’s full superspace is beyond the scope of our present calculation abilities. An (early) inflationary epoch looks very promising. The tremendous expansion during inflation may blow up a small sized region of the world (which was causally connected before inflation) to a size much greater than our current horizon. Therefore, it can be expected that the observable part of the brane looks smooth and flat, regardless of the initial curvature of the brane that inflated.<sup>13</sup><sup>13</sup>13Note that a flat Robertson-Walker Universe requires a total energy density equal to the critical density $`\rho _{\mathrm{cr}}`$, whereas ordinary matter contributes only about a 5% of $`\rho _{\mathrm{cr}}`$. A novel solution to this problem consistent with a large body of observations is the so-called “Manifold Universe” . Furthermore, if we consider conformal matter on the brane the inflationary phase is unstable and could decay into a matter dominated universe with thermalized regions, in agreement with current observations .
Another interesting process which could lead to brane–world reheating is as follows: During inflation trapped regions of false vacuum (within their Schwarzchild radii) caught between bubbles of true vacuum may give rise to the creation of primordial black strings. Now, it is well–known that the black string solution suffers from a Gregory–Laflamme instability leading to the formation of stable black cigars on the brane. In addition, it was shown in that the nucleation of supermassive bulk black holes is highly supressed compared to the above mentioned process. Thus, prompted by the conventional arena , one could speculate that the Hawking–evaporation of primordial black cigars slows down inflation. On the other hand, one could assume the existence of such a bulk black hole. Even in this case, the (brane-world/bulk-black-hole) system evolves towards a configuration of thermal equilibrium as was recently shown in .
Let us now briefly discuss a general $`n`$-dimensional brane-world that falls under the action of a higher dimensional gravitational field. The system can be decomposed into falling shells (which do not interact with each other or with the environment that generates the metric), with trajectories described by the scale factor $`A(\tau )`$. From the above discussion it is clear that the value of $`T`$ will depend on the symmetries of the domain wall. It is easily seen, for instance, that if $`k=1`$ then
$$T<\frac{(d1)L_p^{3d}}{4\pi \mathrm{}},$$
(60)
yielding a closed universe. Roughly speaking, the cosmological constant induced by the conformal anomaly accelerates/slows down the brane to balance the null geodesic congruence in the bulk, shirking the world’s pinch off. We recall that if $`k=1`$, the spacetime has an undesirable event horizon that must be reached by the brane in a finite proper time.
Despite the fact that it is contrary to the spirit of RS-worlds, it would be nice to add “matter fields” in the bulk to study the quantum cosmology and the dual $`CFT`$ coupled to gravity that in this case should be deformed by the insertion of operators.
Even though many kind of interesting phenomena are recognized, brane-world cosmology remains thoroughly non-understood. The lower dimensional model here discussed can hopefully illuminate the “physical $`AdS_5`$ cosmology”.
### Acknowledgements
We would like to thank Vijay Balasubramanian, Thomas Hertog, Finn Larsen, Juan Maldacena, Francisco Villaverde, Matt Visser, and Allan Widom for useful discussions/correspondence. Special thanks go to Harvey Reall for a critical reading of the manuscript and valuable comments. The work of LA and CN was supported by CONICET Argentina, and that of KO by the Danish Natural Science Research Council.
## Appendix A Appendix
Here we present a calculation of the second fundamental form of the metric in Eq. (2) (it should be remarked that this calculation is a direct analog to that of Ref. , and it is included just for the sake of completeness).
Let us start by introducing a Gaussian normal coordinate system in the neighborhood of the brane. We shall denote the one–dimensional surface swept out by the brane by $`\mathrm{\Sigma }`$. Let us introduce a coordinate system $`\varphi _{}`$ on $`\mathrm{\Sigma }`$. Next we consider all the geodesics which are orthogonal to $`\mathrm{\Sigma }`$, and choose a neighborhood $`N`$ around $`\mathrm{\Sigma }`$ so that any point $`pN`$ lies on one, and only one, geodesic. The first coordinate of $`p`$ is determined by the intersection of this geodesic with $`\mathrm{\Sigma }`$. The full set of spatial coordinates is then given by $`(\varphi _{};\eta )`$, while the surface $`\mathrm{\Sigma }`$ under consideration is taken to be located at $`\eta =0`$ so that Eq. (2) can be rewritten as
$$ds^2=\left(1+\frac{y^2}{\mathrm{}^2}\right)dt^2+d\eta ^2+y^2d\varphi ^2,$$
(1)
fixed by the relation, $`dy/d\eta =(1+y^2/\mathrm{}^2)^{1/2}`$. The second fundamental form in such a coordinate–system reads
$$𝔎_{\mu \nu }\frac{1}{2}\frac{g_{\mu \nu }}{\eta }|_{\eta =0,y=A},$$
(2)
and its non-trivial components are
$$𝔎_t^t=\frac{A}{\mathrm{}^2}\left(1+\frac{A^2}{\mathrm{}^2}\right)^{1/2},$$
(3)
$$𝔎_\varphi ^\varphi =\frac{1}{A}\left(1+\frac{A^2}{\mathrm{}^2}\right)^{1/2}.$$
(4)
To analyze the dynamics of the system, we permit the radius of the brane to become a function of time $`AA(\tau )`$. Recall that the symbol $`\tau `$ is used to denote proper time as measured by co–moving observers on the brane–world. Let the position of the brane be described by $`x^\mu (\tau ,\varphi )(t(\tau ),A(\tau ),\varphi )`$, so that the velocity of a piece of stress-energy at the brane $`(u^\mu u_\mu =1)`$ is
$$u^\mu \frac{dx^\mu }{d\tau }=(\frac{dt}{d\tau },\frac{dA}{d\tau },0).$$
(5)
We remind the reader that
$$ds^2=\left(1+\frac{A^2}{\mathrm{}^2}\right)dt^2+\left(\frac{dA}{dt}\right)^2\left(1+\frac{A^2}{\mathrm{}^2}\right)^1dt^2+A^2d\varphi ^2$$
(6)
so,
$$d\tau ^2=dt^2\left\{\left(1+\frac{A^2}{\mathrm{}^2}\right)+\left(\frac{dA}{dt}\right)^2\left(1+\frac{A^2}{\mathrm{}^2}\right)^1\right\}$$
(7)
or equivalently,
$$d\tau ^2=dt^2\left\{\left(1+\frac{A^2}{\mathrm{}^2}\right)^2+\left(\frac{dA}{dt}\right)^2\right\}\left(1+\frac{A^2}{\mathrm{}^2}\right)^1.$$
(8)
Since
$$\frac{dA}{dt}=\frac{dA}{d\tau }\frac{d\tau }{dt},$$
(9)
we first get,
$$\left(1+\frac{A^2}{\mathrm{}^2}\right)\left(\frac{d\tau }{dt}\right)^2=\left(1+\frac{A^2}{\mathrm{}^2}\right)^2+\dot{A}^2\left(\frac{d\tau }{dt}\right)^2$$
(10)
and then,
$$\frac{dt}{d\tau }=\frac{\sqrt{\dot{A}^2+A^2/\mathrm{}^2+1}}{1+A^2/\mathrm{}^2}.$$
(11)
Let us denote by $`\widehat{n}^\mu `$ the unit normal vector to the brane, which satisfies $`u^\mu \widehat{n}_\mu =0`$ and $`\widehat{n}^\mu \widehat{n}_\mu =1`$; its components are $`\widehat{n}^\mu =(\dot{A}/(1+A^2/\mathrm{}^2),(1+A^2/\mathrm{}^2+\dot{A}^2)^{1/2},0)`$, such that the coordinate $`y`$ is increasing in the direction $`\widehat{n}^\mu `$. Thus we obtain
$$𝔎_\varphi ^\varphi =\frac{1}{y}\frac{y}{\eta }|_{y=A}=\frac{1}{A}\left(1+\frac{A^2}{\mathrm{}^2}+\dot{A}^2\right)^{1/2}.$$
(12)
To evaluate $`𝔎_\tau ^\tau `$ one can proceed in two alternative ways. First one can simply use the definition $`𝔎_{\mu \nu }=\frac{1}{2}_{(\mu }\widehat{n}_{\nu )}`$, giving:
$$𝔎_{tt}=\frac{1}{2}_{(t}\widehat{n}_{t)}=\frac{d\widehat{n}_t}{d\tau }\frac{d\tau }{dt}\mathrm{\Gamma }_{tt}^\eta \widehat{n}_\eta =\frac{1+A^2/\mathrm{}^2}{\sqrt{1+\dot{A}^2+A^2/\mathrm{}^2}}(\ddot{A}+A/\mathrm{}^2),$$
(13)
that using
$$𝔎_{\tau \tau }=\frac{x^\mu }{x^\tau }\frac{x^\nu }{x^\tau }𝔎_{\mu \nu },$$
(14)
immediately yields
$$𝔎_\tau ^\tau =𝔎_t^t=\frac{\ddot{A}+A/\mathrm{}^2}{\sqrt{1+\dot{A}^2+A^2/\mathrm{}^2}}.$$
(15)
Alternatively, one can easily check this last result by observing that
$$𝔎_\tau ^\tau 𝔎_{\tau \tau }=u^\mu u^\nu 𝔎_{\mu \nu }=u^\mu u^\nu _\mu \widehat{n}_\nu =u^\mu \widehat{n}_\nu _\mu u^\nu =\widehat{n}_\mu (u^\nu _\nu u^\mu )=\widehat{n}_\mu q^\mu ,$$
(16)
where $`q^\mu `$ is the four acceleration of the brane. Now, by the spherical symmetry of the problem the four acceleration is proportional to the unit normal, $`q^\mu q\widehat{n}^\mu `$, so $`𝔎_\tau ^\tau =q`$. To explicitly evaluate the four acceleration, utilize the fact that $`\xi ^\mu _t^\mu (1,0,0)`$ is a Killing vector for the underlying geometry. At the brane, the components of this vector are $`\xi _\mu =([1+A^2/\mathrm{}^2],0,0)`$, so that $`\xi _\mu \widehat{n}^\mu =\dot{A}`$ and $`\xi _\mu u^\mu =(1+A^2/\mathrm{}^2+\dot{A}^2)^{1/2}`$. With this in mind, comparing
$$\frac{d}{d\tau }(\xi _\mu u^\mu )=\xi _\mu q\widehat{n}^\mu =q\dot{A},$$
(17)
and
$$\frac{d}{d\tau }(\xi _\mu u^\mu )=\dot{A}\frac{A/\mathrm{}^2+\ddot{A}}{\sqrt{1+A^2/\mathrm{}^2+\dot{A}}},$$
(18)
we get
$$𝔎_\tau ^\tau =\frac{A/\mathrm{}^2+\ddot{A}}{\sqrt{1+A^2/\mathrm{}^2+\dot{A}}}=\frac{d}{d\tau }\left\{\mathrm{arcsinh}\left[\frac{\dot{A}}{\sqrt{1+A^2/\mathrm{}^2}}\right]\right\}+\frac{A}{\mathrm{}^2}\frac{dt}{d\tau };$$
(19)
this result agrees with that of Eq. (15). Having calculated the nontrivial components of the second fundamental form we can now derive a simpler expression for the relevant gravity–action (1) in $`AdS_3`$. Since $`\sqrt{g}d^3x2\pi AdAdt`$ and $`\sqrt{\gamma }d^2x2\pi Ad\tau `$ an integration by parts finally leads to
$$S_{\mathrm{gravity}}=\frac{L_p}{2}𝑑\tau \left\{\dot{A}\mathrm{arcsinh}\left[\frac{\dot{A}}{\sqrt{1+A^2/\mathrm{}^2}}\right]+\sqrt{1+\frac{A^2}{\mathrm{}^2}+\dot{A^2}}\right\}.$$
(20)
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# Exact Solutions with w - modes
## I Introduction
It is now well known that sufficiently compact stars (polytropic or uniform density) can support the internal trapping of null geodesics and the ”w-modes” found by Chandrasekhar and Ferrari (1991), . These modes exist both for axial and polar perturbations, though the axial ones have been studied more thoroughly. The w-modes in general have no Newtonian counterparts since they are predominantly modes of the spacetime. In the polar case they couple weakly to the fluid while in the axial case there is no coupling at all. Recent numerical studies of these w-modes have involved the effect of the equation of state and their excitation ,. Whereas the role that these w-modes may play in real astrophysical processes remains open to much further investigation, it is fair to say that little is actually known about the behavior of the governing potential of the wave equation in exact solutions of Einstein’s equations. Such knowledge is important since it is both a route to the physical understanding of relativistic phenomena and a check on numerical procedures. The purpose of this paper is to explore necessary conditions for the internal trapping of null geodesics and the existence of w-modes (when the centrifugal part of the potential dominates) in physically acceptable exact isolated static spherically symmetric perfect fluid solutions of Einstein’s equations. We are able to exhibit physically acceptable exact solutions which have trapping and which could support w-modes.
## II Review of Perfect Fluids
Any metric is an “exact” solution to Einstein’s equations. However, the consequent energy-momentum tensor is almost never of any interest. What is of interest are solutions which might have some contact with reality. Recently a collection of exact isolated static spherically symmetric perfect fluid solutions have been subjected to the following elementary criteria for physical acceptability:
* Isotropy of the pressure ($`p`$).
* Regularity of the origin by way of the scalars polynomial in the Riemann tensor , , .
* Positive definiteness of both $`p`$ and energy density ($`\rho `$) at the origin.
* Isolation by way of the requirement that the pressure reduce to zero at some finite boundary radius $`r_\mathrm{\Sigma }>0`$.
* Monotonicity of both $`p`$ and $`\rho `$ to the boundary.
* Subluminal adiabatic sound speed ($`\text{v}_s^2=\frac{dp}{d\rho }<1`$) .
Perhaps not surprisingly, only about 10 % of the solutions pass these elementary tests. In what follows we take the view that solutions worthy of further consideration must pass all the applicable tests in at least some region . We also take the view that an analytic solution of Einstein’s equations can be expected to approximate only a region of a realistic configuration. That is, an analytic solution could have an interior causal limit ($`\text{v}_s^2=1`$), a circumstance which precludes standard stability arguments , and yet provide an adequate approximation for a region of a realistic configuration.
We begin by setting the notation. The line element in conventional form is (e.g., )
$$ds^2=\frac{dr^2}{1\frac{2m(r)}{r}}+r^2(d\theta ^2+sin(\theta )^2d\varphi ^2)e^{2\mathrm{\Phi }(r)}dt^2$$
(1)
with the coordinates comoving in the sense that the fluid streamlines are given by $`u^a=e^{\mathrm{\Phi }(r)}\delta _t^a`$. In terms of the functions $`\mathrm{\Phi }(r)`$ and $`m(r)`$ the regularity conditions reduce to
$$\mathrm{\Phi }^{}(0)=m(0)=m^{}(0)=0,$$
(2)
with $`{}_{}{}^{}d/dr`$ and $`\mathrm{\Phi }(0)`$ a constant fixed by the scale of $`t`$. Next, in terms of the perfect fluid decomposition ($`T_b^a=(\rho (r)+p(r))u^au_b+p(r)\delta _b^a`$), solving for $`\mathrm{\Phi }^{}(r)`$ from the $`r`$-component of the conservation equations and Einstein’s equations ($`_aT_r^a=0`$ and $`G_r^r8\pi p(r)=0`$) we obtain the Tolman -Oppenheimer-Volkoff (TOV) equation
$$\mathrm{\Phi }^{}(r)=\frac{p^{}(r)}{\rho (r)+p(r)}=\frac{m(r)+4\pi p(r)r^3}{r(r2m(r))},$$
(3)
where, from the $`t`$ component of the Einstein equations ($`G_t^t=8\pi \rho (r)`$),
$$4\pi \rho (r)=\frac{m^{}(r)}{r^2}.$$
(4)
From the TOV equation (here taken to be the right hand members of (3)) we observe that $`p(r)`$ is maximal at $`r=0`$. Moreover, if there is an equation of state ($`p(\rho )`$) then either $`\rho `$ is maximal wrt $`r`$ at $`r=0`$ or $`p`$ is maximal wrt $`\rho `$ at $`r=0`$. Despite that fact that the TOV equation has been known for over sixty years, only recently has its mathematical structure been fully appreciated. For example, we now know that for $`p(r)>0`$ there exists a unique global solution for every $`0<p(0)<\mathrm{}`$. It is not difficult to find “solutions” of the TOV equation. For example, $`m(r)`$ can be chosen in such a way that (3) yields a solution (with $`\rho (r)`$ following from (4)). The simplest choice is clearly $`mr^3`$ but this leads us back to the Schwarzschild interior solution. The metric (1) contains two functions, $`m(r)`$ and $`\mathrm{\Phi }(r)`$, related by (3). The first represents the gravitational energy (effective gravitational mass) (e.g., ). The second is, in the weak field limit $`r2m(r)`$, the Newtonian potential. This interpretation offers no insight into the meaning of $`\mathrm{\Phi }(r)`$ within Einstein’s theory, and is a good point to begin our discussion.
## III Null Geodesic Limit
We start with the “centrifugal” part of the potential for non-radial odd parity perturbations. This governs the evolution of null geodesics. Radial null geodesics of the metric (1) satisfy
$$t=\pm \frac{dr}{e^{\mathrm{\Phi }(r)}\sqrt{1\frac{2m(r)}{r}}}+\mathrm{D},$$
(5)
with $`\theta `$, $`\varphi `$, and D constant. Non-radial null geodesics satisfy $`\theta =\pi /2`$ (by choice),
$$r^4\varphi ^2=1,$$
(7)
$$e^{4\mathrm{\Phi }(r)}t^2=\frac{1}{b^2},$$
(8)
and
$$r^2r^2=(1\frac{2m(r)}{r})(\frac{B(r)^2}{b^2}1),$$
(9)
with
$$B(r)re^{\mathrm{\Phi }(r)},$$
(10)
where $`d/d\lambda `$ for affine $`\lambda `$, and $`b`$ is a constant $`>0`$, the “impact parameter”. The “potential” impact parameter $`B(r)`$ provides, by way of (10), an invariant physical interpretation of $`\mathrm{\Phi }(r)`$. Null geodesics are restricted by the condition $`bB(r)`$ .
From conditions (2) and the definition (10) it follows that
$$B(r)\psi _0r$$
(11)
as $`r0`$ where $`\psi _0`$ is a physically irrelevant scale factor. (The ratio $`B(r)/b`$ is invariant to scale changes in $`t`$.) It follows that the necessary and sufficient condition for the internal trapping of null geodesics (that is the existence of $`r_0`$ such that $`r^{}=0`$ and $`r^{}<0`$ at $`r_0`$) is given by
$$\mathrm{\Phi }^{}(r)>\frac{1}{r}$$
(13)
or, from (3),
$$p^{}(r)<\frac{\rho (r)+p(r)}{r}$$
(14)
which, with an equation of state ($`p(\rho )`$) can be given as
$$\rho ^{}(r)<\frac{\rho (r)+p(r)}{v_s^2r}.$$
(15)
From (13) and the TOV equation it follows that
$$r<3m(r)+4\pi p(r)r^3,$$
(16)
a relation which makes the trapping of null geodesics a manifestly relativistic phenomenon , .
## IV Full Potential
The odd parity (axial) w-modes are non-radial perturbations of the spacetime which do not couple to the fluid at all. In terms of the frequency $`\varpi `$ and mode number $`l2`$ the governing equation is given by
$$(\frac{d^2}{dr_{}^{}{}_{}{}^{2}}+\varpi ^2)Z=V(r_{})Z,$$
(17)
where $`r_{}`$ is the “tortoise” coordinate
$$dr_{}=\frac{e^{\mathrm{\Phi }(r)}}{\sqrt{1\frac{2m(r)}{r}}}dr.$$
(18)
The potential is conveniently expressed in terms of $`r`$ and is given by
$$V(r)=\frac{1}{B(r)^2}(l(l+1)+4\pi r^2(\rho (r)p(r))\frac{6m(r)}{r}).$$
(19)
A necessary condition for the occurrence of resonance scattering of axial gravitational waves by an isolated distribution of fluid is a local minimum in $`V(r)`$ within the boundary of the fluid. (If the centrifugal part of the potential ($`\frac{1}{B(r)^2}(l(l+1)`$) dominates, which is frequently but not always the case (see below), then (16) provides such a condition.) It is the purpose of this paper to explore the occurences of this minimum in physically acceptable exact solutions. It is the shape of the function $`V(r)`$ which is of interest, and since the exterior vacuum has a well known local maximum at $`r3.28M`$ (for $`l=2`$), the boundary conditions associated with the fluid - vacuum interface need careful attention.
## V Boundary Conditions
The Darmois-Israel junction conditions demand the continuity of the first and second fundamental forms at a boundary surface. These conditions are well known (e.g., ) but are usefully reviewed here. We take the “interior” metric to be of the form (1). The “exterior” is the familiar Schwarzschild vacuum (in coordinates ($`\text{r}r,\theta ,\varphi ,Tt`$)):
$$ds^2=\frac{d\text{r}^2}{1\frac{2M}{\text{r}}}+\text{r}^2(d\theta ^2+sin(\theta )^2d\varphi ^2)(1\frac{2M}{\text{r}})dT^2.$$
(20)
At the fluid interface ($`\mathrm{\Sigma }`$), without loss in generality, we take $`\theta `$ and $`\varphi `$ continuous (with intrinsic coordinates $`\theta ,\varphi ,\tau `$, where $`\tau `$ is the proper time). This gives
$$\text{r}_\mathrm{\Sigma }=r_\mathrm{\Sigma }.$$
(21)
The continuity of the first fundamental form is completed by requiring that the particle trajectories at the boundary be timelike. The continuity of the angular components of the second fundamental form (extrinsic curvature) give
$$M=m(r_\mathrm{\Sigma }),$$
(22)
and the continuity of the remaining ($`\tau `$ \- $`\tau `$) component gives
$$\mathrm{\Phi }_\mathrm{\Sigma }^{^{}}=\frac{M}{r_\mathrm{\Sigma }(r_\mathrm{\Sigma }2M)}$$
(23)
which, with the TOV equation, gives
$$p(r_\mathrm{\Sigma })=0.$$
(24)
To summarize, a static spherically symmetric fluid is matched to a vacuum exterior subject to (and only to) (21), (22) and (24). Further restrictions are frequently imposed. In particular, if the coordinates are assumed admissible (the metric and first derivatives assumed continuous across $`\mathrm{\Sigma }`$) then
$$e^{2\mathrm{\Phi }(r_\mathrm{\Sigma })}=12\frac{M}{r_\mathrm{\Sigma }},$$
(25)
and
$$m_\mathrm{\Sigma }^{^{}}=0=\rho (r_\mathrm{\Sigma }).$$
(26)
Whereas (25) can be achieved by a simple change in scale (of $`t`$ or $`T`$), in general, (26) does not hold . Condition (25) is the necessary and sufficient condition for $`B`$ to be continuous and continuously differentiable at $`\mathrm{\Sigma }`$. Similarly, a simple change in scale makes $`V`$ continuous but not continuously differentiable at $`\mathrm{\Sigma }`$. The wave equation (17) is of course invariant to these changes in scale. In summary, $`B`$ can be taken to be continuous and continuously differentiable at $`\mathrm{\Sigma }`$, and $`V`$ can be taken to be continuous.
## VI Examples
Since the uniform density static sphere satisfies (16), one might guess that all static solutions do. This is not the case. For example, the Buchdahl solution does not allow a region which satisfies (16). In contrast, the Tolman VII solution does . (These are useful exact solutions for the study of the equation of state of neutron stars ). In what follows we demonstrate a number of physically acceptable solutions which do satisfy (16). We organize the examples by way of their motivating ansatz.
### A Prescribed form of $`m(r)`$
The Finch-Skea solution is an exact solution which gives reasonable values for the central densities of neutron stars. The solution derives from the ansatz
$$m(r)=\frac{Cr^3}{2(1+Cr^2)},$$
(27)
where $`C`$ is a constant. The line element can be given in the form
$$ds^2=v^2dr^2+r^2d\mathrm{\Omega }^2A^2((C_2C_1v)\mathrm{cos}(v)+(C_1+C_2v)\mathrm{sin}(v))^2dt^2,$$
(28)
where $`v\sqrt{1+\omega ^2}`$, $`\omega ^2Cr^2`$ and $`A,C_1`$ and $`C_2`$ are constants. Clearly $`AC_2`$ can be set by the scale of $`t`$ leaving (say) $`C`$ and $`\beta \frac{C_1}{C_2}`$ as parameters. The latter is conveniently given by
$$\beta =\frac{v_\mathrm{\Sigma }tan(v_\mathrm{\Sigma })1}{tan(v_\mathrm{\Sigma })+v_\mathrm{\Sigma }}$$
(29)
where $`v_\mathrm{\Sigma }\sqrt{1+\omega _\mathrm{\Sigma }^2}`$, or equivalently, in terms of the tenuity $`\alpha \frac{r_\mathrm{\Sigma }}{M}`$,
$$\alpha =\frac{2v_\mathrm{\Sigma }^2}{v_\mathrm{\Sigma }^21}.$$
(30)
The physical restrictions 3 and 6 give, respectively, the following lower and upper bounds to $`\beta `$
$$0.218\beta 5.605,$$
(31)
but the limits which follow from $`B`$ and $`V`$ are more transparently expressed in terms of $`\alpha `$. Up to an irrelevant scale factor, the potential impact parameter follows immediately as,
$$B(\omega )=\frac{w}{(1\beta v)cos(v)+(\beta +v)sin(v)}.$$
(32)
We find that $`B`$ has a local minimum for $`\alpha <3`$ and a local maximum with subluminal sound speed between the local maximum and minimum for $`\alpha >2.768`$. Some typical plots of $`B`$ are shown in Fig. 1. The full potential (up to an irrelevant scale factor) is given by
$$V(\omega )=\frac{l(l+1)+\frac{F(v)\omega ^2}{2}3\frac{\omega ^2}{1+\omega ^2}}{B^2}$$
(33)
where
$$F(v)=\frac{2+v^2}{v^4}+\frac{1}{v^2}\frac{(\beta v+1)+(\beta v)tan(v)}{(\beta v1)(\beta +v)tan(v)}.$$
(34)
Some typical plots of $`V`$ are shown in Fig. 2. We find that there is a local minimum in $`V`$ (with $`l=2`$) for $`\alpha <2.933`$ and the local minimum lies in a region with subluminal sound speed for $`\alpha >2.755`$ . The Finch-Skea solution therefore offers an example of a causal exact solution with trapping .
### B Prescribed form of $`\mathrm{\Phi }(r)`$
A class of models, some of which satisfy conditions 1 through 6, starts with the ansatz
$$e^{2\mathrm{\Phi }(r)}=D(1+Er^2)^n,$$
(35)
where $`D`$ and $`E`$ are constants and $`n`$ is an integer $`1`$. The case $`n=1`$ is known as Tolman IV solution . It follows immediately from (13) that this solution exhibits no trapping ($`B`$ is monotone increasing and $`V`$ monotone decreasing). For $`n=2`$ condition 3 fails. The case $`n=3`$ satisfies conditions 1 through 6. It has been examined by Heintzmann , who gives the solution
$$ds^2=\frac{dr^2}{(1\frac{3ar^2}{2}\frac{1+C(1+4ar^2)^{1/2}}{1+ar^2})}+r^2d\mathrm{\Omega }^2A^2(1+ar^2)^3dt^2.$$
(36)
Again, in terms of the tenuity ($`\alpha \frac{r_\mathrm{\Sigma }}{M}`$), we find that there is a local minimum in $`V`$ (with $`l=2`$) for $`\alpha <2.902`$ and the local minimum lies in a region with subluminal sound speed for $`\alpha >2.788`$. A typical example is shown in Fig. 3 where $`V`$ has been matched onto the vacuum exterior at the boundary $`\mathrm{\Sigma }`$, and the minimum in $`V`$ and sound speed limit have been indicated. The cases $`n=4`$ and $`n=5`$ satisfy conditions 1 through 6 and have been solved by Durgapal following the formulation of Korkina . For $`n=4`$ the solution is given by
$$ds^2=\frac{(1+Cr^2)^2dr^2}{(\frac{710Cr^2C^2r^4}{7}+\frac{KCr^2}{(1+5Cr^2)^{\frac{2}{5}}})}+r^2d\mathrm{\Omega }^2A(1+Cr^2)^4dt^2.$$
(37)
In this case we find that there is a local minimum in $`V`$ (with $`l=2`$) for $`\alpha <2.892`$ and the local minimum lies in a region with subluminal sound speed for $`\alpha >2.780`$. Since $`B(r)`$ has a local maximum up to $`\alpha 3`$, it is clear that the dynamical part of the potential $`V`$ can dominate. For $`n=5`$ the solution is
$$ds^2=\frac{(1+Cr^2)^3dr^2}{(1\frac{Cr^2(309+54Cr^2+8C^2r^4)}{112}+\frac{KCr^2}{\sqrt[3]{1+6Cr^2}})}+r^2d\mathrm{\Omega }^2A(1+Cr^2)^5dt^2.$$
(38)
We find similar results in this case. There is a local minimum in $`V`$ (with $`l=2`$) for $`\alpha <2.886`$ and the local minimum lies in a region with subluminal sound speed for $`\alpha >2.776`$. The solutions (36), (37) and (38) cannot represent the core region where, we note, they are acausal .
## VII Discussion
Condition (16) is proposed as the necessary condition for the internal trapping of null geodesics, and for the occurrence of resonance scattering of axial gravitational waves when the centrifugal term dominates the potential, in static spherically symmetric perfect fluids . This condition is not always satisfied. For example, it is not satisfied in the Buchdahl solution. We have demonstrated some physically acceptable exact solutions for which the condition is satisfied. One, the Finch-Skea solution, offers an example of a complete causal exact solution with trapping. At the very least, these examples can provide a check on numerical procedures which attempt to gauge the role that w-modes may play in real astrophysical processes. In every case studied here we have found that for resonance scattering the tenuity ($`\alpha =r_\mathrm{\Sigma }/M`$) lies in the small range $`2.8<\alpha <2.9`$. (In all cases, as $`\alpha `$ decreases, the causal boundary ($`\text{v}_s^2=1`$), if it exists, moves out and approaches the minimum in $`V`$ for some minimum $`\alpha `$, exactly as expected.) Whereas this range is well above the Buchdahl limit of $`9/4`$ , it is too low for, say, neutron stars . (Unphysical solutions with trapping and $`\alpha >5`$ are known .) It is reasonable to suggest, as has often been done, that an exact solution may reflect only part of a more realistic configuration. Boundary conditions within a distribution are easily derived from the discussion given in section V: $`p(r)`$ and $`m(r)`$ must be continuous to avoid surface layers (shells). In particular, $`m(r)`$ need not be continuously differentiable (which, at least formally, allows first-order phase transitions). All that is needed to raise $`\alpha `$ into a more interesting range (say $`3<\alpha <10`$) is the addition of an envelope. The envelope is constructed subject to the continuity of $`p(r)`$ and $`m(r)`$ at $`r_\sigma `$, where $`r_\sigma `$ is exterior to the minimum in $`V`$, and must allow $`p(r_\mathrm{\Sigma })=0`$ at a finite boundary $`r_\mathrm{\Sigma }`$ where $`r_\sigma <r_\mathrm{\Sigma }`$ .
## Acknowledgments
We thank Kostas Kokkotas, James Lattimer, Eric Poisson and Kjell Rosquist for helpful comments. This work was supported by a grant (to KL) from the Natural Sciences and Engineering Research Council of Canada.
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# 1 Introduction
## 1 Introduction
In deep-inelastic lepton-proton scattering processes, a small excess of positively charged leading particles is expected owing to the valence quark structure of the proton. An interaction involving a valence $`u`$ quark is 8 times more likely than that involving a valence $`d`$ quark. No excess is expected from sea quark or gluon initiated interactions.
Since leading hadrons “remember“ the charge of the struck quark, this gives the experimentalist a probe to investigate the partonic structure of the proton. At HERA this method could be used to perform such studies over several orders of magnitude in $`x`$ and $`Q^2`$ and will complement existing studies based on jet measurements and other hadronic final state variables.
### 1.1 Definitions
To define an event charge, let us first determine a reference direction on which to project momenta. If the DIS scattering were simply an elastic deflection of a quark by a positron, this direction would simply be the final quark direction. But we know that the reaction is more complicated, involving gluons or even more than one quark (boson-gluon-fusion BGF). Nevertheless, from the deflected positron parameters (which, for simulated events are known without error) we define the ’hadronic axis’ by assuming the scattering to be elastic, and using the deflected positron information.
Alternative methods could be to use event shape information, such as the thrust axis, to define the hadronic system. This is beyond the scope of this paper and is not considered here.
Event charge studies have already been made in deep inelastic scattering and in e+e- interactions at LEP.
For this paper, the following definition of event charge is used: $`Q=_iP_i^{||}q_i/_iP_i^{||}`$, where $`q_i`$ is the particle charge (-1 or +1) and $`P_i^{||}`$ is the projected momentum on the hadronic axis .The sums run over $`i`$ all particles which have positive $`P_i^{||}`$’s; Q is between -1 and +1. This was used by Aleph in .
### 1.2 Relation of event charge Q to scattering on valence quarks. <br>The Leading particle and the positive excess $`\epsilon _+`$
For this study, the leading particle is defined as the charged hadron which has the largest momentum projected onto the hadronic axis. If the experimentally measurable $`K_s^0`$ or $`\mathrm{\Lambda }`$ were to be the leading particle then the event is discarded.
The positive excess $`\epsilon _+`$ is defined as the ratio between the number of leading particles which are positively signed $`N_+`$ over those which are negatively signed $`N_{}`$: $`\epsilon _+=\frac{N_+}{N_{}}`$.
#### 1.2.1 Relation of $`\epsilon _+`$ to interactions of valence quarks
In this section, we will use $`\epsilon _+`$ to explain the relation of the positive mean event charge to interactions involving the proton valence quarks. The magnitude of $`\epsilon _+`$ can be related to two quantities:
* the positive excess at the quark level, given by the ratio $`\epsilon _q^+=\frac{M^+}{M^{}}`$ where $`M^{}=N_{\overline{u}}+N_d+N_s`$ and $`M^+=N_u+N_{\overline{d}}+N_{\overline{s}}`$; $`N_u`$, $`N_d`$, $`N_s`$, $`N_{\overline{u}}`$, $`N_{\overline{d}}`$ and $`N_{\overline{s}}`$ are the numbers of those events where the interaction has given a leading u-,d-,s-,$`\overline{u}`$-,$`\overline{d}`$-, respectively $`\overline{s}`$-quark.
* the probability that the charge of the initial quark and the one of the leading particle are the same. Here $`p`$ should be the same irrespective of the quark flavour. This will later be shown a flawed assumption, but we can take mean values.
Then the dependence of $`\epsilon ^+`$ upon $`\epsilon _q^+`$ and $`p`$ is described by the formula:
$$\epsilon ^+=(pM^+(1p)M^{})/(pM^{}+(1p)M^+)=(1p+\epsilon _q^+p)/(p+\epsilon _q^+(1p)).$$
The probability $`p`$, averaged on different flavours, is given by QCD based models, and is found to be of the order of 0.6 to 0.7.
The valence quark interactions from the proton are 8 times more likely to be initiated by an $`u`$ quark than by a $`d`$ quark . Hence, for valence quarks, more $`\pi ^+`$ and $`K^+`$ than $`\pi ^{}`$ and $`K^{}`$ should be detected as leading particles.
For example, in the very high $`x`$ region, the valence quark interactions will dominate the $`ep`$ cross-section. Therefore, $`\epsilon _q^+`$ is equal to 8, and $`\epsilon _+`$ is expected to be 1.39 for $`p=0.6`$.
To quantify this, we take $`Y_\alpha ^+`$ and $`Y_\alpha ^{}`$ as the yields of positive and negative quarks struck by the photon, respectively, and produced by some process $`\alpha `$. To simplify, we could consider only three different processes $`\alpha `$: Boson Gluon fusion (BGF), scattering on a sea quark (Sea) or scattering on valence quarks (V).
Here, we do not distinguish between elastic $`eq`$ diffusion and QCD Compton reactions, where one gluon is added: $`e^++qe^++q^{}+gluon`$. It has been shown with LEPTO, that so far as the event charge is concerned, they are equivalent. This can be simply understood: in QCD compton events, the gluon is often weaker than the final quark q’, and also it has to break down into two quarks (plus gluons). The chance that one out of these two quarks has a larger energy than q’ is very limited; so the leading particle will generally stay related to the q’ quark.
The positive excess at the quark level is thus: $`\epsilon _q^+=(Y_{BGF}^++Y_{sea}^++Y_V^+)/(Y_{BGF}^{}+Y_{sea}^{}+Y_V^{})`$.
For sea quarks and quarks emitted through BGF, the number of quarks is identical to the number of their antiquarks, we can state:
$$Y_{BGF}^++Y_{sea}^+=Y_{BGF}^{}+Y_{sea}^{}=Y_0.$$
Then: $`\epsilon _q^+=(1+\frac{Y^+}{Y_0})/1+\frac{Y_V^{}}{Y_0})`$. Let $`R=\frac{Y_V^+}{Y_V^{}}`$ and $`S=\frac{Y_V^+}{Y_0}`$, so that $`\frac{Y_V^{}}{Y_0}=\frac{S}{R}`$. As was quoted above, the value of R should be 8 (naively), so that : $`\epsilon _q^+=(1+S)/(1+\frac{S}{R})`$ is nearly equal to $`1+S`$. We see then that the positive excess at the quark level is a measure of the ratio of yield by valence quarks over yield by other processes.
We access to it by measuring the mean event charge, but the probability $`p`$ introduced above should be known, using QCD based models. We can perform some tests on these QCD based models, mainly on various fragmentation properties and leading particle studies (mainly $`K_0^S`$ and $`\mathrm{\Lambda }`$ ). After determining the mean event charge as a function of x and $`Q^2`$, we cannot have direct values of $`\epsilon _q^+`$ but we may compare the reconstructed charges with those predicted by QCD based models.
### 1.3 Event selection
The leading particle plays a major role in determination of the mean charge, and so we will perform our kinematic cuts on the leading particle parameters. In the following, we will cut all events having a leading particle angle less than 0.7 rad (where the angle is defined with respect to the initial proton direction ) or larger than $`\pi 0.7`$ rad. In the forward direction, we then avoid particles coming from the proton remnant. In both directions, we have to reserve some room for the next to leading (n.t.l.) and other particles, so that the event charge is well evaluated.
Another cut on all tracks has been imposed by experimental conditions at H1: to avoid beam gas background, to which the mean value of the event Charge Q is very sensitive, we require all tracks to have a transverse momentum larger than 0.6 GeV/c.
The parameter $`f=P_{leading}^{||}/P_{quark}`$ was also used to select events where the leading particle “remembers“ better the initial quark flavour. Typically, it was required to be larger than 0.15.
### 1.4 Our sample of events
There exists different QCD based models, for which we take the following versions: LEPTO(6.5), ARIADNE(410), RAPGAP(2.0), and HERWIG(5.9008). LEPTO, ARIADNE, and RAPGAP have the same hadronisation program, JETSET .
For each of these four QCD based models, we have generated 2 million events in the range $`3.<Q^2<100`$ GeV<sup>2</sup>. The corresponding luminosity is $`6.3`$ pb<sup>-1</sup>, which is below what is now available at HERA. The GRV parton density was used for all. For ARIADNE and RAPGAP, the pomeron was off. For LEPTO, the QCD effects were on.
## 2 Fragmentation studies
We have made comparisons of different QCD based models for the following fragmentation parameters:
* charge correlations in the system of the 2 or 3 leading particles;
* projected momentum of the leading particle on the hadronic direction;
* $`\phi `$, rapidity and momentum ratio correlations of the leading and the struck quark;
* $`\phi `$, rapidity and charge correlations of the leading and n.t.l. particles;
* $`K_S^0`$, $`\mathrm{\Lambda }`$ as leading particles;
* $`\rho `$, $`K^\pm `$ and $`\varphi `$ production in the system of the leading and next to leading particles.
We show here charge correlations in the system of the 2 or 3 leading particles. However, the other parameters are shown in the Internet version of this note.
All these quantities can be experimentally measured, and comparison to M.C’s provides a good test of hadronisation programs.The detailed results are shown in the extended version (on the web). Generally, there is a good agreement between different QCD based models.
### 2.1 Charge correlations for leading and n.t.l. particles
The leading and next to leading particles may either have opposite charges, which is foreseen by the naive model of quark hadronisation, or same charges, which is also foreseen but in less frequent cases. In the following, when we write a pair of signs, the first is for the leading, the second for the n.t.l.. In the next figure 4 different ratios are shown, for channels 1 to 4:
* $`\frac{+and}{++and+}`$: this is the ’negative excess’ of the n.t.l. Naively, this should be a little larger than 1, and it is very close to 1 for all QCD based models.
* $`\frac{++}{}`$: this is the positive excess for leading and n.t.l. when both have the same charge.
* $`\frac{+}{+}`$: positive excess for opposite charge leading and n.t.l, lower than the preceding one.
* $`\frac{+++}{}`$: this is the positive excess when the three first particles carry the same sign.
The agreement between the three JETSET QCD based models and on the other hand HERWIG is rather poor.
In the next figure, we show 3 other ratios:
* $`\frac{\mathrm{different}\mathrm{signs}}{\mathrm{same}\mathrm{signs}}`$ for the leading and the n.t.l..It should be much larger than 1, and it is of the order of 2.
* $`\frac{\mathrm{different}\mathrm{signs}}{\mathrm{same}\mathrm{signs}}`$ for the leading and the next to next to leading (n.t.n.t.l.). Naively, this would be expected to be less than 1 but is larger than 1. This is due to the fact that events with signs (+++) or (- - -) are very seldom, as is shown in the next channel.
* $`\frac{(++)+(++)+(+)+(+)+(+)+(++)}{3((+++)+())}`$.
We see that the value of this ratio is between 3 to 4: it is very seldom that the three first particles carry the same sign.
## 3 QCD model predictions for different parton induced quantities
As was written in the introduction, the event charge Q depends upon two quantities: the mean probability $`p`$ that the initial quark charge and that of the leading particle are the same and $`\epsilon _q^+`$: the positive excess at the quark level. Before presenting the results on the event charge, we will present the predictions of different QCD based models for these quantities.
### 3.1 The sign excess depending on the flavour
The sign excess related to a quark is defined as the ratio of the number of events originated from that quark, which gives a leading particle having its sign, to the number of those which originate from $`q`$ as well, but give a leading particle of opposite sign. For instance, for the u quark initiating the reaction, it is: $`S_u=\frac{N_+}{N_{}}`$, whereas for the $`\overline{u}`$:$`S_{\overline{u}}=\frac{N_{}}{N_+}`$. The sign excess is related directly to the probability $`p`$: for example, if we have a sign excess of 3, then the probability $`p`$ for that quark is 75 percent. Here we compare the $`S_q{}_{}{}^{}s`$ for the 4 QCD based models.
Each box of the preceding figure has 6 channels, from -3 to 3. We use the standard quark identity of the Particle Data Group :-3 is $`\overline{s}`$, -2 is $`\overline{u}`$, -1 is $`\overline{d}`$, 1 is $`d`$, 2 is $`u`$ and 3 is $`s`$. The sign excesses are shown for these 6 abscissa of the boxes.
It is easily seen that the $`d`$,$`\overline{d}`$ have less sign excess than $`u`$,$`\overline{u}`$; there are at least 2 reasons: first, in the naive model, the $`d`$ can give a proton as leading particle,which then has the wrong sign; second, for what concerns charged particles, the $`u`$ can give a $`\pi ^+`$ or a $`K^+`$, the $`d`$ only a $`\pi ^{}`$ and not a $`K^{}`$ , and the $`u`$ has there also an advantage.
It is seen that the sign excesses for a quark and its antiparticle are equal (within statistical errors), which is a good consistency check of the models.
Generally, HERWIG gives a greater sign excess than RAPGAP and LEPTO, which in turn gives more than ARIADNE. Let us point out that the sign excess reflects on one hand the hadronisation, (production of mesons and baryons), and on the other hand the production of gluons and quarks from an initial quark (the one from the proton) . As ARIADNE uses the same hadronisation program as RAPGAP and LEPTO, we see that there is a sensitivity to the underlying differences in the physics of these models.
### 3.2 The positive excess at the level of initial quarks
The positive excess at the level of quarks: $`\epsilon _q^+=(N_u+N_{\overline{d}}+N_{\overline{s}})/(N_{\overline{u}}+N_d+N_s)`$, where $`N_q`$ is the number of events where a quark $`q`$ from the proton initiates the reaction. On the next plot is shown how this quantity varies with $`x`$ and $`Q^2`$. We can see that, as expected, $`\epsilon _q^+`$ increases with $`x`$ (more valence at higher $`x`$), and that, at large $`Q^2`$, it begins to take important values.
The different QCD based models show similar behaviours, except for HERWIG for$`Q^2<10`$ GeV<sup>2</sup>.
### 3.3 Final result: mean event charge as a function of x and $`Q^2`$
To recap, the charge that we consider here is the following: $`Q=\frac{_iP_i^{||}q_i}{_iP_i^{||}}`$, where the sum over $`i`$ extends to all particles having a positive momentum projection on the struck quark direction.
The found mean event charges are shown in the next figure. HERWIG predicts a far greater charge than the other models. The charge increases with $`x`$ for the two ranges in $`Q^2`$ considered.
## 4 Conclusions
Using four QCD based models the charge properties of inclusively produced hadrons in $`ep`$ DIS processes has been investigated. Deviations of these models have been seen, particularly between those models employing string and cluster hadronisation for the total event charge. Especially, HERWIG seems to give very different results than LEPTO, ARIADNE and RAPGAP, which all three use the hadronisation program JETSET.
The mean event charge may also be measured in special cases of deep inelastic scattering:
* for charge current events, where it should have the sign of the incoming electron or positron in the beam. Due to limited statistics, this could only be an interesting check.
* for diffractive events. In the Pomeron model, the mean event charge is obviously zero. If we find a mean positive charge in these events, this means either that there is another mechanism involved in the diffraction or that the cuts selecting these diffractive events are not severe enough.
Michel Jaffré and Ursula Berthon are acknowledged for their help in processing data.
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# Untitled Document
TAUP 2550-99
22 June, 2000
Classical Radiation Reaction Off-Shell
Corrections to the Covariant Lorentz Force
O. Oron and L.P. Horwitz<sup>*</sup><sup>*</sup>Also at Department of Physics, Bar Ilan University, Ramat Gan 529000, Israel
School of Physics and Astronomy
Raymond and Beverly Sackler Faculty of Exact Sciences
Tel Aviv University, Ramat Aviv 69978, Israel
Abstract:
It has been shown by Gupta and Padmanabhan that the radiation reaction force of the Abraham-Lorentz-Dirac equation can be obtained by a coordinate transformation from the inertial frame of an accelerating charged particle to that of the laboratory. We show that the problem may be formulated in a flat space of five dimensions, with five corresponding gauge fields in the framework of the classical version of a fully gauge covariant form of the Stueckelberg-Feynman-Schwinger covariant mechanics (the zero mode fields of the $`0,1,2,3`$ components correspond to the Maxwell fields). Without additional constraints, the particles and fields are not confined to their mass shells. We show that in the mass-shell limit, the generalized Lorentz force obtained by means of the retarded Green’s functions for the five dimensional field equations provides the classical Abraham-Lorentz-Dirac radiation reaction terms (with renormalized mass and charge). We also obtain general coupled equations for the orbit and the off-shell dynamical mass during the evolution. The theory does not admit radiation if the particle remains identically on-shell. The structure of the equations implies that mass-shell deviation is bounded when the external field is removed.
PACS: 41.60.-m, 03.50.De, 03.30.+p, 11.10.Ef
Gupta and Padmanabhan<sup>4</sup> have shown that the motion of a charged particle in an electromagnetic field can be described in the inertial frame of the particle with a time varying non-trivial background metric. Using the general covariant form of the Maxwell equations and transforming back to the inertial frame of the laboratory, they obtained the Abraham-Lorentz-Dirac radiation reaction term as a consequence of this geometrical picture. This result demonstrates that the description of the motion of a charged particle in acceleration must include the radiation terms of the Abraham-Lorentz-Dirac equation<sup>5</sup>.
Alternatively, one can develop the mechanics in a flat space of higher dimension, an approach that we shall take. We shall work with the manifestly covariant mechanics of Stueckelberg<sup>6</sup>, which provides a description of dynamical systems under the influence of forces (which may be represented in terms of potentials or gauge fields) in a framework which is Lorentz covariant. This theory admits, on a classical level, deviations from the particle’s mass shell during interaction, as in quantum field theory. A similar approach was used by Mendonça and Oliveira e Silva<sup>7</sup>, who studied the motion of a relativistically kicked oscillator in the $`E,t`$ plane using what they called a “super Hamiltonian.” One can, in fact derive the relativistic Lorentz force
$$m\ddot{x}^\mu =F_\nu ^\mu \dot{x}^\nu $$
$`(1)`$
from such a Hamiltonian.
Consider the Hamiltonian<sup>5,6</sup>(we take $`c=1`$ henceforth)
$$K=\frac{(p^\mu eA^\mu (x))(p_\mu eA_\mu (x))}{2M}$$
$`(2),`$
where $`xx^\mu `$. The Hamilton equations (generalized to the four-dimensional symplectic mechanics<sup>6</sup>) are
$$\begin{array}{cc}\hfill \frac{dx^\mu }{d\tau }& =\frac{K}{p_\mu }=\frac{p^\mu eA^\mu (x)}{M}\hfill \\ \hfill \frac{dp^\mu }{d\tau }& =\frac{K}{x_\mu }=e\frac{A^\lambda (x)}{x_\mu }\frac{p^\lambda eA^\lambda (x)}{M},\hfill \end{array}$$
$`(3)`$
where $`\tau `$ is the absolute (universal) invariant time parametrizing the path of the particle in spacetime<sup>6</sup>. Computing $`\frac{dp^\mu }{d\tau }`$ from the first of these, one finds Eq. $`(1)`$. It moreover follows from the first of Eqs.$`(3)`$ that
$$\frac{dx^\mu }{d\tau }\frac{dx_\mu }{d\tau }=\frac{(p^\mu eA^\mu (x))(p_\mu eA_\mu (x))}{M^2};$$
$`(4)`$
this quantity is absolutely conserved, since $`K`$ does not depend explicitly on $`\tau `$. It follows, since the square of the proper time $`ds^2=dx^\mu dx_\mu `$, that $`ds`$ is proportional to $`d\tau `$, independently of the acceleration of the particle. The numerator of $`(4)`$ is the mass-squared of the particle; we infer that this result is associated with the restriction of the particle to a sharp mass shell.
Taking into account full $`U(1)`$ gauge invariance, the Stueckelberg-Schrödinger equation<sup>6</sup> (including a compensation field for the $`\tau `$-derivative) is
$$\left(i\frac{}{\tau }+e_0a_5\right)\psi _\tau (x)=\frac{(p^\mu e_0a^\mu (x,\tau ))(p_\mu e_0a_\mu (x,\tau ))}{2M}\psi _\tau (x),$$
$`(5)`$
where the gauge fields may depend on $`\tau `$ and $`e_0`$ is a dimensionless coupling. The corresponding classical Hamiltonian then has the form
$$K=\frac{(p^\mu e_0a^\mu (x,\tau ))(p_\mu e_0a_\mu (x,\tau ))}{2M}e_0a_5(x,\tau ).$$
$`(6)`$
The equations of motion for the field variables are given (for both the classical and quantum theories) by<sup>8</sup>
$$\lambda _\alpha f^{\beta \alpha }(x,\tau )=e_0j^\beta (x,\tau ),$$
$`(7)`$
where $`\alpha ,\beta =0,1,2,3,5`$, the last corresponding to the $`\tau `$ index, and $`\lambda `$, of dimension $`\mathrm{}^1`$, is a factor on the terms $`f^{\alpha \beta }f_{\alpha \beta }`$ in the Lagrangian associated with $`(6)`$ (including degrees of freedom of the fields), required by dimensionality, as we shall see below. The field strengths are
$$f^{\alpha \beta }=^\alpha a^\beta ^\beta a^\alpha ,$$
$`(8)`$
and the current satisfies the conservation law<sup>8,10</sup>
$$_\alpha j^\alpha (x,\tau )=0.$$
$`(9)`$
Writing out $`(9)`$ explicitly ($`j^5\rho `$, the density of events in spacetime),
$$_5\rho +_\mu j^\mu =0;$$
$`(10)`$
integrating over $`\tau `$ on $`(\mathrm{},\mathrm{})`$, and assuming that $`j^5(x,\tau )`$ vanishes<sup>8</sup> at $`|\tau |\mathrm{}`$, one finds that
$$_\mu J^\mu (x)=0,$$
where (for some dimensionless $`\eta `$)
$$J^\mu (x)=\eta _{\mathrm{}}^{\mathrm{}}𝑑\tau j^\mu (x,\tau ).$$
$`(11)`$
We identify this $`J^\mu (x)`$ with the Maxwell conserved current. In ref. 9, for example, this expression occurs with
$$j^\mu (x,\tau )=\dot{x}^\mu (\tau )\delta ^4(xx(\tau )),$$
$`(12)`$
and $`\tau `$ is identified with the proper time of the particle (an identification which can be made for the motion of a free particle).
Integrating the $`\mu `$-components of Eq. $`(7)`$ over $`\tau `$ (assuming $`f^{\mu 5}(x,\tau )0`$ for $`\tau \pm \mathrm{}`$), we obtain the Maxwell equations with the Maxwell charge $`e=e_0/\eta `$ and the Maxwell fields given by
$$A^\mu (x)=\lambda _{\mathrm{}}^{\mathrm{}}a^\mu (x,\tau )𝑑\tau .$$
$`(13)`$
The Hamiltonian of Stueckelberg<sup>6</sup> and Mondonça and Oliveira e Silva<sup>8</sup> can be recovered in the limit of the zero mode of the fields
$$a^\mu (x,\tau )=𝑑s\widehat{a}^\mu (x,s)e^{is\tau }.$$
$`(14)`$
In the zero mode limit, when the Fourier transform of the fields have support only in the neighborhood $`\mathrm{\Delta }s`$ of $`s=0`$, the vector potential takes on the form $`a^\mu (x,\tau )\mathrm{\Delta }s\widehat{a}^\mu (x,0)=(\mathrm{\Delta }s/2\pi \lambda )A^\mu (x)`$, and we identify $`e=(\mathrm{\Delta }s/2\pi \lambda )e_0`$. The zero mode therefore emerges when the inverse correlation length of the field satisfies the relation $`\eta \mathrm{\Delta }s=2\pi \lambda `$. We remark that in this limit, the fifth equation obtained from $`(7)`$ decouples; the zero mode of the $`\tau `$ derivative of $`a^\mu (x,\tau )`$ vanishes. If the parameter $`\lambda `$ is independent of the dynamical structure of the fields, then the effective width of $`\widehat{a}^\mu (x,s)`$, when it is well-defined, affects the value of the charge $`e`$, as well as the relation between the effective Maxwell current and the microscopic current $`j^\mu `$.This effect, occurring when a Maxwell type theory is a good approximation, can be understood as a classical analog of charge renormalization, where the effective charge is a function of momentum transfer.
Again, writing the Hamilton equations for the Hamiltonian $`(6)`$, we find the generalized Lorentz force<sup>10</sup>
$$M\ddot{x}^\mu =e_0f_\nu ^\mu \dot{x}^\nu +f_5^\mu =e_0\left(f_{self}^{}{}_{\nu }{}^{\mu }\dot{x}^\nu +f_{self}^{}{}_{5}{}^{\mu }+f_{ext}^{}{}_{\nu }{}^{\mu }\dot{x}^\nu +f_{ext}^{}{}_{5}{}^{\mu }\right).$$
$`(15)`$
Multiplying this equation by $`\dot{x}_\mu `$, one obtains
$$M\dot{x}_\mu \ddot{x}^\mu =e_0\dot{x}_\mu f_5^\mu =e_0\left(\dot{x}_\mu f_{self}^{}{}_{5}{}^{\mu }+\dot{x}_\mu f_{ext}^{}{}_{5}{}^{\mu }\right);$$
$`(16)`$
this equation therefore does not necessarily lead to the trivial relation between $`ds`$ and $`d\tau `$ discussed above in connection with Eq. $`(4)`$. The $`f_5^\mu `$ term has the effect of moving the particle off-shell.
In the following we use the Green’s functions for $`(7)`$ to calculate the radiation reaction force directly, as, for example, in the derivation of Sokolov and Ternov<sup>11</sup>. In the limit for which the particle stays on its mass shell during the interaction, we show that this formula reduces to the known Abraham-Lorentz-Dirac formula<sup>9,12</sup> for the Maxwell self-interaction problem. We furthermore show that the deviation from mass shell is stable. We shall use the retarded Green’s function and treat divergences by renormalization of charge and the mass parameter $`M`$.
Choosing the generalized Lorentz gauge $`_\alpha a^\alpha =0`$, Eq. $`(7)`$ becomes
$$\lambda _\alpha ^\alpha a^\beta (x,\tau )=(\sigma _\tau ^2_t^2+^2)a^\beta =e_0j^\beta (x,\tau ),$$
$`(17)`$
where $`\sigma =\pm 1`$ corresponds to the possible choices of metric for the symmetry $`\mathrm{O}(4,1)`$ or $`\mathrm{O}(3,2)`$ of the homogeneous field equations.
The Green’s functions for Eq. $`(17)`$ can be constructed from the inverse Fourier transform
$$G(x,\tau )=\frac{1}{(2\pi )^5}d^4k𝑑\kappa \frac{e^{i(k^\mu x_\mu +\sigma \kappa \tau )}}{k_\mu k^\mu +\sigma \kappa ^2}.$$
$`(18)`$
Integrating this expression over all $`\tau `$ gives the Green’s function for the standard Maxwell field. Assuming that the radiation reaction acts causally in $`\tau `$, we shall restrict our attention here to the $`\tau `$-retarded Green’s function. In his calculation of the radiation corrections to the Lorentz force, Dirac<sup>12</sup> used the difference between advanced and retarded Green’s functions in order to cancel the singularities that they contain. One can, alternatively<sup>11</sup>, use the retarded Green’s function and “renormalize” the mass in order to eliminate the singularity. In our analysis, we follow the latter procedure.
The $`\tau `$\- retarded Green’s function is given by multiplying the principal part of the integral $`(18)`$ by $`\theta (\tau )`$. Carrying out the integrations (on a complex contour in $`\kappa `$; we consider the case $`\sigma =+1`$ in the following), one finds (this Green’s function differs from that used in ref. 13, constructed on a complex contour in $`k^0`$)
$$G(x,\tau )=\frac{2\theta (\tau )}{(2\pi )^3}\{\begin{array}{cc}\frac{\mathrm{tan}^1\left(\frac{\sqrt{x^2\tau ^2}}{\tau }\right)}{(x^2\tau ^2)^{\frac{3}{2}}}\frac{\tau }{x^2(x^2+\tau ^2)}\hfill & x^2+\tau ^2<0\text{;}\hfill \\ \frac{1}{2}\frac{1}{(\tau ^2+x^2)^{\frac{3}{2}}}\mathrm{ln}\left|\frac{\tau \sqrt{\tau ^2+x^2}}{\tau +\sqrt{\tau ^2+x^2}}\right|\frac{\tau }{x^2(\tau ^2+x^2)}\hfill & x^2+\tau ^2>0\text{.}\hfill \end{array}$$
$`(19)`$
With the help of this Green’s function, the solutions of Eq. $`(17)`$ for the self-fields can be written,
$$\begin{array}{cc}\hfill a_{self}^{}{}_{}{}^{\mu }(x,\tau )& =\frac{e_0}{\lambda }d^4x^{}𝑑\tau ^{}G(xx^{},\tau \tau ^{})\dot{x}^\mu (\tau ^{})\delta ^4(x^{}x(\tau ^{}))\hfill \\ & =\frac{e_0}{\lambda }𝑑\tau ^{}\dot{x}^\mu (\tau ^{})G(xx(\tau ^{}),\tau \tau ^{})\hfill \\ \hfill a_{self}^{}{}_{}{}^{5}(x,\tau )& =\frac{e_0}{\lambda }d^4x^{}𝑑\tau ^{}G(xx^{},\tau \tau ^{})\delta ^4(x^{}x(\tau ^{}))\hfill \\ & =\frac{e_0}{\lambda }𝑑\tau ^{}G(xx(\tau ^{}),\tau \tau ^{})\hfill \end{array}$$
$`(20)`$
where we have used (12) (along with $`j^5(x,\tau )=\delta ^4(xx(\tau )))`$. We have written this Green’s function as a scalar, acting in the same way on all five components of the source $`j^\alpha `$; to assure that the resulting field is in Lorentz gauge, however, it should be written as a five by five matrix, with the factor $`\delta _\beta ^\alpha k^\alpha k_\beta /k^2`$ ($`k_5=\kappa `$) included in the integrand. Since we are computing only the gauge invariant field strengths here, this extra term will not influence any of the results.
From $`(8)`$ and $`(15)`$, it then follows that the generalized Lorentz force for the self-action (the force of the fields generated by the world line on a point $`x^\mu (\tau )`$ of the trajectory) is
$$\begin{array}{cc}\hfill M\ddot{x}^\mu & =\frac{e_0^2}{\lambda }𝑑\tau ^{}(\dot{x}^\nu (\tau )\dot{x}_\nu (\tau ^{})^\mu \dot{x}^\nu (\tau )\dot{x}^\mu (\tau ^{})_\nu )G(xx(\tau ^{}))|_{x=x(\tau )}\hfill \\ & +\frac{e_0^2}{\lambda }𝑑\tau ^{}(^\mu \dot{x}^\mu (\tau ^{})_\tau )G(xx(\tau ^{}))|_{x=x(\tau )}\hfill \\ & +e_0\left(f_{ext}^{}{}_{\nu }{}^{\mu }\dot{x}^\nu +f_{ext}^{}{}_{5}{}^{\mu }\right)\hfill \end{array}$$
$`(21)`$
We define $`u(x_\mu (\tau )x_\mu (\tau ^{}))(x^\mu (\tau )x^\mu (\tau ^{}))`$, so that
$$_\mu =2(x_\mu (\tau )x_\mu (\tau ^{}))\frac{}{u}.$$
$`(22)`$
Eq. (21) then becomes
$$\begin{array}{cc}\hfill M\ddot{x}^\mu & =2\frac{e_0^2}{\lambda }d\tau ^{}\{\dot{x}^\nu (\tau )\dot{x}_\nu (\tau ^{})(x^\mu (\tau )x^\mu (\tau ^{}))\hfill \\ & \dot{x}^\nu (\tau )\dot{x}^\mu (\tau ^{})(x_\nu (\tau )x_\nu (\tau ^{}))\}\frac{}{u}G(xx(\tau ^{}),\tau \tau ^{})|_{x=x(\tau )})\hfill \\ & +\frac{e_0^2}{\lambda }d\tau ^{}\{2(x^\mu (\tau )x^\mu (\tau ^{}))\frac{}{u}\dot{x}^\mu (\tau ^{})_\tau \}G(xx(\tau ^{}),\tau \tau ^{})|_{x=x(\tau )}).\hfill \\ & +e_0\left(f_{ext}^{}{}_{\nu }{}^{\mu }\dot{x}^\nu +f_{ext}^{}{}_{5}{}^{\mu }\right)\hfill \end{array}$$
$`(23)`$
We now expand the integrands in Taylor series around the most singular point $`\tau =\tau ^{}`$. In this neighborhood, keeping the lowest order terms in $`\tau ^{\prime \prime }=\tau \tau ^{}`$, the variable $`u`$ reduces to $`u\dot{x}^\mu \dot{x}_\mu \tau _{}^{\prime \prime }{}_{}{}^{2}`$. We shall also use the following definition;
$$\epsilon 1+\dot{x}^\mu \dot{x}_\mu ,$$
$`(24)`$
a quantity that vanishes on the mass shell of the particle (as we have pointed out above). In this case the derivatives of $`(19)`$ take the form
$$\begin{array}{cc}\hfill \frac{G}{u}& \frac{\theta (\tau ^{\prime \prime })f_1(ϵ)}{(2\pi )^3\tau _{}^{\prime \prime }{}_{}{}^{5}}\hfill \\ \hfill \frac{G}{\tau ^{\prime \prime }}& \frac{\theta (\tau ^{\prime \prime })f_2(ϵ)}{(2\pi )^3\tau _{}^{\prime \prime }{}_{}{}^{4}}+\frac{\delta (\tau ^{\prime \prime })f_3(ϵ)}{(2\pi )^3\tau _{}^{\prime \prime }{}_{}{}^{3}}\hfill \end{array}$$
$`(25)`$
where we have used the following definitions:
$`\epsilon <0`$:
$$\begin{array}{cc}\hfill f_1(\epsilon )& =\frac{3\mathrm{tan}^1(\sqrt{\epsilon })}{(\epsilon )^{\frac{5}{2}}}\frac{3}{\epsilon ^2(1\epsilon )}+\frac{2}{\epsilon (1\epsilon )^2}\hfill \\ \hfill f_2(\epsilon )& =\frac{3\mathrm{tan}^1(\sqrt{\epsilon })}{(\epsilon )^{\frac{5}{2}}}\frac{1}{\epsilon ^2}\frac{2\epsilon }{\epsilon ^2(1\epsilon )}\hfill \\ \hfill f_3(\epsilon )& =\frac{\mathrm{tan}^1(\sqrt{\epsilon })}{(\epsilon )^{\frac{3}{2}}}+\frac{1}{\epsilon (1\epsilon )}\hfill \end{array}$$
$`(26a)`$
$`\epsilon >0`$:
$$\begin{array}{cc}\hfill f_1(\epsilon )& =\frac{\frac{3}{2}\mathrm{ln}\left|\frac{1+\sqrt{\epsilon }}{1\sqrt{\epsilon }}\right|}{(\epsilon )^{\frac{5}{2}}}\frac{3}{\epsilon ^2(1\epsilon )}+\frac{2}{\epsilon (1\epsilon )^2}\hfill \\ \hfill f_2(\epsilon )& =\frac{\frac{3}{2}\mathrm{ln}\left|\frac{1+\sqrt{\epsilon }}{1\sqrt{\epsilon }}\right|}{(\epsilon )^{\frac{5}{2}}}\frac{1}{\epsilon ^2}\frac{2\epsilon }{\epsilon ^2(1\epsilon )}\hfill \\ \hfill f_3(\epsilon )& =\frac{\frac{1}{2}\mathrm{ln}\left|\frac{1+\sqrt{\epsilon }}{1\sqrt{\epsilon }}\right|}{(\epsilon )^{\frac{3}{2}}}+\frac{1}{\epsilon (1\epsilon )}\hfill \end{array}$$
$`(26b)`$
For either sign of $`\epsilon `$, when $`\epsilon 0`$,
$$\begin{array}{cc}\hfill f_1(\epsilon )& \frac{8}{5}+\frac{24}{7}\epsilon ,\hfill \\ \hfill f_2(\epsilon )& \frac{2}{5}\frac{4}{7}\epsilon ,\hfill \\ \hfill f_3(\epsilon )& \frac{2}{3}+\frac{4}{5}\epsilon \hfill \end{array}$$
$`(26c)`$
One sees that the derivatives in (25) have no singularity in $`\epsilon `$ at $`\epsilon =0`$.
From $`(8)`$ and $`(20)`$, we have
$$\begin{array}{cc}& f_{self}^{}{}_{5}{}^{\mu }(x(\tau ),\tau )=\hfill \\ & e𝑑\tau ^{}\{2(x^\mu (\tau )x^\mu (\tau ^{}))\frac{}{u}\dot{x}^\mu (\tau ^{})_\tau \}G(xx(\tau ^{}),\tau \tau ^{})|_{x=x(\tau )},\hfill \end{array}$$
$`(27)`$
We see (from $`(25)`$) that the main contributions to the integrals come from small $`\tau ^{\prime \prime }`$. We may therefore expand $`x^\mu (\tau )x^\mu (\tau ^{})`$ and $`\dot{x}^\mu (\tau )\dot{x}^\mu (\tau ^{})`$ in $`(27)`$ in power series in $`\tau ^{\prime \prime }`$, and write the integrals formally with infinite limits.
Substituting $`(27)`$ into $`(16)`$, we obtain (note that $`x^\mu `$ and its derivatives are evaluated at the point $`\tau `$, and are not subject to the $`\tau ^{\prime \prime }`$ integration), after integrating by parts using $`\delta (\tau ^{\prime \prime })=\frac{}{\tau ^{\prime \prime }}\theta (\tau ^{\prime \prime })`$,
$$\begin{array}{cc}\hfill M\dot{x}_\nu \ddot{x}^\nu =\frac{2e_0^2}{\lambda (2\pi )^3}_{\mathrm{}}^{\mathrm{}}d\tau ^{\prime \prime }\{\frac{(f_1f_23f_3)}{\tau ^{\prime \prime 4}}\dot{x}_\nu \dot{x}^\nu & \frac{(\frac{1}{2}f_1f_22f_3)}{\tau ^{\prime \prime 3}}\dot{x}_\nu \ddot{x}^\nu \hfill \\ \hfill +\frac{(\frac{1}{6}f_1\frac{1}{2}f_2\frac{1}{2}f_3)}{\tau ^{\prime \prime 2}}\dot{x}_\nu \stackrel{\mathrm{}}{𝑥}^\nu \}\theta (\tau ^{\prime \prime })+e_0\dot{x}_\mu f_{ext}^\mu _5.& \end{array}$$
$`(28)`$
The integrals are divergent at the lower bound $`\tau ^{\prime \prime }=0`$ imposed by the $`\theta `$-function; we therefore take these integrals to a cut-off $`\mu >0`$. Eq.$`(28)`$ then becomes
$$\begin{array}{cc}\hfill \frac{M}{2}\dot{\epsilon }=\frac{2e_0^2}{\lambda (2\pi )^3}\{\frac{(f_1f_23f_3)}{3\mu ^3}(\epsilon 1)& \frac{(\frac{1}{2}f_1f_22f_3)}{4\mu ^2}\dot{\epsilon }\hfill \\ \hfill +\frac{(\frac{1}{6}f_1\frac{1}{2}f_2\frac{1}{2}f_3)}{\mu }\dot{x}_\nu \stackrel{\mathrm{}}{𝑥}^\nu \}+e_0\dot{x}_\mu f_{ext}^\mu _5.& \end{array}$$
$`(29)`$
Following a similar procedure, we obtain from $`(23)`$
$$\begin{array}{cc}\hfill M\ddot{x}^\mu & =\frac{2e_0^2}{\lambda (2\pi )^3}\{\frac{f_1}{4\mu ^2}((1\epsilon )\ddot{x}^\mu +\frac{1}{2}\dot{\epsilon }\dot{x}^\mu )+\frac{f_1}{3\mu }(\dot{x}_\nu \stackrel{\mathrm{}}{𝑥}^\nu \dot{x}^\mu +(1\epsilon )\stackrel{\mathrm{}}{𝑥}^\mu )\hfill \\ & +\frac{(f_1f_23f_3)}{3\mu ^3}\dot{x}^\mu \frac{(\frac{1}{2}f_1f_22f_3)}{2\mu ^2}\ddot{x}^\mu +\frac{(\frac{1}{6}f_1\frac{1}{2}f_2\frac{1}{2}f_3)}{\mu }\stackrel{\mathrm{}}{𝑥}^\mu \}\hfill \\ & +e_0\left(f_{ext}^{}{}_{\nu }{}^{\mu }\dot{x}^\nu +f_{ext}^{}{}_{5}{}^{\mu }\right).\hfill \end{array}$$
$`(30)`$
Using $`(29)`$ to substitute for the coefficient of the $`\frac{1}{\mu ^3}`$ term in $`(30)`$ , we obtain (for $`\epsilon 1`$)
$$\begin{array}{cc}\hfill M(\epsilon )\ddot{x}^\mu & =\frac{1}{2}\frac{M(\epsilon )}{(1\epsilon )}\dot{\epsilon }\dot{x}^\mu +\frac{2e_0^2}{\lambda (2\pi )^3\mu }F(\epsilon )\left\{\stackrel{\mathrm{}}{𝑥}^\mu +\frac{1}{(1\epsilon )}\dot{x}_\nu \stackrel{\mathrm{}}{𝑥}^\nu \dot{x}^\mu \right\}\hfill \\ & +\frac{e_0\dot{x}^\mu \dot{x}_\nu f_{ext}^{}{}_{5}{}^{\nu }}{1\epsilon }+e_0f_{ext}^{}{}_{\nu }{}^{\mu }\dot{x}^\nu +e_0f_{ext}^{}{}_{5}{}^{\mu },\hfill \end{array}$$
$`(31)`$
where
$$F(\epsilon )=\frac{f_1}{3}(1\epsilon )+(\frac{1}{6}f_1\frac{1}{2}f_2\frac{1}{2}f_3).$$
$`(32)`$
Here, the coefficients of $`\ddot{x}^\mu `$ have been grouped into a renormalized (off-shell) mass term, defined (as in the procedure of Sokolov and Ternov<sup>11</sup>) as
$$M(\epsilon )=M+\frac{e^2}{2\mu }\left[\frac{f_1(1\epsilon )}{2}+\frac{1}{2}f_1f_22f_3\right]$$
$`(33)`$
where, as we shall see below,
$$e^2=\frac{2e_0^2}{\lambda (2\pi )^3\mu },$$
$`(34)`$
can be identified with the Maxwell charge by studying the on-shell limit.
We now obtain, from $`(31)`$,
$$\begin{array}{cc}\hfill M(\epsilon )\ddot{x}^\mu & =\frac{1}{2}\frac{M(\epsilon )}{1\epsilon }\dot{\epsilon }\dot{x}^\mu +F(\epsilon )e^2\left\{\stackrel{\mathrm{}}{𝑥}^\mu +\frac{1}{1\epsilon }\dot{x}_\nu \stackrel{\mathrm{}}{𝑥}^\nu \dot{x}^\mu \right\}\hfill \\ & +e_0f_{ext}^{}{}_{\nu }{}^{\mu }\dot{x}^\nu +e_0\left(\frac{\dot{x}^\mu \dot{x}_\nu }{1\epsilon }+\delta _\nu ^\mu \right)f_{ext}^{}{}_{5}{}^{\nu }.\hfill \end{array}$$
$`(35)`$
We remark that when one multiplies this equation by $`\dot{x}_\mu `$, it becomes an identity (all of the terms except for $`e_0f_{ext}^{}{}_{\nu }{}^{\mu }\dot{x}^\nu `$ may be grouped to be proportional to $`\left(\frac{\dot{x}^\mu \dot{x}_\nu }{1\epsilon }+\delta _\nu ^\mu \right)`$); one must use Eq. $`(29)`$ to compute the off-shell mass shift $`\epsilon `$ corresponding to the longitudinal degree of freedom in the direction of the four velocity of the particle. Eq. $`(35)`$ determines the motion orthogonal to the four velocity. Equations $`(29)`$ and $`(35)`$ are the fundamental dynamical equations governing the off-shell orbit.
We now show that the standard relativistic Lorentz force, with radiation corrections, can be obtained from these equations when $`\mu \dot{\epsilon }<<\epsilon <<1`$ and $`\ddot{\epsilon }`$ and $`f_{ext}^\mu _5`$ are small. In this case, Eq. $`(29)`$ becomes
$$\left(M\frac{1}{15\mu }\right)\frac{\dot{\epsilon }}{2}e^2\left\{\frac{8\epsilon }{15\mu ^2}+\frac{2}{15}\dot{x}_\nu \stackrel{\mathrm{}}{𝑥}^\nu \right\}$$
$`(36)`$
The left hand side can be neglected if
$$\left[M/(\frac{e^2}{\mu })\right](\mu \dot{\epsilon })<<\epsilon .$$
$`(37)`$
We shall see below that we must have $`0.68e^2/\mu <M`$ for stability of $`\epsilon `$, but if $`e^2/\mu `$ is not too small, the inequality $`(37)`$ is consistent with our assumed inequalities, and it then follows that
$$4\epsilon /\mu ^2\dot{x}_\nu \stackrel{\mathrm{}}{𝑥}^\nu .$$
$`(38)`$
If, furthermore, $`\ddot{\epsilon }`$ is small, then
$$\dot{x}_\mu \stackrel{\mathrm{}}{𝑥}^\mu =\ddot{\epsilon }\ddot{x}_\mu \ddot{x}^\mu \ddot{x}_\mu \ddot{x}^\mu ,$$
$`(39)`$
the known expression associated with radiation. Since $`\epsilon /\mu ^2`$ may be appreciable even if $`\epsilon `$ is small, the inequalitites we have assumed can admit a significant contribution of this type. Under these conditions equation $`(34)`$ becomes,
$$M_{ren}\ddot{x}^\mu =\frac{2}{3}e^2\{\stackrel{\mathrm{}}{𝑥}^\mu \ddot{x}_\nu \ddot{x}^\nu \dot{x}^\mu \}+e_0f_{ext}^{}{}_{\nu }{}^{\mu }\dot{x}^\nu ,$$
$`(40)`$
where $`M_{ren}=M(\epsilon )|_{\epsilon =0}=M+e^2/3\mu `$.
This result is of the form of the standard relativistic Lorentz force with radiation reaction.<sup>9,11,12,14</sup>
We now study the stability of the variations of the off-shell parameter $`\epsilon `$ when the external field is removed. First, we construct an equation of motion for $`\epsilon `$. We define the functions
$$\begin{array}{cc}\hfill F_1(\epsilon )& =\frac{1}{3\mu ^2}(\epsilon 1)(f_1f_23f_3)\hfill \\ \hfill F_2(\epsilon )& =\frac{1}{4\mu }(\frac{1}{2}f_1f_22f_3)\hfill \\ \hfill F_3(\epsilon )& =\frac{1}{6}f_1\frac{1}{2}f_2\frac{1}{2}f_3\hfill \end{array}$$
$`(41)`$
so equation $`(29)`$, in the absence of external fields, becomes:
$$\frac{M}{2}\dot{\epsilon }=e^2\left\{F_1(\epsilon )+F_2(\epsilon )\dot{\epsilon }+F_3(\epsilon )\dot{x}_\mu \stackrel{\mathrm{}}{𝑥}^\mu \right\}.$$
$`(42)`$
Solving for the explicit $`x`$ derivatives in $`(42)`$ and differentiating with respect to $`\tau `$, one obtains
$$\begin{array}{cc}\hfill \dot{x}_\mu \stackrel{\mathrm{}.}{𝑥}^\mu +\ddot{x}_\mu \stackrel{\mathrm{}}{𝑥}^\mu =& \frac{1}{F_3}\left\{F_2^{}\dot{\epsilon }^2+\ddot{\epsilon }\left(\frac{M}{2e^2}+F_2\right)F_1^{}\dot{\epsilon }\right\}\hfill \\ & \frac{F_3^{}}{F_{3}^{}{}_{}{}^{2}}\left\{F_2+\frac{M}{2e^2}\dot{\epsilon }F_1\right\}\dot{\epsilon }H.\hfill \end{array}$$
$`(43)`$
Together with
$$\dot{x}_\mu \stackrel{\mathrm{}.}{𝑥}^\mu +3\ddot{x}_\mu \stackrel{\mathrm{}}{𝑥}^\mu =\frac{1}{2}\stackrel{\mathrm{}}{𝜀}$$
one finds, from $`(43)`$,
$$\ddot{x}_\mu \stackrel{\mathrm{}}{𝑥}^\mu =\frac{1}{4}\stackrel{\mathrm{}}{𝜀}\frac{1}{2}H(\epsilon ,\dot{\epsilon },\ddot{\epsilon })$$
$`(44)`$
Multiplying Eq.$`(35)`$ by $`\ddot{x}_\mu `$ (with no external fields) and using $`(42)`$ and $`(44)`$, we obtain
$$\stackrel{\mathrm{}}{𝜀}A(\epsilon )\ddot{\epsilon }+B(\epsilon )\dot{\epsilon }^2+C(\epsilon )\dot{\epsilon }D(\epsilon )=0,$$
$`(45)`$
where
$$\begin{array}{cc}\hfill A(\epsilon )& =\frac{2}{F_3}\left(\frac{M}{2e^2}+F_2\right)+\frac{2M(\epsilon )}{e^2F(\epsilon )},\hfill \\ \hfill B& =\frac{2F_3^{}}{F_3^2}\left(\frac{M}{2e^2}+F_2\right)\frac{2F_2^{}}{F_3}+\frac{2}{1\epsilon }\frac{1}{F_3}\left(\frac{M}{2e^2}+F_2\right)\frac{M(\epsilon )}{e^2F(\epsilon )}\frac{1}{1\epsilon },\hfill \\ \hfill C& =\frac{4M(\epsilon )}{e^2F(\epsilon )}\frac{1}{F_3}\left(\frac{M}{2e^2}+F_2\right)\frac{2}{F_3^2}F_1F_3^{}\hfill \\ & \frac{2F_1}{(1\epsilon )F_3}+\frac{2}{F_3}F_1^{},\hfill \\ \hfill D& =\frac{4M(\epsilon )}{e^2F(\epsilon )}\frac{F_1}{F_3}.\hfill \end{array}$$
$`(46)`$
We first study the possibilty of having a solution of the form $`\epsilon \epsilon _0`$, a constant. In this case $`\ddot{\epsilon }=0`$ implies,
$$\dot{x}^\mu \stackrel{\mathrm{}}{𝑥}_\mu =\ddot{x}_\mu \ddot{x}^\mu .$$
Since all the derivatives of $`\epsilon `$ are zero we also find from $`(44)`$,
$$\ddot{x}_\mu \stackrel{\mathrm{}}{𝑥}^\mu =0$$
. Multiplying eq.$`(35)`$ by $`\ddot{x}_\mu `$ and substituting these last two results we get
$$\ddot{x}_\mu \ddot{x}^\mu =\dot{x}^\mu \stackrel{\mathrm{}}{𝑥}_\mu =0$$
From $`(42)`$ we find then that
$$\dot{x}_\mu \stackrel{\mathrm{}}{𝑥}^\mu =\frac{F_1}{F_3}=0.$$
From $`(41)`$ and $`(26b)`$, one sees that this equation can be satisfied only if $`\epsilon =0`$ ($`F_1=0`$) or $`\epsilon =1`$ ($`F_3=\mathrm{}`$).
Since $`\dot{\epsilon }=0`$ we find that $`\dot{t}^2\ddot{t}^2=|\dot{𝐱}|^2|\ddot{𝐱}|^2cos^2\theta `$. Together with $`\ddot{x}_\mu \ddot{x}^\mu =\ddot{t}^2|\ddot{𝐱}^2|=0`$ this implies
$$\dot{t}^2\ddot{t}^2=|\dot{x}|^2\ddot{t}^2\mathrm{cos}^2\theta $$
The solution $`\ddot{t}=|\ddot{x}|=0`$ implies $`\dot{x}_\mu =\mathrm{const}`$ . The other solution $`\dot{t}^2=|\dot{x}|^2\mathrm{cos}^2\theta `$ implies that $`|\dot{𝐱}|^2(1\mathrm{cos}^2\theta )=\epsilon 1`$; since the left hand side is positive, $`\epsilon `$ cannot be zero, and the only possibility for a constant solution is then $`\epsilon =1`$, motion on the light cone. We shall show below that the trajectory cannot reach this bounday.
The mass shell condition $`\epsilon =0`$, in the theoretical framework we have given here, implies that the particle motion must be with constant velocity, and that no radiation ($`\ddot{x}_\mu \ddot{x}^\mu =0`$) is possible, i.e., in order to radiate, the particle must be off-shell. This result is also true in the presence of an external field. In particular, it follows from Eq. $`(29)`$ that for $`\epsilon 0`$,
$$\begin{array}{cc}\hfill \frac{2}{15\mu }\dot{x}^\mu \stackrel{\mathrm{}}{𝑥}_\mu & \frac{2}{15\mu }\ddot{x}_\mu \ddot{x}^\mu \hfill \\ & =e_0\dot{x}_\mu f_{ext}^{}{}_{5}{}^{\mu }.\hfill \end{array}$$
$`(47)`$
From Eq.$`(15)`$, however, it follows (in case $`\dot{\epsilon }=0`$) that $`\dot{x}_\mu f_{ext}^{}{}_{5}{}^{\mu }=\dot{x}_\mu f_{self}^{}{}_{5}{}^{\mu }`$, so that the nonvanishing value of $`\ddot{x}_\mu \ddot{x}^\mu `$ corresponds only to a self-acting field $`f_{self}^{}{}_{5}{}^{\mu }`$ (driven by $`f_{ext}^{}{}_{5}{}^{\mu }`$), and not to radiation.
We now show that, in general, $`\epsilon `$ is bounded when the external fields are turned off. For the case $`\epsilon <0`$ the function $`F_3`$ is zero at $`\epsilon =0.735`$. In this case eq.$`(42)`$ becomes
$$\dot{\epsilon }(0.735)=\frac{F_1(0.735)}{\frac{M}{2e^2}+F_2(0.735)}.$$
$`(48)`$
If $`\dot{\epsilon }>0`$ at this value of $`\epsilon `$, then $`\epsilon `$ cannot cross this boundary. Since $`F_1(0.735)=\frac{.624}{\mu ^2}`$ , $`F_2(0.375)=\frac{.259}{\mu }`$, this condition implies that
$$\mu >0.68\frac{e^2}{M}.$$
Setting $`M,e`$ equal to the electron mass (the lowest mass charged particle) and charge one finds that $`\mu >10^{23}`$sec, a cut-off of reasonable size for a classical theory.
We now show that $`ϵ`$ is bounded from above by unity. The full classical Hamiltonian, obtained by adding the contribution of the fields to the expression on the right hand side of $`(6)`$, is a conserved quantity. In the absence of external fields, all the field quantities are related to the source particle through the Green’s functions. In the absence of external fields, as the particle motion approaches the light cone, there are infinite contributions arising from the fields evaluated on the particle trajectory. In this case, it follows from $`(4)`$ that $`\frac{(p^\mu e_0a^\mu (x,\tau ))(p_\mu e_0a_\mu (x,\tau ))}{2M}=0`$. The $`a_5`$ self-field term is less singular than the $`f_{\mu \nu }f^{\mu \nu }`$ and $`f_{\mu 5}f_5^\mu `$ terms, which involve derivatives of the Green’s functions, as in $`(25)`$, squared. As seen from $`(26b)`$, the most singular contribution arises from $`f_1^2`$. Since the total Hamiltonian $`K`$ is conserved, the coefficient of this singularity must vanish. The coefficients involve just $`\dot{\epsilon }`$ (and its square) and $`\ddot{\epsilon }`$; one finds a simple nonlinear differential equation for which only $`\dot{\epsilon }=0`$ can be a solution. It follows that the conservation law restricts the evolution of $`\epsilon `$ to values less than unity, i.e. the particle trajectory cannot pass through the light cone.
This bound manifests itself in the structure of the differential equation $`(45)`$ for $`\epsilon `$. In the limit that $`\epsilon 1`$, the coefficients $`A,B,C,D`$ are all finite; however the behavior of the linearized solution depends on the derivates of these coefficients, and, in this limit, $`B^{}`$ is singular, driving the solution away from the light cone.
Numerical studies are under way to follow the motion of this highly nonlinear system both in the presence and absence of external fields.
Acknowledgements: We wish to thank Y. Ashkenazy, J. Bekenstein, C. Piron and F. Rohrlich and S.L. Adler for helpful discussions, and Z. Schuss and J. Schiff for explaining some things about non-linear equations to us. One of us (L.H.) also wishes to thank S.L. Adler for his hospitality at the Institute for Advanced Study where much of this work was done.
References
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5. Y. Ashkenazy and L.P. Horwitz, chao-dyn/9905013, submitted.
6. E.C.G. Stueckelberg, Helv. Phys. Acta 14, 322 (1941); 14, 588 (1941); JR.P. Feynman, Rev. Mod. Phys. 20, 367 (1948); R.P. Feynman, Phys. Rev. 80, 440(1950); J.S. Schwinger, Phys. Rev. 82, 664 (1951);L.P. Horwitz and C. Piron, Helv. Phys. Acta 46, 316 (1973).
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8. D. Saad, L.P. Horwitz and R.I. Arshansky, Found. of Phys. 19, 1125 (1989); M.C. Land, N. Shnerb and L.P. Horwitz, Jour. Math. Phys. 36, 3263 (1995); N. Shnerb and L.P. Horwitz, Phys. Rev A48, 4068 (1993). We use a different convention for the parameters here.
9. See, for example, J.D. Jackson, Classical Electrodynamics, 2nd edition, John Wiley and Sons, New York(1975); F. Rohrlich, Classical Charged Particles, Addison-Wesley, Reading, (1965); S. Weinberg, Gravitation and Cosmology: Principles and Applications of the General Theory of Relativity, Wiley, N.Y. (1972).
10. M.C. Land and L.P. Horwitz, Found. Phys. Lett. 4, 61 (1991); M.C. Land, N. Shnerb and L.P. Horwitz, Jour. Math. Phys. 36, 3263 (1995).
11. For example, A.A. Sokolov and I.M. Ternov, Radiation from Relativistic Electrons, Amer. Inst. of Phys. Translation Series, New York (1986).
12. P.A.M. Dirac, Proc. Roy. Soc. London Ser. A, 167, 148(1938).
13. M.C. Land and L.P. Horwitz, Found. Phys. 21, 299 (1991).
14. L.D. Landau and E.M. Lifshitz,The Classical Theory of Fields 4th ed., (Pergamon Pr., Oxford, 1975).
15. A.O. Barut and Nuri Unal, Phys. Rev A40, 5404 (1989) found non-vanishing contributions of the type $`\dot{x}_\nu \ddot{x}^\nu `$ to the Lorentz-Dirac equation in the presence of spin.
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# Signatures of Quark-Gluon Plasma Phase Transition in High-Energy Nuclear Collisions
## 1 INTRODUCTION
A recent news release from CERN stated that a new state of matter, the quark-gluon plasma (QGP), was created by high-energy heavy-ion collisions at CERN. The evidence for such an observation was summarized by Heinz and Jacobs . The news release has generated a great deal of excitement. It is useful to review and re-examine this evidence from various perspectives to see whether such an excitement is justified or not.
In this review and re-examination, we shall discuss the nature of the order of the hadron-QGP phase transition, the dynamics of how such a transition takes place in high-energy heavy-ion collisions, and the signatures for the phase transition. Prominent signatures include the large pressure in the quark-gluon plasma phase which leads to an anomalous increase in the freeze-out volume, the emission of direct photons when the matter is in the quark-gluon plasma phase , the strong enhancement of strangeness , and the suppression of $`J/\psi `$ production . Other evidence can also be found by mapping out the phase boundary in the phase diagram from the distribution and the yield of hadron products, which will be reported by Braun-Munzinger in these Proceedings . Our re-assessment of these signatures provides additional support to indicate that, in line with the summary of Heinz and Jacobs , the experimental data appear to be consistent with the occurrence of the quark-gluon plasma.
## 2 THE QUARK-GLUON PLASMA PHASE TRANSITION
A phase transition is classified as a first- or second-order transition depending on how its free energy varies with the order parameter. The phase transition is first order if the free energy as a function of the order parameter exhibits two local minima at different values of the order parameter, and the order parameter of the lowest minimum changes from one local minimum to the other discontinuously as the temperature passes the critical temperature $`T_{\mathrm{crit}}`$ as shown schematically in Fig. 1$`a`$. A phase transition is second order if the free energy as a function of the order parameter exhibits a single minimum, and the location for the minimum changes continuously for different values of the order parameter as the temperature passes the critical temperature $`T_{\mathrm{crit}}`$ (Fig. 1$`b`$). (See also Fig. 11.3 and 11.4 of Ref. .)
Fig. 1 The free energy divided by the temperature as a function of the order parameter for different temperature $`T`$. Fig. 1$`a`$ shows a first-order phase transition, and Fig. 1$`b`$ shows a second-order phase transition. An arrow indicates the location of a local minimum.
The order parameter for stronglyinteracting matter is the Wilson line parameter, which is related to the interaction energy $`V`$ when an isolated quark is placed in the medium. It is proportional to $`e^{Vt}`$ for a Wilson line of temporal length $`t`$. In the hadron phase, because of the linear interaction between a quark and an antiquark, the interaction energy $`V`$ is infinite when an isolated quark is placed in themedium and the Wilson line parameter is zero. The Wilson line parameter is non-zero in the deconfined phase.
Lattice gauge calculations show that the transition is first order for a quark-gluon plasma with two flavors, butchanges to a second-order phase transition for three flavors with the physical mass of the strange quark .These theoretical predictions need to confront experimental data to deduce the nature of the phase transition.
How does one visualize the difference between the hadron matter and the quark-gluon plasma? One can envisage a spatial distribution of the link variables which contain the gluon degrees of freedom and the generators of SU(3). The link variables in strongly interacting matter are the analogue of spins in a spin gas. The hadron phase is the low-temperature phase in which neighboring link variables (in a plaquette) are all correlated in order to reside in the lowest energy state. The correlation can be considered roughly as some generalized ‘alignment’ of the link variables, in analogy with the alignment of spins in the spin lattice gas. In such an ‘aligned’ state of the link variables in all regions of space, the interaction between a quark and an antiquark at large distances becomes a linear function of the distance, which leads to the confinement of quarks and gluons. In contrast, in the high-temperature quark-gluon plasma phase, the temperature is so high that the tendency to align the link variables due to the QCD interaction is overwhelmed by the tendency to disalign the link variables due to the thermal fluctuation. The link variables are no longer aligned. The interaction between a quark and an antiquark above the critical temperature at large distances is no longer governed by the linear interaction, and the quarks and gluons become deconfined.
The deconfinement phase is further accompanied by the restoration of chiral symmetry in which the light quarks in the hadron phase are restored to their nearly massless current quark mass values. Thus, the phase transition from hadron matter to the quark-gluon plasma can be described as a change of the link variables from an ordered aligned state to a state of disorientation, accompanied by an abundant production of nearly massless quark pairs and additional gluons. As the link variables representing the gluons are now free to orient randomly and the nearly massless light quark pairs are copiously produced, the degree of freedom is greatly enlarged and the energy density of the deconfined phase is much greater as a result.
## 3 HEAVY-ION COLLISIONS TO PRODUCE THE PHASE TRANSITION
High-energy heavy-ion collisions can be used to produce a phase transition from the hadron matter to the quark-gluon plasma. Because of the Lorentz contraction, the nucleon-nucleon collisions in a nucleus-nucleus collision occur at about the same time and at nearly the same spatial proximity. As a consequence, a region of very high energy density can be created.
Fig. 2. The space-time diagram of a row of $`N_B`$ projectile nucleons colliding with a row of $`N_A`$ target nucleons.
The dynamics of such a nucleus-nucleus collision can be best described in the nucleon-nucleon center-of-mass system. For a given impact parameter, we can divide thetransverse area of the colliding nuclei into rows with an area of a nucleon-nucleon cross section $`\sigma _{in}`$. Within each row, each projectile nucleon will collide with each target nucleon. By assuming straight-line space-time trajectories, the space-time diagramof these nucleons is shown in Fig. 2.
We can focus our attention at one spatial point at $`z_0`$ in this row. A nucleon-nucleon collision at this point (represented by an open circle in Fig. 2) will lead to the deposition of energy about that point. At a proper time of about 1 fm/c, the deposited energy will materialize as field quanta of matter in the form of hadrons if the matter is favored to be in the hadronic state. The field quanta will be quarks and gluons if the matter is favored to be in the quark-gluon plasma state.
There are many other nucleon-nucleon collisions which take place sequentially at $`z_0`$ at a time interval of $`d/\gamma `$, as represented by the open circles in Fig. 2. Here, $`d2.5`$ fm is the average nucleon-nucleon spacing in nuclear matter at rest and $`\gamma `$ is the relativistic factor for the motion of the nuclei in the nucleon-nucleon center-of-mass system. Each later nucleon-nucleon collision at the same collision point deposits additional energy. The local energy density increases as a function of time and is approximately proportional to the number of nucleon-nucleon collisions $`N`$ occurring at that point:
$`ϵN{\displaystyle \frac{dn}{dy}}{\displaystyle \frac{m_t}{\sigma _{in}d/\gamma }},`$ (1)
where $`dn/dy`$ is the average multiplicity per unit of rapidity at $`y_{_{CM}}`$$`=`$0 for a nucleon-nucleon collision, and $`m_t`$$``$$`0.35`$ GeV is the transverse mass of a produced pion. For collisions at 158A GeV per nucleon, the energy density deposited per nucleon-nucleon collision as given by Eq. (1) is approximately $`N\times `$1 GeV/fm<sup>3</sup>.
In actual nucleus-nucleus collisions, nucleons lose energy as they collide and the multiplicity distribution in later collisions will be slightly lower than those from earlier collisions. Equation (1), which is defined in terms of an average multiplicity value, is a simplifying, but useful, relation which gives an approximate estimate of the local energy density.
As the size of the colliding nuclei increases, the energy density deposited at a collision point in a central collision will increase. When the number of nucleon-nucleon collisions at a spatial point exceeds a critical number $`N_c`$, the local energy density will increase beyond the critical energy density $`ϵ_c`$ for transition to the quark-gluon plasma phase . Then, the state of lowest free energy becomes the QGP phase with field quanta of quarks, antiquarks, and gluons. If the number of nucleon-nucleon collisions at a point is lower than the critical collision number $`N_c`$ , it will not reach the critical energy density and the field quanta will be hadrons consisting mostly of pions.
## 4 ANOMALOUS FREEZE-OUT VOLUME
In a first-order phase transition, the energy density and pressure change abruptly from the hadron phase to the quark-gluon plasma phase. Consequently, the pressure changes abruptly as a function of the local collision number when the collision number passes the critical value, $`N_c`$.
The pressure in the quark-gluon plasma is given approximately by $`P_{\mathrm{QGP}}37\pi ^2T^4/90`$ which is much greater than the pressure of a pion gas, $`P_{\mathrm{pion}}=6\pi ^2T^4/90`$. Thus, the presence of a quark-gluon plasma is characterized by a region of very high pressure. The large pressure of the quark-gluon plasma can be detected by a large freeze-out volume.
We would like to examine the relation between the initial pressure and the freeze-out volume in Bjorken hydrodynamics . The matter produced in high-energy heavy-ion collisions is subject to a strong longitudinal expansion. We consider a plateau distribution of matter within rapidity range $`|y|<y_0`$ and longitudinal coordinates $`|z|<|z_0|`$. The initial proper time is then $`\tau _0=|z_0|/\mathrm{sinh}y_0`$. The solution for the pressure $`P`$ in Bjorken hydrodynamics is
$`P(\tau ,y)=P(\tau )\theta (y_0|y|),`$ (2)
where
$`{\displaystyle \frac{P(\tau )}{P(\tau _0)}}=\left({\displaystyle \frac{\tau _0}{\tau }}\right)^{4/3}.`$ (3)
(See Eq. (13.17$`b`$) of Ref. .) At the proper time $`\tau `$, the boundary of the matter is expanded up to $`z_\tau =\pm \tau \mathrm{sinh}y_0`$, and thus the ratio of $`|z_\tau |/|z_0|`$ is
$`{\displaystyle \frac{|z_\tau |}{|z_0|}}=\left({\displaystyle \frac{P(\tau _0)}{P(\tau )}}\right)^{3/4}.`$ (4)
As the matter expands, there will be a pressure $`P(\tau _f)`$, the freeze-out pressure, at which particles will no longer interact and will exhibit freeze-out characteristics. Because the volume of matter $`V`$ is proportional to $`|z_\tau |`$, the ratio of the freeze-out volume $`V(\tau _f)`$ at $`\tau _f`$ to the initial volume $`V(\tau _0)`$ is
$`{\displaystyle \frac{V(\tau _f)}{V(\tau _0)}}=\left({\displaystyle \frac{P(\tau _0)}{P(\tau _f)}}\right)^{3/4}.`$ (5)
Hence, the freeze-out volume $`V(\tau _f)`$ varies with the initial pressure as $`[P(\tau _0)]^{3/4}`$.
Equation (5) is applicable to hadron matter. It is also applicable approximately to the QGP undergoing strong longitudinal expansion, for which the temperature of the QGP is lowered so rapidly that the expansion drives the system below the critical temperature. When the energy density and pressure of the supercooled quark-gluon plasma drop down to the same level as those of the hadron matter, spontaneous transition to hadron matter will take place at proper time $`\tau _H`$ resulting in hadron matter of volume $`V(\tau _H)`$. The hadron matter then undergoes further longitudinal expansion to reach the freeze-out volume $`V(\tau _f)`$ at proper time $`\tau _f`$. For the quark-gluon plasma subject to a strong longitudinal expansion, we have
$`{\displaystyle \frac{V(\tau _f)}{V(\tau _0)}}={\displaystyle \frac{V(\tau _f)}{V(\tau _H)}}{\displaystyle \frac{V(\tau _H)}{V(\tau _0)}}\left({\displaystyle \frac{P(\tau _H)}{P(\tau _f)}}\right)^{3/4}\left({\displaystyle \frac{P(\tau _0)}{P(\tau _H)}}\right)^{3/4}=\left({\displaystyle \frac{P(\tau _0)}{P(\tau _f)}}\right)^{3/4}.`$ (6)
The freeze-out volume varies approximately with the initial QGP pressure as $`[P(\tau _0)]^{3/4}`$.
Based on the above, the large pressure of the quark-gluon plasma will lead to a large freeze-out volume and a first-order phase transition will be indicated by an anomalous increase of the freeze-out volume due to the increase in the pressure in the QGP phase. We expect that the anomalous increase in the freeze-out volume should occur when the average number of nucleon-nucleon collisions passes through the critical value $`N_c`$. The average nucleon-nucleon collision number in a collision is a function of centrality or multiplicity, so the freeze-out volume should increase anomalously as a function of centrality or multiplicity in a first-order phase transition. In contrast, the increase should be much more gradual in a second-order phase transition.
In actual heavy-ion collisions, the freeze-out volume also increases with centrality because of the increase in the volume of the overlapping region. Consequently, there will be the systematics of such a normal increase of the freeze-out volume, as one goes from the peripheral region to the central region. However, the increase in the freeze-out volume due to the large pressure of the quark-gluon plasma is an anomalous addition to these systematics. By looking at the systematics of the freeze-out volume as a function of centrality or multiplicity, one should be able to separate out these two components: the normal component due to the increase in the volume of the collision region for peripheral and semi-central collisions, and an anomalous part due to the high pressure of the quark-gluon plasma which becomes important for central collisions.
Recently the NA44 Collaboration observed that the freeze-out volume increased anomalously from semi-central collisions to central collisions, for Pb+Pb collisions at 158A GeV (Fig. 3). This interesting observation by the NA44 Collaboration may be evidence for the quark-gluon plasma. In semi-central collisions, the producded matter are likely to be hadron matter. Therefore, the expansion and the freeze-out volume is governed mainly by the expanding hadron matter with $`P(\tau _0)P_{\mathrm{pion}}`$ and $`V_{\mathrm{freeze}\mathrm{out}}P_{\mathrm{pion}}^{3/4}`$. For central collisions, if quark-gluon plasma is produced, the expansion and cooling for these central collisions is governed essentially by the initial pressure of the quark-gluon plasma, $`P(\tau _0)`$$``$$`P_{QGP}`$,which is a very high pressure, and$`V_{\mathrm{freeze}\mathrm{out}}P_{\mathrm{QGP}}^{3/4}`$. Therefore, this expansion will lead to an anomalous increase in the freeze-out volume for central collisions, much greater than what one would expect from the systematics for semi-central collisions. The rapid increases of the freeze-out volume as a function of multiplicity, as observed by the NA44 Collaboration, may indicate that the hadron-QGP phase transition is a first-order phase transition. Further investigations are needed to shed more light on this interesting observation of the NA44 Collaboration.
Fig. 3. The NA44 data of the freeze-out volume as a function of $`dN^\pm /d\eta `$ . The curves joining the data points are parametrizations in the form of $`(dN^\pm /d\eta )^\alpha `$ from the NA44 Collaboration.
## 5 DIRECT PHOTON PRODUCTION
As the quark-gluon plasma expands, its energy density and temperature decreases. At the point when its temperature decreases below the critical temperature for hadron-QGP phase transition, matter will undergo a phase transition from the quark-gluon plasma to the hadron phase.
During the time when the matter is in the quark-gluon plasma phase, it will emit particles. Photons arising from the electromagnetic interactions of the constituents of the plasma will provide information on the properties of the plasma at the time of their emission. Since photons are hardly absorbed by the medium, they form a relatively ‘clean’ probe to study the state of the quark-gluon plasma. The presence of these photons in high-energy heavy-ion collisions can also possibly provide evidence for the production of the quark-gluon plasma .
Photons are also produced by many other processes in heavy-ion reactions. They can come from the decay of $`\pi ^0`$ and $`\eta ^0`$. As $`\pi ^0`$ particles are copiously produced in strong interactions between nucleons, photons coming from the decay of $`\pi ^0`$ are much more abundant than photons produced by electromagnetic interactions of the constituents of the quark-gluon plasma. The photons from the decay of $`\pi ^0`$ and $`\eta ^0`$ can be subtracted out by making a direct measurement of their yields, obtained by combining pairs of photons. Because of the large number of $`\pi ^0`$ produced, this subtraction is a laborious task, but much progress has been made to provide meaningful results after the subtraction of the photons from the $`\pi ^0`$ and $`\eta ^0`$ backgrounds . Photon measurements obtained after the subtraction of the photons from meson decays are conveniently called measurements of “direct photons”.
Direct photons are produced from the interaction of matter in the QGP phase, a mixed QGP and hadron phase, a pure hadron gas, and hard QCD processes. Different processes give rise to photons in different momentum regions. One may wish to go to the region of photon transverse momentum $`p_T>`$ 2 GeV/c to minimize the contributions from hadrons. If a hot quark-gluon plasma is formed initially, clear signals of photons from the plasma could be visible by examining photons with $`p_T`$ in the range 2$``$3 GeV/c . On the other hand, photons in this region of transverse momenta are also produced by the collision of partons of the projectile nucleons with partons of the target nucleons. Such a contribution must be subtracted in order to infer the net photons from the quark-gluon plasma sources.
Fig. 4. The spectrum of direct photons. The data are from the WA98 Collaboration , and the curves are from Ref. .
Recently, the WA98 Collaboration has measured the photon spectrum for central Pb+Pb collisions at 158A GeV . The photon spectrum from the hard scattering of the nucleons has been calculated previously by taking into account the next-to-leading order contributions and the intrinsic transverse momentum of partons. The results scaled to the Pb+Pb collisions are shown as the dashed curve in Fig. 4. One finds that the WA98 experimental data is in excess of the contributions from the hard scattering of the nucleons.
We can calculate the total photon spectrum including the contributions from the nucleon hard scattering and the quark-gluon plasma. If we assume a quark-gluon plasma with a transverse area of the overlapping area appropriate for the corresponding impact parameter and an initial plasma formation time at $`\tau `$$`=`$1 fm/c, with a temperature of 300 MeV cooling to a critical temperature of 200 MeV, we obtain the photon spectrum in Fig. 9 of Ref. shown here as the solid curve. The good agreement of the model with the WA98 data provides additional support for the production of the quark-gluon plasma in high-energy central Pb+Pb collisions.
## 6 Strangeness enhancement
The strangeness content is enhanced in hadron matter as the temperature increases, but the strangeness is enhanced to an even greater extent in a quark-gluon plasma. The greater strangeness enhancement arises from a higher temperature in the quark-gluon plasma and from its lower effective light quark masses because of the restoration of chiral symmetry. As the strangeness content is greatly enhanced, the probability for the production of multi-strange hyperons will also be greatly enhanced in a quark-gluon plasma .
In a recent measurement of the WA97 Collaboration , the production of multi-strange hyperons is found to be substantially enhanced. In particular, the production of $`\mathrm{\Omega }^{}+\overline{\mathrm{\Omega }^+}`$ in Pb+Pb collisions at 158A GeV is enhanced by up to a factor of 15 relative to that of p+Be. A more detailed description of the strangeness enhancement will be presented by Odyniec in these Proceedings .
Multi-strange hyperons can also be produced by secondary collisions of hadrons. It is known that the collision of the produced pions with nucleons leads to the enhancement of kaons and $`\mathrm{\Lambda }`$ particles. Repeated collisions of kaons in the medium with a $`\mathrm{\Lambda }`$ particle can raise the strangeness of the hyperon by one or more units. The enhancement of hyperons by these secondary collisions increases with the increase in the size of the colliding nuclei, and it is important to take this increase into account. Recent comparison of the hyperon yields as obtained by theoretical RQMD cascade model calculations shows that the RQMD calculations can reproduce the yields of strange particles with one or two units of strangeness, but it underpredicts the yield of $`\mathrm{\Omega }^{}`$ by a factor of 3, and underpredicts the yield of $`\overline{\mathrm{\Omega }^+}`$ by about 40 % .
The discrepancies of the yields of the $`\mathrm{\Omega }`$ hyperons with the RQMD cascade model may be additional evidence for the production of the quark-gluon plasma in Pb+Pb collisions. This may be the case, but such a conclusion relies heavily on the assumed input of many strangeness-raising cross sections, for which no experimental data are available. The evidence will be substantially strengthened if the strangeness-raising cross sections can be better determined by a reliable theoretical model, such as the quark-interchange model of Barnes and Swanson . As emphasized by Ko , the strangeness-raising cross section for the reaction such as $`K^0+\mathrm{\Xi }^{}\pi ^0+\mathrm{\Omega }^{}`$ can proceed through the interchange of a strange quark from $`K`$ to $`\mathrm{\Xi }`$, which need not be OZI-suppressed as in the $`\pi +NK+\mathrm{\Lambda }`$ reaction, where an intermediate strange quark pair is produced. Future theoretical evaluation of these strangeness-raising cross sections will be of great interest in clarifying the origin of the enhancement of the $`\mathrm{\Omega }`$ hyperons in Pb+Pb collisions.
## 7 $`J/\psi `$ Suppression
In a quark-gluon plasma the screening of the charm quark and antiquark will make the $`J/\psi `$ unbound, and its production will be suppressed. The occurrence of $`J/\psi `$ suppression has been suggested as a signature for the quark-gluon plasma . The experimental observation by the NA50 Collaboration of an anomalous $`J/\psi `$ suppression in Pb+Pb collisions has led to the suggestion that this anomalous suppression arises from the production of the quark-gluon plasma . The phenomenon of anomalous $`J/\psi `$ suppression has also been studied by many authors .
The production of $`J/\psi `$ is suppressed not only by the quark-gluon plasma but also by the collision of the $`J/\psi `$ (or its precursor) with nucleons and produced particles. The absorption of $`J/\psi `$ by these particles was considered in a simple analytical model . One divides the transverse area into rows of nucleons of cross section $`\sigma _{in}`$. The projectile nucleons in each row collide with the target nucleons within the same row. One assumes simple straight-line space-time trajectories for the colliding nucleons as in Fig. 2. In this space-time diagram, each nucleon-nucleon collision is a possible source of $`J/\psi `$ production and also a source of produced particles which will absorb $`J/\psi `$. By parametrizing the dissociation cross section of $`J/\psi `$ in its collision with nucleons or produced particles, the $`J/\psi `$ production cross section in a nucleus-nucleus collision can be evaluated.
It was found that the simple model of $`J/\psi `$ absorption by nucleons and produced particles can explain the $`p`$+A, O+A, and S+U data, but not the Pb+Pb data. The deviation of Pb+Pb data from the theoretical extrapolations suggests that there is a transition to a new phase of strong absorption, the quark-gluon plasma, which sets in when the local energy density exceeds a certain threshold. As we remarked before, the energy density at a particular spatial point is approximately proportional to the number of nucleon-nucleon collisions which have taken place at that point. We postulate that the matter at a point makes a transition to the new phase of strong $`J/\psi `$ absorption if there have been $`N_c`$ or more nucleon-nucleon collisions at that point.
Fig. 5. The ratio of the $`J/\psi `$ cross section to the Drell-Yan cross section at different values of transverse energy $`E_T`$. The NA50 data are from Ref. .
We assume that the deconfinement temperature is greater than the dissociation temperatures for the dissociation of both $`\chi `$ and $`J/\psi `$. Then, $`N_c`$$`=`$4 gives very good agreement with the pA, O+A, S+U, and Pb+Pb data. The critical nucleon-nucleon collision number, $`N_c`$$`=`$4, corresponds to a critical energy density of $`ϵ_c`$$`=`$4.2 GeV/fm<sup>3</sup>. The results for the ratio of the $`J/\psi `$ cross section to the Drell-Yan cross section is shown as the solid curve in Fig. 5. The good agreement of this critical energy density model with the WA50 data provides evidence for the production of the quark-gluonplasma in central high-energy Pb+Pb collisions.
The analytical model of Ref. for $`J/\psi `$ is useful to provide new insight into the main characteristics of the suppression process. It is desirable to make a more refined calculation to include many improvements in order to see whether the main characteristics may depend on these refinements. The cross section for the dissociation of $`J/\psi `$ by hadrons should be reliably calculated and their dependence on hadron types and hadron energies included. The population of the $`\rho `$ meson should also be obtained in a dynamical model where the $`\rho `$ mesons are both formed by pion collisions and decay back into pions. The effect of the degradation of the energy of the colliding nucleons and their subsequent lower hadron production should also be taken into account. These improvements can be included in a Monte Carlo cascade model calculation as shown below.
To obtain reliable $`J/\psi `$ dissociation cross sections, we first determine the effective quark-gluon potential from hadron masses and use these interactions in the quark-interchange model of Barnes and Swanson to obtain the dissociation cross section. The Barnes and Swanson model has been tested and found to give good agreement with experimental phase shifts and cross sections in many meson-meson reactions.
By using the Barnes and Swanson model, the cross section for the dissociation of $`J/\psi `$ by a pion is found to be small, with a peak cross section of about 1 mb above the threshold of 640 MeV. On the other hand, the cross section for the dissociation by a $`\rho `$ meson is large because it is above the $`D+\overline{D}`$ threshold. The cross section decreases rapidly as a function of energy. The dependence of the cross sections on the mass of the $`\rho `$ meson have also been calculated .
Having obtained the dissociation cross sections, we incorporate them in a heavy-ion Monte Carlo cascade program to study the survival probability of a produced $`J/\psi `$. We use different basic numerical programs in our Monte Carlo cascade calculations . In one of our programs for which we shall report our results, we base our calculations on the MARCO program, which is a multiple collision model with a stopping law to govern the degradation of the energy of nucleons and a particle production law to specify the distribution of produced particles after anucleon-nucleon collision . The MARCO program has been found to give a good description of the gross features of hadron production .The rescattering of the produced particle is now incorporated to follow the cascade of the hadrons.
Fig. 6. $`J/\psi `$ cross section as a function of $`A^{1/3}(A1)/A+B^{1/3}(B1)/B`$. .
In Fig. 6, we plot the $`J/\psi `$ production cross section as a function of $`A^{1/3}(A1)/A+B^{1/3}(B1)/B`$, where $`A`$ and $`B`$ are the mass numbers of the colliding nuclei. The results from the Monte Carlo cascade calculations are given as the dashed curve. The inclusion of the $`\rho `$ mesons increases the absorption slightly for S+U and Pb+Pb but does not lead to good agreement with the Pb+Pb data. Again, when we postulate that a new phase of $`J/\psi `$ absorption is produced when the local energy density exceeds a critical value, then we obtain the results as shown in the solid curve. The qualitative features of the $`J/\psi `$ absorption are not changed with a refined calculation for $`J/\psi `$ production.
## 8 CONCLUSION AND SUMMARY
In high-energy heavy-ion collisions, the occurrence of the quark-gluon plasma will be accompanied by an anomalous increase in pressure, an excess of direct photon production, an excess of strangeness production, and an anomalous $`J/\psi `$ suppression. We reviewed these signatures and find that they are present in central Pb+Pb collisions at 158A GeV.
For the indicator of the anomalous increase in pressure, the NA44 data exhibits an anomalous increase in the freeze-out volume as a function of multiplicity. Superimposing on the normal systematic increase of the freeze-out volume due to the increase in the volume of the interacting region may be the additional increase due to the high pressure of the quark-gluon plasma.
For the question of direct photon production, the WA98 Collaboration has measured the direct photon spectrum in central Pb+Pb collision, and has found that the photon yield is in excess of the expected yield from nucleon-nucleon hard scattering processes. The excess direct photons can be understood as arising from a quark-gluon plasma with an initial temperature of 300 MeV.
For the enhancement of strangeness, the WA97 Collaboration observed a large enhancement in the yield of $`\mathrm{\Omega }`$ hyperons which cannot be explained by RQMD cascade calculations. Such a discrepancy may support the enhancement as originating from the quark-gluon plasma. This evidence needs to be confirmed by calculating the strangeness-raising cross sections with a reliable model.
In the case of $`J/\psi `$ suppression, the anomalous suppression in central Pb+Pb collisions cannot be explained by the absorption by hadrons. This was demonstrated in a simple analytical model and in a Monte Carlo cascade model using dissociation cross sections obtained from the model of Barnes and Swanson. On the other hand, a critical energy density model can explain the anomalous $`J/\psi `$ suppression, suggesting the production of the quark-gluon plasma in central Pb+Pb collisions at 158A GeV.
In conclusion, our re-assessment of these signatures provides additional support to indicate that, in line with the summary of Heinz and Jacobs , the experimental data appear to be consistent with the occurrence of the quark-gluon plasma in central Pb+Pb collisions at 158A GeV.
The author would like to thank Drs. T. Awes, C. M. Ko, M. Murray, and N. Xu for helpful discussions. The author would like to thank Dr. H. Boggild for his permission to show the data of the NA44 Collaboration (Fig. 3). He would also like to thank his collaborators, Drs. T. Barnes, E. Swanson, S. Sorensen, and B. H. Sa, for their collaboration to study $`J/\psi `$ dissociation cross sections and $`J/\psi `$ suppression.
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# Reeh-Schlieder Defeats Newton-Wigner: On alternative localization schemes in relativistic quantum field theory
## 1 Introduction
Relativistic quantum theory presents us with us a set of peculiar interpretive difficulties over and above the traditional ones of elementary quantum mechanics. For example, while the notion of a “localized object” has a transparent mathematical counterpart in elementary quantum mechanics, it appears that not every aspect of our common-sense notion of localization can be maintained in the context of relativistic quantum theory (cf. Malament 1996). Many of the thorny issues involving localization in relativistic quantum *field* theory have a common formal root in the so-called “Reeh-Schlieder theorem.” Thus, it is of particular philosophical interest that I.E. Segal (1964) and, more recently, G. Fleming (2000) have claimed to be able to avoid the Reeh-Schlieder theorem—and thereby its “counterintuitive” consequences—by means of a judicious reworking of the standard association between observables and regions of space.<sup>1</sup><sup>1</sup>1Saunders 1992 provides an extensive discussion of Segal’s approach, although with different points of emphasis than the current presentation. Although Fleming (1996, 12) appears to dismiss Saunders’ comparison of his and Segal’s approaches, Fleming’s proposal for NW-local algebras (Fleming 2000) (prior to “covariant generalization”) is mathematically identical to Segal’s (1964) proposal.
I am not convinced, however, that Segal and Fleming’s “Newton-Wigner” localization scheme offers any satisfying resolution for the “problem” of localization in relativistic quantum field theory. In particular, the Newton-Wigner localization scheme is not completely immune from the consequences of the Reeh-Schlieder theorem; and neither Segal nor Fleming has offered a conceptually coherent description of the physical meaning of Newton-Wigner localization.
The context of the Reeh-Schlieder theorem is the axiomatic (or algebraic) approach to quantum field theory. This approach singles out a family of postulates that apply quite generally to “physically reasonable” quantum field models, and these postulates are used as a starting point for further structural investigations. One might expect, then, that Segal and Fleming would attempt to undercut the Reeh-Schlieder theorem by questioning one of the assumptions it makes concerning which models are “physically reasonable.” However, Segal and Fleming do not discuss the Reeh-Schlieder theorem at this level of generality; rather, their discussion of the Reeh-Schlieder theorem is restricted to a concrete field model, viz., the free Bose field.
I begin then in Section 2 with a brief review of the global structure of the free Bose field model. In Section 3, I present the standard recipe for assigning observables to regions in space, and I explicate the counterintuitive consequences—stemming from the Reeh-Schlieder theorem—of this standard localization scheme. In Section 4, I present the Newton-Wigner localization scheme and show how it “avoids” the counterintuitive consequences of the Reeh-Schlieder theorem. Finally, in Sections 5 and 6, I argue that Reeh-Schlieder has the final word against the Newton-Wigner localization scheme.
## 2 The free Bose field
In this section, I briefly review the mathematical formalism for the quantum theory of the free Bose field. Although my presentation differs from Fleming’s (2000) in being more abstract and in its emphasis on mathematical rigor, I take it that all parties agree concerning the *global* structures of the free field model (at least in the absence of measurement interactions). That is, we agree on our answers to the following four questions:
1. What is the state space?
2. What are the observables (i.e., physical quantities)?
3. When no measurements are being made, how does the system evolve in time? In other words, what is the (free) Hamiltonian?
4. What is the ground (i.e., vacuum) state?
Disputes arise only at the level of the *local* structure of the free field model; e.g., which states are “localized” in this region of space? In this section, I spell out the answers to questions 1–4. In Section 3, I take up questions concerning localization.
Recall that in its heuristic formulation, the free scalar quantum field is described by an “operator-valued field” $`\mathrm{\Phi }`$ on Minkowski spacetime that solves the Klein-Gordon equation
$$\frac{^2\mathrm{\Phi }}{t^2}+m^2\mathrm{\Phi }=^2\mathrm{\Phi },$$
(1)
and that satisfies the appropriate (equal-time) canonical commutation relations. As is well-known, however, there are mathematical difficulties with understanding $`\mathrm{\Phi }`$ as an operator-valued function. A more rigorous approachtakes $`\mathrm{\Phi }`$ as an “operator-valued distribution.” That is, for each smooth, real-valued test-function $`f`$ on Minkowski spacetime, $`\mathrm{\Phi }(f)`$ can be defined as an operator on some Hilbert space.
For my purposes here, it will be more convenient to turn to another (mathematically equivalent) representation of the field $`\mathrm{\Phi }`$. Let $`C_0^{\mathrm{}}(^3)`$ denote the vector space of smooth, compactly supported functions from $`^3`$ into $``$, and let
$$S=C_0^{\mathrm{}}(^3)C_0^{\mathrm{}}(^3).$$
(2)
Recall now that a scalar-valued solution $`\varphi `$ of the Klein-Gordon equation is uniquely determined by its Cauchy data (i.e., its values, and the values of its first derivative) at any fixed time. Thus, there is a one-to-one correspondence between elements of $`S`$ and (a certain subset of) the space of solutions of the Klein-Gordon equation. Moreover, the conserved four-vector current $`\varphi \underset{\mu }{\overset{}{}}\psi `$ gives rise to a symplectic form $`\sigma `$ on $`S`$:
$$\sigma (u_0u_1,v_0v_1)=_^3(u_0v_1u_1v_0)d^3𝐱.$$
(3)
We let $`D_t`$ denote the natural (inertial) symplectic flow on $`S`$; i.e., $`D_t`$ maps the time-zero Cauchy data of $`\varphi `$ to the time-$`t`$ Cauchy data of $`\varphi `$. The triple $`(S,\sigma ,D_t)`$ contains the essential information specifying the classical theory of the scalar field of mass $`m`$.
A representation of the Weyl form of the canonical commutation relations (CCRs) is a mapping $`fW(f)`$ of $`S`$ into unitary operators acting on some Hilbert space $`𝒦`$ such that $`W(0)=I`$ and
$$W(f)W(g)=e^{i\sigma (f,g)}W(f+g).$$
(4)
I will now sketch the construction of the unique (up to unitary equivalence) “Minkowski vacuum representation” of the CCRs. This construction proceeds in two steps. In *first quantization*, we “Hilbertize” the classical phase space $`S`$, and we “unitarize” the classical dynamical group $`D_t`$. More precisely, suppose that $``$ is a Hilbert space, and that $`U_t`$ is a weakly continuous one-parameter group of unitary operators acting on $``$. Suppose also that the infinitesimal generator $`A`$ of $`U_t`$ is a positive operator; i.e., $`(f,Af)0`$ for all $`f`$ in the domain of $`A`$. If there is a one-to-one real-linear mapping $`K`$ of $`S`$ into $``$ such that
1. $`K(S)+iK(S)`$ is dense in $``$,
2. $`2\mathrm{I}\mathrm{m}(Kf,Kg)=\sigma (f,g)`$,
3. $`U_tK=KD_t`$,
then we say that the triple $`(K,,U_t)`$ is a *one-particle structure* over $`(S,\sigma ,D_t)`$. Constructing a one-particle structure over $`(S,\sigma ,D_t)`$ is a mathematically rigorous version of “choosing the subspace of positive frequency solutions” of the space of complex solutions to the Klein-Gordon equation.
If there is a one-particle structure over $`(S,\sigma ,D_t)`$, then it is unique up to unitary equivalence (Kay 1979). That is, suppose that $`(K,,U_t)`$ and $`(L,\stackrel{~}{},\stackrel{~}{U}_t)`$ are one-particle structures over $`(S,\sigma ,D_t)`$. Then, $`LK^1`$ extends uniquely to a unitary mapping $`V`$ from $``$ onto $`\stackrel{~}{}`$.
It is also not difficult to see that $`V`$ intertwines the unitary groups on the respective Hilbert spaces, i.e., $`VU_t=\stackrel{~}{U}_tV`$. This uniqueness result can be interpreted as showing that the choice of time evolution in the classical phase space suffices to determine uniquely the (first) quantization of the classical system.
I will construct two (unitarily equivalent) versions of the one-particle structure over $`(S,\sigma ,D_t)`$. First, we may complete $`S`$ relative to the unique Hilbert space norm in which time-evolution (given by $`D_t`$) is an isometry. Specifically, let $`H`$ denote the linear operator $`(^2+m^2)^{1/2}`$ on $`C_0^{\mathrm{}}(^3)`$,<sup>2</sup><sup>2</sup>2The mathematically rigorous definition of $`H`$ is as follows: Define the operator $`A=^2+m^2`$ on $`C_0^{\mathrm{}}(^3)`$. Then, $`A`$ is essentially self-adjoint, and the self-adjoint closure $`\overline{A}`$ of $`A`$ is a positive operator with spectrum in $`[m^2,\mathrm{})`$. Using the functional calculus for unbounded operators, we may define $`H=\overline{A}^{1/2}`$, and it follows that the spectrum of $`H`$ is contained in $`[m,\mathrm{})`$. and define a real inner-product $`\mu `$ on $`S`$ by
$`\mu (u_0u_1,v_0v_1)`$ $`=`$ $`(1/2)\left((u_0,Hv_0)+(u_1,H^1v_1)\right)`$ (5)
$`=`$ $`(1/2)\left({\displaystyle _^3}u_0(Hv_0)d^3𝐱+{\displaystyle _^3}u_1(H^1v_1)d^3𝐱\right).`$ (6)
Now let $`_\mu `$ denote the completion of $`S`$ relative to the inner-product $`\mu `$.<sup>3</sup><sup>3</sup>3If $`^\pm (^3)`$ denotes the completion of $`C_0^{\mathrm{}}(^3)`$ relative to the inner product $`(,H^{\pm 1})`$, then $`_\mu =^+(^3)^{}(^3)`$. Define an operator $`J`$ on $`_\mu `$ by setting
$$J(u_0u_1)=H^1u_1Hu_0,$$
(7)
on the dense subset $`S`$ of $`_\mu `$. Clearly $`J^2=I`$, i.e., $`J`$ is a “complex structure” on $`_\mu `$. Thus, $`_\mu `$ becomes a complex vector space when we define scalar multiplication by $`(a+ib)f=af+J(bf)`$, and is a complex Hilbert space relative to the inner-product
$`(f,g)_\mu `$ $`=`$ $`\mu (f,g)+i\mu (Jf,g)`$ (8)
$`=`$ $`\mu (f,g)+(i/2)\sigma (f,g).`$ (9)
Finally, it can be shown that $`[J,D_t]=0`$, so that $`D_t`$ extends uniquely to a weakly continuous one-parameter group of *unitary* operators (denoted again by $`D_t`$) on the complex Hilbert space $`_\mu `$. Therefore, $`(\iota ,_\mu ,D_t)`$, with $`\iota `$ the identity mapping, is a one-particle structure over $`(S,\sigma ,D_t)`$.
It may not be immediately obvious—especially to those accustomed to non-relativistic quantum mechanics—how to tie the physics of localization to the mathematical structure of the Hilbert space $`_\mu `$. (For example, which vectors in $`_\mu `$ are localized in a given spatial region?) The Newton-Wigner one-particle structure brings us back to familiar territory by using the space $`L^2(^3)`$ as the concrete representation of the one-particle space. In particular, define the mapping $`K:SL^2(^3)`$ by
$$K(u_0u_1)=\mathrm{\hspace{0.25em}2}^{1/2}(H^{1/2}u_0+iH^{1/2}u_1).$$
(10)
It is then straightforward to check that the complex-linear span of $`K(S)`$ is dense in $`L^2(^3)`$, and that $`K`$ preserves (modulo a factor of $`2`$) the symplectic form $`\sigma `$. Moreover, it can be shown that $`K`$ intertwines $`D_t`$ with the one parameter unitary group $`U_t=e^{itH}`$ on $`L^2(^3)`$. Therefore, $`(K,L^2(^3),U_t)`$ is a one-particle structure over $`(S,\sigma ,D_t)`$.
Since $`(\iota ,_\mu ,D_t)`$ and $`(K,L^2(^3),U_t)`$ are one-particle structures over $`(S,\sigma ,D_t)`$, it follows that $`(K\iota ^1)=K`$ extends uniquely to a unitary operator $`V`$ from $`_\mu `$ onto $`L^2(^3)`$:
Thus, the one-particle spaces $`(_\mu ,D_t)`$ and $`(L^2(^3),U_t)`$ are mathematically, and hence physically, equivalent. On the other hand, the two spaces certainly *suggest* different notions of localization.
### 2.1 Second quantization
Once we have a one-particle space $`(,U_t)`$ in hand, the movement to a quantum *field* theory (i.e., “second quantization”) is mathematically straightforward and uniquely determined.<sup>4</sup><sup>4</sup>4For a more detailed exposition, see Bratteli and Robinson 1997, Section 5.2. In particular, let $`()`$ denote the “Fock space” over $``$. That is,
$$()=^2^3\mathrm{},$$
(11)
where $`^n`$ is the $`n`$-fold symmetric tensor product of $``$. As usual we let
$$\mathrm{\Omega }=100\mathrm{}$$
(12)
denote the vacuum vector in $`()`$. For each $`f`$, we define the creation $`a^+(f)`$ and annihilation $`a(f)`$ operators on $`()`$ as usual, and we let $`\mathrm{\Phi }(f)`$ denote the self-adjoint closure of the unbounded operator
$$2^{1/2}(a(f)+a^+(f)).$$
(13)
If we let $`W(f)=\mathrm{exp}\{i\mathrm{\Phi }(f)\}`$, then the $`W(f)`$ satisfy the Weyl form of the canonical commutation relations:
$$W(f)W(g)=e^{i\mathrm{Im}(f,g)/2}W(f+g),$$
(14)
and vacuum expectation values are given explicitly by
$$\mathrm{\Omega },W(f)\mathrm{\Omega }=\mathrm{exp}\left(f^2/4\right).$$
(15)
The dynamical group on $`()`$ is given by the “second quantization” $`\mathrm{\Gamma }(U_t)=e^{itd\mathrm{\Gamma }(H)}`$ of the dynamical group $`U_t=e^{itH}`$ on $``$, and the vacuum vector $`\mathrm{\Omega }`$ is the unique eigenvector of the Hamiltonian $`d\mathrm{\Gamma }(H)`$ with eigenvalue $`0`$.
## 3 Local algebras and the Reeh-Schlieder theorem
To this point we have only discussed the global structure of the free Bose field model. The physical observables for the free Bose field are given by the self-adjoint operators on Fock space $`()`$. We equip this model with a *local structure* when we define a correspondence between regions in space and “subalgebras” of observables. This labelling may be done for various purposes, but the traditional motivation was to indicate those observables that can (in theory) be measured in that region of space.
Now, each real-linear subspace $`E`$ of the one-particle space $``$ gives rise naturally to a subalgebra of operators, viz., the algebra generated by the Weyl operators $`\{W(f):fE\}`$. Thus, a localization scheme needs only to determine which real-linear subspace of $``$ should be taken as corresponding to a region $`G`$ in physical space. *It is on this point that the Newton-Wigner localization scheme disagrees with the standard localization scheme.* In the remainder of this section, I discuss the standard localization scheme and its consequences.
The standard localization scheme assigns to the spatial region $`G`$ the subset $`S(G)_\mu `$ of Cauchy data localized in $`G`$. That is, if $`C^{\mathrm{}}(G)`$ denotes the subspace of $`C_0^{\mathrm{}}(^3)`$ of functions with support in $`G`$, then
$$S(G)=C^{\mathrm{}}(G)C^{\mathrm{}}(G),$$
(16)
is a real-linear subspace of $`_\mu `$. (Note that $`S(G)`$ is not closed nor, as we shall soon see, complex-linear.) Thus, in the Newton-Wigner representation, the classical localization scheme assigns $`G`$ to the real-linear subspace $`V(S(G))`$ of $`L^2(^3)`$. When no confusion can result, I will suppress reference to the unitary operator $`V`$ and simply use $`S(G)`$ to denote the pertinent subspace in either concrete version of the one-particle space.
Note that the correspondence $`GS(G)`$ is monotone; i.e., if $`G_1G_2`$ then $`S(G_1)S(G_2)`$. Moreover, if $`G_1G_2=\mathrm{}`$, then $`S(G_1)`$ and $`S(G_2)`$ are “symplectically orthogonal.” That is, if $`fS(G_1)`$ and $`gS(G_2)`$, then $`\mathrm{Im}(f,g)=0`$. Indeed, if $`u_0u_1S(G_1)`$ and $`v_0v_1S(G_2)`$, then
$$\sigma (u_0u_1,v_0v_1)=_^3(u_0v_1u_1v_0)d^3𝐱=0,$$
(17)
since the $`u_i`$ and $`v_i`$ have disjoint regions of support.
Now, we say that a Weyl operator $`W(f)`$ acting on $`()`$ is *classically localized* in $`G`$ just in case $`fS(G)`$. (“Classically” here refers simply to the fact that our notion of localization is derived from the local structure of the classical phase space $`S`$.) Let $`𝐁(())`$ denote the algebra of bounded operators on $`()`$. We then define the subalgebra $`(G)𝐁(())`$ of operators classically localized in $`G`$ to be the “von Neumann algebra” generated by the Weyl operators classically localized in $`G`$. That is, $`(G)`$ consists of arbitrary linear combinations and “weak limits” of Weyl operators classically localized in $`G`$.<sup>5</sup><sup>5</sup>5Since $`fW(f)`$ is weakly continuous, $`(G)`$ contains $`W(f)`$ for all $`f`$ in the closure of $`S(G)`$.
If $`𝐁(())`$, we let $`^{}`$ denote all operators in $`𝐁(())`$ that commute with every operator in $``$. If $``$ contains $`I`$ and is closed under taking adjoints, then von Neumann’s “double commutant theorem” entails that $`(^{})^{}`$ is the von Neumann algebra generated by $``$. Thus, we have
$$(G)=\{W(f):fS(G)\}^{\prime \prime }.$$
(18)
In order also to associate unbounded operators with local regions, we say that an unbounded operator $`A`$ is *affiliated* with the local algebra $`(G)`$ just in case $`U^1AU=A`$ for any unitary operator $`U(G)^{}`$. It then follows that $`\mathrm{\Phi }(f)`$ is affiliated with $`(G)`$ just in case $`W(f)(G)`$.
The correspondence $`G(G)`$ clearly satisfies isotony. That is, if $`G_1G_2`$ then $`(G_1)(G_2)`$. Moreover, the local algebras also satisfy fixed-time microcausality. That is, if $`G_1G_2=\mathrm{}`$ then all operators in $`(G_1)`$ commute with all operators in $`(G_2)`$. (This follows directly from Eq. (14) and the fact that $`S(G_1)`$ and $`S(G_2)`$ are symplectically orthogonal.)
### 3.1 Anti-locality and the Reeh-Schlieder theorem
Let $``$ be some subalgebra of $`𝐁(())`$. We say that a vector $`\psi ()`$ is *cyclic* for $``$ just in case $`[\psi ]=()`$, where $`[\psi ]`$ denotes the closed linear span of $`\{A\psi :A\}`$. Of course, every vector in $`()`$, including the vacuum vector $`\mathrm{\Omega }`$, is cyclic for the global algebra $`𝐁(())`$ of all bounded operators on $`()`$. The Reeh-Schlieder theorem, however, tells us that the vacuum vector $`\mathrm{\Omega }`$ is cyclic for any *local* algebra $`(G)`$.
The first version of the Reeh-Schlieder theorem I will present is a restricted version of the theorem—due to Segal and Goodman—applicable only to the free Bose field model. The key concept in this version of the theorem is the notion of an “anti-local” operator.
###### Definition.
An operator $`A`$ on $`L^2(^3)`$ is said to be anti-local just in case: For any $`fL^2(^3)`$ and for any open subset $`G`$ of $`^3`$, $`\mathrm{supp}(f)G=\mathrm{}`$ and $`\mathrm{supp}(Af)G=\mathrm{}`$ only if $`f=0`$.
Thus, in particular, an anti-local operator maps any wavefunction with support inside a bounded region to a wavefunction with infinite “tails.”
The following lemma may be the most important lemma for understanding the local structure of the free Bose field model:
###### Lemma (Segal and Goodman 1965).
The operator $`H=(^2+m^2)^{1/2}`$ is anti-local.
This lemma has the important consequence that for any non-empty open subset $`G`$ of $`^3`$, the *complex*-linear span of $`S(G)`$ is dense in $``$ (cf. Segal and Goodman 1965, Corollary 1). However, for any real-linear subspace $`E`$ of $``$, $`\mathrm{\Omega }`$ is cyclic for the algebra generated by $`\{W(f):fE\}`$ if and only if the complex-linear span of $`E`$ is dense in $``$ (cf. Petz 1990, Proposition 7.7). Thus, the anti-locality of $`H`$ entails that $`\mathrm{\Omega }`$ is cyclic for every local algebra.
###### Reeh-Schlieder Theorem.
Let $`G`$ be any nonempty open subset of $`^3`$. Then, $`\mathrm{\Omega }`$ is cyclic for $`(G)`$.
What is the significance of this cyclicity result? Segal (1964, 140) claims that the theorem is “striking,” since it entails that
> “…the entire state vector space of the field could be obtained from measurements in an arbitrarily small region of space-time!”
He then goes on to claim that the result is, “quite at variance with the spirit of relativistic causality” (143). Fleming also sees the cyclicity result as counterintuitive, apparently because it does not square well with our understanding of relativistic causality. For example (cf. Fleming 2000, 4), the Reeh-Schlieder theorem entails that for any state $`\psi ()`$, and for any predetermined $`ϵ`$, there is an operator $`A(G)`$ such that $`A\mathrm{\Omega }\psi <ϵ`$. In particular, $`\psi `$ may be a state that differs from the vacuum only in some region $`G^{}`$ that is disjoint (and hence spacelike separated) from $`G`$. If, then, $`A`$ is interpreted as an “operation” that can be performed in the region $`G`$, it follows that operations performed in $`G`$ can result in arbitrary changes of the state in the region $`G^{}`$. This, then, is taken by Fleming to show that, “the local fields allow the possibility of arbitrary space-like distant effects from arbitrary localized actions” (Fleming 2000, 20).
Fleming’s use of “actions” and “effects” seems to construe a local operation—represented by an operator $`A(G)`$—as a purely *physical* disturbance of the system; i.e., the operation here is a *cause* with an *effect* at spacelike separation. If this were the only way to think of local operations, then I would grant that the Reeh-Schlieder theorem is counterintuitive, and indeed very contrary to the spirit of relativisitic causality. However, once one makes the crucial distinction between selective and nonselective local operations, local cyclicity does not obviously conflict with relativistic causality (cf. Clifton and Halvorson 2000, Section 2). Rather than dwell on that here, however, I will proceed to spell out some of the further “counterintuitive” consequences of the Reeh-Schlieder theorem.
1. Let $`G_1`$ and $`G_2`$ be disjoint subsets of $`^3`$. Suppose that $`W(f)`$ is classically localized in $`G_1`$ and $`W(g)`$ is classically localized in $`G_2`$. Then, $`\mathrm{Im}(f,g)=0`$ and therefore $`W(f)W(g)=W(f+g)`$. Thus,
$`\mathrm{\Omega },W(f)W(g)\mathrm{\Omega }`$ $`=`$ $`\mathrm{exp}(f+g^2/4)`$ (19)
$`=`$ $`\mathrm{\Omega },W(f)\mathrm{\Omega }\mathrm{\Omega },W(g)\mathrm{\Omega }e^{\mathrm{Re}(f,g)/2}.`$ (20)
However, $`S(G_1)`$ and $`S(G_2)`$ are not orthogonal relative to the real part of the inner product $`(,)`$. Indeed, if $`f=u_0u_1`$ and $`g=v_0v_1`$, then
$`\mathrm{Re}(f,g)`$ $`=`$ $`(u_0,Hv_0)+(u_1,H^1v_1)`$ (21)
$`=`$ $`{\displaystyle _^3}u_0(Hv_0)d^3𝐱+{\displaystyle _^3}u_1(H^1v_1)d^3𝐱.`$ (22)
But since $`H`$ and $`H^1`$ are anti-local, the two integrals in (22) will not generally vanish. Therefore, the vacuum state is not a product state across $`(G_1)`$ and $`(G_2)`$.
It should be noted, however, that the above argument does not entail that the vacuum state is “entangled”—since it could still be a *mixture* of product states across $`(G_1)`$ and $`(G_2)`$. However, it can be shown directly from the cyclicity of the vacuum vector $`\mathrm{\Omega }`$ that the vacuum state is not even a mixture of product states across $`(G_1)`$ and $`(G_2)`$ (Halvorson and Clifton 2000). Moreover, the vacuum predicts a maximal violation of Bell’s inequality relative to the algebras $`(G)`$ and $`(G^{})`$, where $`G^{}=^3\backslash G`$ (Summers and Werner 1985). (Bell correlation, however, is not entailed by cyclicity.)
2. The cyclicity of the vacuum combined with (equal-time) microcausality entails that the vacuum vector is *separating* for any local algebra $`(G)`$, where $`G^{}`$ has non-empty interior. That is, for any operator $`A(G)`$, if $`A\mathrm{\Omega }=0`$ then $`A=0`$. In particular, for any local event—represented by projection operator $`P(G)`$—the probability that event will occur in the vacuum state is nonzero. Thus, the vacuum is “seething with activity” at the local level.
Since the vacuum is entangled across $`(G)`$ and $`(G^{})`$, it follows that the vacuum is a mixed state when restricted to the local algebra $`(G)`$. In fact, when we restrict the vacuum to $`(G)`$, it is *maximally mixed* in the sense that the vacuum may be written as a mixture with any one of a dense set of states of $`(G)`$ (Clifton and Halvorson 2000). Intuitively speaking, then, the vacuum state provides minimal information about local states of affairs. This is quite similar to the singlet state, which restricts to the maximally mixed state $`(1/2)I`$ on either one-particle subsystem (cf. Redhead 1995a).
3. For any annihilation operator $`a(f)`$, we have $`a(f)\mathrm{\Omega }=0`$. Thus, $`a(f)`$ cannot be affiliated with the local algebra $`(G)`$. Since the family of operators affiliated with $`(G)`$ is closed under taking adjoints, it also follows that no creation operators are affiliated with $`(G)`$.
The concreteness of the model we are dealing with allows a more direct understanding of why, mathematically speaking, local algebras do not contain creation and annihilation operators. Inverting the relation in (13), and using the fact that $`fa^+(f)`$ is linear and $`fa(f)`$ is anti-linear, it follows that
$`a^+(f)`$ $`=`$ $`2^{1/2}(\mathrm{\Phi }(f)i\mathrm{\Phi }(if)),`$ (23)
$`a(f)`$ $`=`$ $`2^{1/2}(\mathrm{\Phi }(f)+i\mathrm{\Phi }(if)).`$ (24)
Thus, an algebra generated by the operators $`\{W(f):fE\}`$, will contain the creation and annihilation operators $`\{a^+(f),a(f):fE\}`$ only if $`E`$ is a *complex*-linear subspace of $``$. This is *not* the case for a local algebra $`(G)`$ where $`E=S(G)`$ is a *real*-linear subspace of $``$. In fact, referring to the concrete one-particle space $`_\mu `$ allows us to see clearly that $`S(G)`$ is not invariant under the complex structure $`J`$. If $`u_0u_1S(G)`$, then
$$J(u_0u_1)=H^1u_1Hu_0.$$
(25)
But since $`H`$ and $`H^1`$ are anti-local, it is not the case that $`Hu_0C^{\mathrm{}}(G)`$ or $`H^1u_1C^{\mathrm{}}(G)`$. Thus, $`JfS(G)`$ when $`fS(G)`$. What is more, since the complex span of $`S(G)`$ is dense in $`_\mu `$, if $`S(G)`$ *were* a complex subspace, then it would follow that $`(G)=𝐁(())`$.
3. Number operators also annihilate the vacuum. Since the vacuum is separating for local algebras, no number operator is affiliated with any local algebra. Thus, an observer in the region $`G`$ cannot count the number of particles in $`G`$!
How should we understand the inability of local observers to count the number of particles in their vicinity? According to Redhead (1995b), a heuristic calculation shows that the local number density operator $`N_G`$ does not commute with the density operator $`N_G^{}`$ (where $`G^{}`$ is the complement of $`G`$). Thus, he claims that
> “…it is usual in axiomatic formulations of quantum field theory to impose a microcausality condition on physically significant local observables, *viz* that the associated operators *should* commute at space-like separation. The conclusion of this line of argument is that number densities are not physical observables, and hence we do not have to bother about trying to interpret them.” (81)
While Redhead’s conclusion is correct, it is instructive to note that his reasoning cannot be reproduced in a mathematically rigorous fashion. That is, there are *no* local number density operators—in particular, neither $`N_G`$ nor $`N_G^{}`$ exist—and so it cannot be literally true that $`N_G`$ and $`N_G^{}`$ fail to commute.
In order to see this, consider first the (single wavefunction) number operator $`N_f=a^+(f)a(f)`$, where $`f`$ is “classically localized” in $`G`$, i.e., $`fS(G)`$. Since $`fa^+(f)`$ is linear, and $`fa(f)`$ is anti-linear, it follows that $`N_f=N_{(e^{it}f)}`$ for all $`t`$. That is, a single wavefunction number operator $`N_f`$ is invariant under phase tranformations of $`f`$. However, classical localization of a wavefunction is *not* invariant under phase transformations. Thus, it is not possible to formulate a well-defined notion of classical localization for a single wavefunction number operator.
How, though, do we define a number density operator $`N_G`$? Heuristically, one sets
$$N_G=_GN(𝐱)d^3𝐱,$$
(26)
where $`N(𝐱)=a^+(𝐱)a(𝐱)`$. Since, however, $`N(𝐱)`$ is not a well-defined mathematical object, Eq. (26) is a purely formal expression. Thus, we replace $`N(𝐱)`$ with the single wavefunction number operator $`N_f`$ and we set,
$$N_G=\underset{i}{}N_{f_i},$$
(27)
where $`f_i`$ is a basis of the real-linear subspace $`S(G)`$ of $``$.<sup>6</sup><sup>6</sup>6Actually, this infinite sum is also a formal expression, since it sums unbounded operators. A technically correct definition would define $`N_G`$ as an upper bound of quadratic forms (see Bratteli and Robinson 1997). Using the fact that $`N_f=N_{if}`$ for each $`f`$, it follows then that $`N_G=N_{[G]}`$, where $`N_{[G]}`$ is the number operator for the closed complex-linear span $`[S(G)]`$ of $`S(G)`$ in $``$; and the anti-locality of $`H`$ entails that $`[S(G)]=`$. Therefore, the operator we defined in Eq. (27) turns out to be the *total* number operator $`N`$.
4. The Reeh-Schlieder theorem also has implications for the *internal* structure of the local algebra $`(G)`$. In particular, the local algebra $`(G)`$ is what is called a “type III” von Neumann algebra (Araki 1964). (The algebra $`𝐁(())`$ of all bounded operators on $`()`$ is called a type I von Neumann algebra.) From a physical point of view, this is significant since type III algebras contain only infinite-dimensional projections—which entails that there are strict limits on our ability to “isolate” a local system from outside influences (Clifton and Halvorson 2000). Type III algebras also have *no* pure (normal) states.
## 4 Newton-Wigner localization
In the previous section, we saw that the standard localization scheme $`G(G)`$ has a number of “counterintuitive” features, all of which follow from the Reeh-Schlieder theorem. These counterintuitive features prompted Segal (1964) and Fleming (2000) to suggest a reworking of the correspondence between spatial regions and subalgebras of observables. In this section I give a mathematically rigorous rendering of the Segal-Fleming proposal, and I show how it avoids both the Reeh-Schlieder theorem and its consequences. (Here I deal only with Fleming’s first proposal, prior to his generalization to “covariant fields.”)
Recall that a localization scheme defines a correspondence between regions in space and real-linear subspaces of the one-particle space $``$. The Newton-Wigner localization scheme defines this correspondence in precisely the way it is done in elementary quantum mechanics: A region $`G`$ in $`^3`$ corresponds to the subspace $`L^2(G)L^2(^3)`$ of wavefunctions with probability amplitude vanishing (almost everywhere) outside of $`G`$. We may then use the unitary mapping $`V`$ between $`_\mu `$ and $`L^2(^3)`$ to identify the subspace $`V^1L^2(G)`$ of Newton-Wigner localized wavefunctions in $`_\mu `$. Hereafter, I will suppress reference to $`V^1`$ and use $`L^2(G)`$ to denote the pertinent subspace in either concrete version of the one-particle space.
Note that the correspondence $`GL^2(G)`$ is monotone; i.e., if $`G_1G_2`$ then $`L^2(G_1)L^2(G_2)`$. Moreover, if $`G_1G_2=\mathrm{}`$, then $`L^2(G_1)`$ and $`L^2(G_2)`$ are *fully* orthogonal—a key difference between NW localization and classical localization.
Now, we say that a Weyl operator $`W(f)`$ acting on $`()`$ is *NW-localized* in $`G`$ just in case $`fL^2(G)`$. We then define the algebra $`_{NW}(G)`$ of NW-localized operators on $`()`$ as the von Neumann algebra generated by the Weyl operators NW-localized in $`G`$. That is,
$$_{NW}(G)=\{W(f):fL^2(G)\}^{\prime \prime }.$$
(28)
Clearly, the correspondence $`G_{NW}(G)`$ satisfies isotony. Moreover, since $`G_1G_2=\mathrm{}`$ entails that $`L^2(G_1)`$ and $`L^2(G_2)`$ are orthogonal subspaces of $``$, the correspondence $`G_{NW}(G)`$ satisfies fixed-time microcausality. Thus, at least in this fixed-time formulation, the NW localization scheme appears to have all the advantages of the classical localization scheme. I will now proceed to spell out some features of the NW localization scheme that may make it seem *more* attractive than the standard localization scheme.
If $`G`$ is an open subset of $`^3`$, then
$$L^2(^3)=L^2(GG^{})=L^2(G)L^2(G^{}).$$
(29)
Accordingly, if we let $`_G=(L^2(G))`$ and $`_G^{}=(L^2(G^{}))`$ then it follows that
$`()`$ $`=`$ $`_G_G^{}.`$ (30)
(Here the equality sign is intended to denote that there is a natural isomorphism between $`()`$ and $`_G_G^{}`$.) Moreover, the vacuum vector $`\mathrm{\Omega }()`$ is the product $`\mathrm{\Omega }_G\mathrm{\Omega }_G^{}`$ of the respective vacuum vectors in $`_G`$ and $`_G^{}`$. By definition, $`\mathrm{\Phi }(f)`$ is affiliated with $`_{NW}(G)`$ when $`fL^2(G)`$. Since $`L^2(G)`$ is a *complex*-linear subspace of $``$, it follows that $`\mathrm{\Phi }(if)`$ is also affiliated with $`_{NW}(G)`$, and hence that $`a^+(f),a(f),`$ and $`N_f`$ are all affiliated with $`_{NW}(G)`$. If we let $`U`$ denote the unitary operator that maps $`_G_G^{}`$ naturally onto $`()`$, then it is not difficult to see that
$$U^1a^+(f)U=a_G^+(f)I,$$
(31)
where $`a_G^+(f)`$ is the creation operator on $`_G`$. Thus, we also have $`U^1a(f)U=a_G(f)I`$, and since the creation and annihilation operators $`\{a_G^\pm (f):fL^2(G)\}`$ form an irreducible set of operators on $`_G`$, it follows that
$`_{NW}(G)`$ $`=`$ $`𝐁(_G)I,`$ (32)
$`_{NW}(G^{})`$ $`=`$ $`I𝐁(_G^{}).`$ (33)
(Again, equality here means there is a natural isomorphism.)
It follows then that acting on $`\mathrm{\Omega }=\mathrm{\Omega }_G\mathrm{\Omega }_G^{}`$ with elements from $`_{NW}(G)`$ results only in vectors of the form $`\psi \mathrm{\Omega }_G^{}`$ for some $`\psi _G`$. Thus, the vacuum is *not* cyclic for the local algebra $`_{NW}(G)`$.
1. It is obvious from the preceding that the vacuum is a product state across $`_{NW}(G)`$ and its complement $`_{NW}(G^{})`$. This also follows directly from the fact that $`L^2(G)`$ and $`L^2(G^{})`$ are fully orthogonal subspaces of $``$. Indeed, let $`W(f)_{NW}(G)`$ and $`W(g)_{NW}(G^{})`$. Then since $`f+g^2=f^2+g^2`$, it follows that
$`\mathrm{\Omega },W(f)W(g)\mathrm{\Omega }`$ $`=`$ $`\mathrm{\Omega },W(f+g)\mathrm{\Omega }`$ (34)
$`=`$ $`\mathrm{exp}(f+g^2/4)`$ (35)
$`=`$ $`\mathrm{\Omega },W(f)\mathrm{\Omega }\mathrm{\Omega },W(g)\mathrm{\Omega }.`$ (36)
2. Restricting the vacuum state $`\mathrm{\Omega }`$ to $`_{NW}(G)`$ is equivalent to restricting the product state $`\mathrm{\Omega }_G\mathrm{\Omega }_G^{}`$ to $`𝐁(_G)I`$. Thus, the restriction of $`\mathrm{\Omega }`$ to $`_{NW}(G)`$ is pure, and the global vacuum provides a “maximally specific” description of local states of affairs.
3. If $`\{f_i\}`$ is an orthonormal basis of $`L^2(G)`$, then the number operator $`N_G=_iN_{f_i}`$ is affiliated with $`_{NW}(G)`$. Moreover, the number operator $`N_G^{}`$ is affiliated with $`_{NW}(G^{})`$, and by microcausality we have $`[N_G,N_G^{}]=0`$. We may also see this by employing the correspondence between $`()`$ and $`_G_G^{}`$. The Fock space $`_G`$ has its own total number operator $`\stackrel{~}{N}_G`$. Similarly, $`_G^{}`$ has its own total number operator $`\stackrel{~}{N}_G^{}`$. Obviously then, $`\stackrel{~}{N}_GI`$ is affiliated with $`𝐁(_G)I`$, and $`I\stackrel{~}{N}_G^{}`$ is affiliated with $`I𝐁(_G^{})`$. Just as obviously, $`\stackrel{~}{N}_GI`$ commutes with $`I\stackrel{~}{N}_G^{}`$.
4. As can be seen from Eq. (32), the local algebra $`_{NW}(G)`$ is a type I von Neumann algebra. According to Segal (1964, 140), this is precisely the structure of local algebras that is “suggested by considerations of causality and empirical accessibility.”
## 5 The full strength of Reeh-Schlieder
The results of the previous two sections speak for themselves: The Newton-Wigner localization scheme results in a mathematical structure that appears to be much more in accord with our a priori physical intuitions than the structure obtained from the standard localization scheme. In this section, however, I show that the NW localization scheme “avoids” the Reeh-Schlieder theorem in only a trivial sense, and I show that the NW localization scheme has its own counterintuitive features without parallel in the standard localization scheme.
First, while the NW-local algebras avoid cyclicity of the vacuum vector, they still have a dense set of cyclic vectors.<sup>7</sup><sup>7</sup>7Cf. Fleming’s claim that, “…it is remarkable that *any state* can have enough structure within an arbitrarily small region, $`O`$, to enable even the mathematical reconstituting of essentially the whole state space” (Fleming 2000, 5).
###### Theorem 1.
$`_{NW}(G)`$ has a dense set of cyclic vectors in $`()`$.
###### Proof.
Since the Hilbert spaces $`_G`$ and $`_G^{}`$ have the same (infinite) dimension, it follows from Theorem 4 of (Clifton et al. 1998) that $`_{NW}(G)=𝐁(_G)I`$ has a dense set of cyclic vectors in $`()=_G_G^{}`$. ∎
Thus, if the worry about the Reeh-Schlieder theorem is about cyclicity in general, adopting the NW localization scheme does nothing to alleviate this worry.
Perhaps, however, the worry about the Reeh-Schlieder theorem is specifically a worry about cyclicity of the *vacuum* state. (One wonders, though, why this would be worse than cyclicity of any other state.) Even so, I argue now that the NW localization scheme does not avoid the “vacuum-specific” consequences of the full Reeh-Schlieder theorem.
Let $`𝒦`$ be an arbitrary Hilbert space, representing the state space of some quantum field theory. (For example, $`𝒦=()`$ in the case of the free Bose field.) Suppose also that there is a representation $`𝐚U(𝐚)`$ of the spacetime translation group in the group of unitary operators on $`𝒦`$. Given such a representation, there is a “four operator” $`𝐏`$ on $`𝒦`$ such that $`U(𝐚)=e^{i𝐚𝐏}`$. We say that the representation $`𝐚U(𝐚)`$ satisfies the *spectrum condition* just in case the spectrum of $`𝐏`$ is contained in the forward light cone. From a physical point of view, the spectrum condition corresponds to the assumption that (a) all physical effects propagate at velocities at most the speed of light, and (b) energy is positive. Note, consequently, that the spectrum condition is a purely global condition, and so is not likely to be a source of dispute between proponents of differing localization schemes.
A *net of local observable algebras* is an assigment $`O𝔄(O)`$ of open regions in Minkowski spacetime to von Neumann subalgebras of $`𝐁(𝒦)`$. (Note that this definition is not immediately pertinent to the localization schemes presented in Sections 3 and 4, since they gave an assignment of algebras to open regions in space at a fixed time.) The full Reeh-Schlieder theorem will apply to this net if it satisfies the following postulates:
1. *Isotony:* If $`O_1O_2`$, then $`𝔄(O_1)𝔄(O_2)`$.
2. *Translation Covariance:* $`U(𝐚)^1𝔄(O)U(𝐚)=𝔄(O+𝐚)`$.
3. *Weak Additivity:* For any open $`OM`$, the set
$$\underset{𝐚M}{}U(𝐚)^1𝔄(O)U(𝐚)$$
of operators is irreducible (i.e., leaves no subspace of $`𝒦`$ invariant).
In this general setting, a vacuum vector $`\mathrm{\Omega }`$ can be taken to be any vector invariant under all spacetime translations $`U(𝐚)`$.
###### Full Reeh-Schlieder Theorem.
Suppose that $`\{𝔄(O)\}`$ is a net of local observable algebras satisfying postulates 1–3. Then, for any open region $`O`$ in Minkowski spacetime, $`\mathrm{\Omega }`$ is cyclic for $`𝔄(O)`$.
Note that the Reeh-Schlieder theorem does *not* require the postulate of microcausality (i.e., if $`A𝔄(O_1)`$ and $`B𝔄(O_2)`$, where $`O_1`$ and $`O_2`$ are spacelike separated, then $`[A,B]=0`$).<sup>8</sup><sup>8</sup>8To see that microcausality is logically independent from postulates 1–3, take the trivial localization scheme: $`𝔄(O)=𝐁(𝒦)`$, for each $`O`$.
For the standard localization scheme, there is a straightforward connection between the full Reeh-Schlieder theorem and the fixed-time version given in Section 3. In particular, there is an alternative method for describing the standard localization scheme that involves appeal to spacetime regions rather than space regions at a fixed time (see Horuzhy 1988, Chapter 4). It then follows that $`(G)=𝔄(O_G)`$, where $`O_G`$ is the “domain of dependence” of the spatial region $`G`$. Thus, the fixed-time version of the Reeh-Schlieder theorem may be thought of as corollary of the full Reeh-Schlieder theorem in connection with the fact that $`(G)=𝔄(O_G)`$.
Segal and Fleming avoid the fully general version of the Reeh-Schlieder theorem only by remaining silent about how we ought to assign algebras of observables to open regions of *spacetime*.<sup>9</sup><sup>9</sup>9It is essential for the proof of the full Reeh-Schlieder theorem that the region $`O`$ has some “temporal extension”: The theorem uses the fact that if $`A𝔄(O_1)`$ where $`O_1O`$, then $`U(𝐚)^1AU(𝐚)𝔄(O)`$ for sufficiently small $`𝐚`$ in four independent directions. Since, however, the typical quantum field theory cannot be expected to admit a fixed-time ($`3+1`$) formulation (cf. Haag 1992, 59), it is not at all clear that they have truly avoided the Reeh-Schlieder theorem in any interesting sense. It would certainly be interesting to see which, *if any*, of the full Reeh-Schlieder theorem’s three premises would be rejected by a more general NW localization scheme.
However, we need not speculate about the possibility that the full Reeh-Schlieder theorem will apply to some generalization of NW localization scheme: The Reeh-Schlieder theorem already has counterintuitive consequences for the fixed-time NW localization scheme. In particular, although the vacuum $`\mathrm{\Omega }`$ is not cyclic under operations NW-localized in some spatial region $`G`$ at a single time, $`\mathrm{\Omega }`$ *is* cyclic under operators NW-localized in $`G`$ *within an arbitrary short time interval*. Before I give the precise version of this result, I should clarify some matters concerning the relationship between the dynamics of the field and local algebras.
In the standard localization scheme, the dynamics of local algebras may be thought of two ways. On the one hand, we may think of the assignment $`G(G)`$ as telling us, once and for all, which observables are associated with the region $`G`$, in which case the state of $`(G)`$ (i.e., the reduced state of the entire field) changes via the unitary evolution $`U(t)`$ (Schrödinger picture). On the other hand, we may think of the state of the field as fixed, in which case the algebra $`(G)`$ evolves over time to the algebra $`U(t)^1(G)U(t)`$ (Heisenberg picture). Thus, $`U(t)^1(G)U(t)`$ gives those operators classically localized in $`G`$ at time $`t`$. The Schrödinger picture is particularly intuitive in this case, since it mimics the dynamics of a classical field where quantities associated with points in space change their values over time.
Now, neither Segal nor Fleming explain how we should think of the dynamics of the NW-local algebras. Presumably, however, we are to think of the dynamics of the NW-local algebras in precisely the same way as we think of the dynamics of the standard local algebras.<sup>10</sup><sup>10</sup>10It is conceivable that Segal or Fleming have some different idea concerning the relationship between NW-local algebras at different times. For example, perhaps even in the Schrödinger picture, the map $`G_{NW}(G)`$ should be thought of as time-dependent. Although this is surely a formal possibility, it is exceedingly difficult to understand what it might mean, physically, to have a time-dependent association of physical magnitudes with regions in space. In particular, we may suppose that the state of the field is, at all times, the vacuum state $`\mathrm{\Omega }`$, and that $`U(t)^1_{NW}(G)U(t)`$ gives those operators NW-localized in $`G`$ at time $`t`$.
Now for any $`\mathrm{\Delta }`$ let
$$𝐒_\mathrm{\Delta }=\{U(t)^1AU(t):A_{NW}(G),t\mathrm{\Delta }\}.$$
(37)
That is, $`𝐒_\mathrm{\Delta }`$ consists of those operators NW-localized in $`G`$ at some time $`t\mathrm{\Delta }`$.
###### Theorem 2.
For any interval $`(a,b)`$ around $`0`$, $`\mathrm{\Omega }`$ is cyclic for $`𝐒_{(a,b)}`$.
###### Sketch of proof:.
Let $`[𝐒_{(a,b)}\mathrm{\Omega }]`$ denote the closed linear span of $`\{A\mathrm{\Omega }:A𝐒_{(a,b)}\}`$. Since the infinitesimal generator $`d\mathrm{\Gamma }(H)`$ of the group $`U(t)`$ is positive, Kadison’s “little Reeh-Schlieder theorem” (1970) entails that $`[𝐒_{(a,b)}\mathrm{\Omega }]=[𝐒_{}\mathrm{\Omega }]`$. However, $`[𝐒_{}\mathrm{\Omega }]=()`$; i.e., $`\mathrm{\Omega }`$ is cyclic under operators NW-localized in $`G`$ over all times (Segal 1964, 143). Therefore, $`\mathrm{\Omega }`$ is cyclic for $`𝐒_{(a,b)}`$. ∎
In Fleming’s language, then, the NW-local fields “allow the possibility of arbitrary space-like distant effects” from actions localized in an arbitrarily small region of space over an arbitrarily short period of time. Is this any less “counterintuitive” than the instantaneous version of the Reeh-Schlieder theorem for the standard localization scheme?<sup>11</sup><sup>11</sup>11One may, however, reject the interpretation of elements of $`_{NW}(G)`$ as operations that can be performed in $`G`$. I return to this point in the next section.
Finally, we are in a position to see explicitly a “counterintuitive” feature of the NW localization scheme that is not shared by the standard localization scheme: NW-local operators fail to commute at spacelike separation. For this, choose mutually disjoint regions $`G_1`$ and $`G_2`$ in $`^3`$, and choose an interval $`(a,b)`$ around $`0`$ so that $`O_1:=_{t(a,b)}(G_1+t)`$ and $`O_2:=_{t(a,b)}(G_2+t)`$ are spacelike separated. Let $`𝔄_{NW}(O_i)`$ be the von Neumann algebra generated by
$$\underset{t(a,b)}{}U(t)^1_{NW}(G_i)U(t).$$
(38)
Then it follows from Theorem 2 that the vacuum is cyclic for $`𝔄_{NW}(O_2)`$. However, since $`𝔄_{NW}(O_1)_{NW}(G)`$ contains annihilation operators and number operators, it follows that $`𝔄_{NW}(O_1)`$ and $`𝔄_{NW}(O_2)`$ do not satisfy microcausality. (Microcausality, in conjunction with cyclicity of the vacuum vector, would entail that the vacuum vector is separating.) More specifically, while the algebras $`U(t)^1_{NW}(G_1)U(t)`$ and $`U(t)^1_{NW}(G_2)U(t)`$ do satisfy microcausality for any fixed $`t`$, microcausality does not generally hold for the algebras $`U(t)^1_{NW}(G_1)U(t)`$ and $`U(s)^1_{NW}(G_2)U(s)`$ when $`ts`$ (despite the fact that $`G_1+t`$ and $`G_2+s`$ are spacelike separated).
It would be naive at this stage to claim that failure of generalized microcausality provides a simple reductio on the NW-localization scheme. As I will argue in the next section, however, the failure of generalized microcausality for the NW-local algebras leaves little room for making any physical sense of the NW localization scheme.
## 6 Local properties and local measurements
The assignment $`G(G)`$ was originally taken to have the operational meaning that $`(G)`$ consists of those observables that are *measurable* in the region $`G`$. What is the intended meaning of the alternative assignment $`G_{NW}(G)`$? Segal (1964, 142) presents the NW localization scheme as a contrasting claim about what can be measured in the spatial region $`G`$:
> “From an operational viewpoint it is these variables \[i.e., $`\mathrm{\Phi }(f)`$ with $`fL^2(G)`$\]…that appear as the localized field variables, and the ring $`_{NW}(G)`$…appears as the appropriate ring of local field observables, rather than the ring $`(G)`$…” (notation adapted)
According to this interpretation, the standard localization scheme and the NW localization scheme present us with two empirically *in*equivalent versions of quantum field theory. (For example, the vacuum displays Bell correlations relative to the algebras $`(G)`$ and $`(G^{})`$, while the vacuum is a product state across $`_{NW}(G)`$ and $`_{NW}(G^{})`$.) Thus, deciding which localization scheme is “correct” would be a matter of experiment, not a matter of interpretation.
There is, however, also a conceptual difficulty with interpreting $`_{NW}(G)`$ as the algebra of observables measurable in $`G`$. In particular, if $`[A,B]0`$, then measurement of $`A`$ can affect the statistics for outcomes of $`B`$ and vice versa. Thus, if $`_{NW}(G)`$ is what is measurable in $`G`$, then the failure of generalized microcausality for the NW-local algebras would pose the threat of causal anomalies.
This difficulty posed by the failure of generalized microcausality is clear to Fleming. In response, he and Butterfield (1999, 158) note that
> “…one naturally assumes that one can interpret the *association* of an operator with a spacetime region as implying that one can *measure it by performing operations confined* to that region,”
and they assert that they, “…question \[this\] interpretive assumption” (159). But if not local measurability, what does association of an observable with a spatial region mean? Fleming (2000, 21) claims that NW position operators,
> “…are more closely related than the local field coordinate to assessments of *where*, on hyperplanes and in space-time, objects, systems, their localizable properties and phenomena are located.”
It seems then that Fleming intends something along the lines of:
> $`()`$ The projections in $`_{NW}(G)`$ correspond to the *properties* of the system that are *localized in $`G`$.*
But what does Fleming mean by saying that a property is localized in a spatial region $`G`$? And why would the properties localized in $`G`$ differ from what can be measured in $`G`$?
Although Fleming has not offered a “philosophical account” of localized properties, he has provided analogies from classical mechanics in order to prime our intuitions about physical quantities that may “pertain to” a region, without being measurable in that region (cf. Fleming 2000; Fleming and Butterfield 1999). For example, take the center of mass $`C`$ of a spatially extended system. At a given time, $`C`$ is located at a point $`𝐱`$ in space, but $`C`$ is not measurable at $`𝐱`$ or even in spatial regions immediately surrounding $`𝐱`$. Perhaps then we can think of NW-localized quantities as similar to center of mass, center of charge, and their ilk.
This analogy, however, conceals an equivocation in the meaning of “localized.” To see this clearly, let me distinguish two types of localized quantities in classical mechanics. On the one hand, a physical quantity $`Q`$ may be permanently attached to some point $`𝐱`$ in space, in which case the values of $`Q`$ are given by vectors (or more generally, tensors) in the tangent space $`T_𝐱`$ over $`𝐱`$ (e.g., magnetic field strength at $`𝐱`$). I will refer to this first type of localized quantity as *fixedly-localized*. On the other hand, some quantities take vectors in *physical space* as their values (e.g., center of mass of a spatially extended system, or position of a point particle). I will refer to this latter type of localized quantity as *variably-localized*.
I grant that there is a sense in which v-localized quantities are “localized,” despite the fact that they may not be locally measurable (e.g., center of mass). However, analogies to v-localized quantities go no distance in clarifying the localization map $`G_{NW}(G)`$, since this is a permanent assignment (f-localization) of physical quantities to regions in space. What we really need, then, is an example of a f-localized quantity that is not locally measurable.
One might claim that examples of such quantities are readily forthcoming: Let $`C`$ be center of mass, and let $`Q`$ be the quantity that assumes the value $`1`$ if $`C=𝐱`$, and $`0`$ if $`C𝐱`$. Obviously, $`Q`$ will not typically be measurable in the vicinity of $`𝐱`$. But shouldn’t we say that $`Q`$ always “pertains to” $`𝐱`$, or is f-localized at $`𝐱`$?
The intuition behind thinking that $`Q`$ “pertains to” $`𝐱`$ seems to be based on the fact that $`Q`$ tells us something about $`𝐱`$, viz., whether it is $`C`$’s value. However, if this a sufficient condition for $`Q`$’s pertaining to $`𝐱`$, then $`Q`$ pertains to *every* point in space. Indeed, let $`𝐲`$ be another vector in $`^3`$ and introduce the new quantity $`C^{}=C+(𝐲𝐱)`$. Then $`Q=1`$ if and only if $`C^{}=𝐲`$, and so the previous line of reasoning would imply that $`Q`$ pertains to $`𝐲`$ (since $`Q`$ tells us whether $`𝐲`$ is the value of $`C^{}`$). What we should conclude, then, is that $`Q`$ is *not* f-localized at $`𝐱`$ in the same sense that a field quantity may be localized at $`𝐱`$. Thus, we have yet to find an example of a f-localized quantity that is not locally measurable.
In summary, while it is clear what it means for a physical quantity to be v-localized in $`G`$ but not measurable in $`G`$, it is by no means clear what it would mean for a physical quantity to be f-localized in $`G`$ but not measurable in $`G`$. As a result, it is also not clear how we should interpret the localization map $`G_{NW}(G)`$.
## 7 Conclusion
Introduction of the NW localization scheme into quantum field theory was an ingenious move. By means of one deft transformation, it appears to thwart the Reeh-Schlieder theorem and to restore the “intuitive” picture of localization from non-relativistic quantum mechanics. However, there are many reasons to doubt that Newton-Wigner has truly spared us of the counterintuitive consequences of the Reeh-Schlieder theorem. First, NW-local algebras still have a dense set of cyclic vectors. Second, since general quantum field theories cannnot be expected to admit a fixed-time formulation, it is not clear that the NW localization scheme has any interesting level of generality. Third, NW-local operations on the vacuum over an arbitrarily short period of time do generate the state space of the entire field. And, finally, the failure of generalized microcausality for the NW localization scheme leaves us without any natural physical interpretation of the correspondence $`G_{NW}(G)`$.
After showing that the Reeh-Schlieder theorem fails for NW-local algebras, Fleming (2000, 11) states that, “Now it is clear why it would be worthwhile to see the NW fields as covariant structures.” While there may be very good reasons for seeing the NW fields as covariant structures, avoiding the Reeh-Schlieder theorem is *not* one of them.
Acknowledgments: I am grateful for helpful correspondence from Jeremy Butterfield and Bernard Kay, and I am extremely grateful to Rob Clifton for his input throughout this project.
References:
Araki, Huzihiro (1964), “Types of von Neumann algebras of local observables for the free scalar field”, Progress of Theoretical Physics 32: 956-961.
Bratteli, Ola, and Derek Robinson (1997), Operator Algebras and Quantum Statistical Mechanics, Vol. 2. NY: Springer.
Clifton, Rob, David Feldman, Hans Halvorson, Michael Redhead, and Alex Wilce (1998), “Superentangled states”, Physical Review A 58: 135-145.
Clifton, Rob, and Hans Halvorson (2000), “Entanglement and open systems in algebraic quantum field theory”, Studies in the History and Philosophy of Modern Physics, forthcoming.
Fleming, Gordon (1996), “Just how radical is hyperplane dependence?”, in Rob Clifton (ed.), Perspectives on Quantum Reality. Dordrecht: Kluwer.
Fleming, Gordon (2000), “Reeh-Schlieder meets Newton-Wigner”, in PSA 1998, v. 2, forthcoming.
Fleming, Gordon, and Jeremy Butterfield (1999), “Strange positions”, in Jeremy Butterfield and Constantine Pagonis (eds.), From Physics to Philosophy, NY: Cambridge University Press.
Haag, Rudolf (1992), Local Quantum Physics. NY: Springer.
Halvorson, Hans, and Rob Clifton (2000), “Generic Bell correlation between arbitrary local algebras in quantum field theory”, Journal of Mathematical Physics 41: 1711-1717.
Horuzhy, Sergei (1988), Introduction to Algebraic Quantum Field Theory. Dordrecht: Kluwer.
Kadison, Richard V. (1970), “Some analytic methods in the theory of operator algebras”, in C. T. Taam (ed.), Lectures in Modern Analysis and Applications, Vol. II. NY: Springer, 8-29.
Kay, Bernard (1979), “A uniqueness result in the Segal-Weinless approach to linear Bose fields”, Journal of Mathematical Physics 20: 1712-1713.
Malament, David (1996), “In defense of dogma: Why there cannot be a relativistic quantum mechanics of (localizable) particles”, in Rob Clifton (ed.), Perspectives on Quantum Reality. Dordrecht: Kluwer, 1-10.
Petz, Dénes (1990), An Invitation to the Algebra of Canonical Commutation Relations. Leuven University Press.
Redhead, Michael (1995a), “More ado about nothing”, Foundations of Physics 25: 123-137.
Redhead, Michael (1995b), “The vacuum in relativistic quantum field theory”, in David Hull, Micky Forbes, and Richard M. Burian (eds.), PSA 1994, v. 2. East Lansing, MI: Philosophy of Science Association, 77-87.
Saunders, Simon (1992), “Locality, complex numbers, and relativistic quantum theory”, in David Hull, Micky Forbes, and Kathleen Okruhlik (eds.), PSA 1992, v. 1: East Lansing, MI: Philosophy of Science Association, 365-380.
Segal, Irving E. (1964), “Quantum fields and analysis in the solution manifolds of differential equations”, in William T. Martin and Irving E. Segal, (eds.), Proceedings of a Conference on the Theory and Applications of Analysis in Function Space. Cambridge: MIT Press, 129-153.
Segal, Irving E., and Roe W. Goodman (1965), “Anti-locality of certain Lorentz-invariant operators”, Journal of Mathematics and Mechanics 14: 629-638.
Summers, Stephen J., and Reinhard Werner (1985), “The vacuum violates Bell’s inequalities”, Physics Letters 110A: 257-259.
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# 1. Introduction
## 1. Introduction
Analysis of hadronization effects to the final states in $`\text{e}^+\text{e}^{}`$annihilation has became the subject of active QCD studies . There exist infrared and collinear safe event shape variables for which perturbative QCD can be applied at large center-of-mass energies $`s=Q^2`$ to calculate their differential distributions and mean values as series in $`\alpha _s(Q)`$. It has been observed many years ago that for some shape variables like thrust, $`t`$, and heavy jet mass, $`\rho `$, perturbative QCD predictions deviate from the data by corrections suppressed by powers of the large energy scale $`1/Q^p`$, with the exponent $`p`$ depending on the variable and $`p=1`$ for $`t`$ and $`\rho `$variables. Such hadronization corrections were measured experimentally over a wide energy interval $`14\sqrt{s}/\mathrm{GeV}189`$ and were found to have a different form for the differential event shape distributions, $`d\sigma /de`$, as compared to their mean values, $`e=\sigma _{\mathrm{tot}}^1𝑑ee𝑑\sigma /𝑑e`$. For the mean value $`e`$ the leading power correction is parameterized by a nonperturbative scale $`\lambda _p`$ of dimension $`p`$, while hadronization corrections to the differential distribution are described by a function $`f_{\mathrm{hadr}}(Q,e)`$ depending on both the shape variable and the center-of-mass energy
$$e=e_{_{\mathrm{PT}}}+\lambda _p/Q^p,\frac{1}{\sigma _{\mathrm{tot}}}\frac{d\sigma }{de}=\frac{d\sigma _{_{\mathrm{PT}}}}{de}+f_{\mathrm{hadr}}(Q,e)$$
(1.1)
with $`e`$ denoting a general event shape variable $`(e=t,\rho ,C,\mathrm{})`$ and the subscript PT referring to perturbative contribution, $`e_{_{\mathrm{PT}}}=𝑑ee𝑑\sigma _{_{\mathrm{PT}}}/𝑑e`$. Obviously, the hadronization corrections to the differential distributions have a richer structure then those to the mean values. For instance, nonperturbative scales $`\lambda _p`$ parameterizing power corrections to $`e`$ are defined by the moment $`𝑑eef_{\mathrm{hadr}}(Q,e)`$.
Power corrections in (1.1) are associated with hadronization effects in $`\mathrm{e}^+\mathrm{e}^{}`$final states and, as a consequence, the magnitude of the scales $`\lambda _p`$ and the function $`f_{\mathrm{hadr}}(Q,e)`$ cannot be calculated within perturbative QCD approach. However it was recognized some time ago , that analysis of infrared renormalon ambiguities of perturbative QCD series suggests the value of dimensionless exponents $`p`$ as well as the dependence of the function $`f_{\mathrm{hadr}}(Q,e)`$ on the large scale $`Q`$. Namely, perturbative QCD series generate power corrections of the form (1.1) through IR renormalons contribution but fail to predict uniquely their values – it is only the sum of perturbative and nonperturbative contributions that becomes well-defined . To give a meaning to the perturbative series in (1.1) one has to regularize IR renormalon singularities. This can be done in two different ways: one can specify a particular prescription for integrating IR renormalon singularities like principal value prescription . Alternatively, one can avoid IR renormalon ambiguities by introducing an explicit IR cut-off $`\mu `$ on momenta of soft particles in perturbative expressions. In this case, one can either impose a “hard” IR cut-off on momenta of soft particles in the Feynman integrals, $`k_{}>\mu `$, or replace QCD coupling constant by a effective IR finite coupling constant which coincides with $`\alpha _s(k_{})`$ at large scale $`k_{}`$ and deviates from it at $`k_{}<\mu `$ . Following each of these ways, one specifies perturbative ($`\mu `$dependent) contribution to (1.1) including perturbatively induced power corrections. Still, there exists a genuine nonperturbative contribution to the event shapes coming from the QCD dynamics at scales below $`\mu `$. This contribution cannot be determined from the analysis of perturbative QCD series while its magnitude depends on the choice of the IR regularization and, as a consequence, on the IR cut-off $`\mu `$.
For some hadronic observables like mean values of the event shapes and their differential distributions away from the end-point region, leading nonperturbative power corrections can be parameterized using different IR renormalon inspired phenomenological models . Their predictions agree well with the experimental data and the extracted values of phenomenological nonperturbative parameters exhibit approximate universality. Despite a phenomenological success of these models, it remains still unclear what is the physical meaning of new nonperturbative QCD scales and what is the origin of the universality property within QCD. In the present paper we address these problems using the factorization properties of the event shape distributions established in . We shall argue that nonperturbative power corrections to the thrust, heavy jet mass and $`C`$parameter distributions are described by the universal shape function which is a new nonperturbative QCD distribution measuring the energy flow in the two-jet final states in $`\text{e}^+\text{e}^{}`$annihilation.
The paper is organized as follows. In Sect. 2 we discuss the general properties of power corrections to the event shape distributions in the end-point region. In Sect. 3 we formulate the factorization procedure and define the shape function. In Sect. 4 we show that the differential event shape distributions are given by the convolution of the resummed perturbative cross-sections with universal shape function. Choosing a simple ansatz for this function we compare QCD predictions with the existing data. In Sect. 5 we apply the obtained expressions to calculate the power corrections to the first two moments of the distributions. Concluding remarks are given in Sect. 6.
## 2. Event shape distributions
In this paper we shall consider three event shape variables: thrust $`T`$, heavy jet mass $`\rho `$ and $`C`$parameter. They are defined in the standard way as
$$T=\underset{\stackrel{}{n}_T}{\mathrm{max}}\frac{_k|\stackrel{}{p}_k\stackrel{}{n}_T|}{_k|\stackrel{}{p}_k|},\rho =\mathrm{max}(\frac{M_R^2}{Q^2},\frac{M_L^2}{Q^2}),$$
(2.1)
where $`M_R^2`$ and $`M_L^2`$ denote the total invariant masses flowing into the right and left hemispheres with respect to the plane orthogonal to the thrust axis $`\stackrel{}{n}_T`$. The $`C`$parameter is given by
$$C=3(\theta _1\theta _2+\theta _2\theta _3+\theta _3\theta _1)$$
(2.2)
with $`\theta _j`$ being eigenvalues of space-like part of the energy-momentum tensor $`\mathrm{\Theta }^{\alpha \beta }=_kp_k^\alpha p_k^\beta /|p_k|/_j|p_j|`$.
Introducing the new variable $`t=1T`$ one notices that thus defined event shapes $`e=(t,\rho ,C)`$ have a number of common features. Lowest order perturbative QCD calculation leads in all three cases to the following expression for the differential distribution for $`e>0`$
$$\frac{d\sigma _{_{\mathrm{PT}}}}{de}=\frac{\alpha _s(Q)}{2\pi }A_e(e)\theta (e_{\mathrm{max}}e)+\left(\frac{\alpha _s(Q)}{2\pi }\right)^2B_e(e)+𝒪(\alpha _s^3),$$
(2.3)
where $`A_e`$ and $`B_e`$ are known coefficient functions and normalization is chosen as $`𝑑e\frac{d\sigma _{_{\mathrm{PT}}}}{de}=1`$. Lowest order correction $`A_e`$ gets contribution only from the three-particle final state which populates the kinematic region $`0ee_{\mathrm{max}}`$ with $`t_{\mathrm{max}}=\rho _{\mathrm{max}}=1/3`$ and $`C_{\mathrm{max}}=3/4`$. Away from the end-point region, $`e\mathrm{\Lambda }_{_{\mathrm{QCD}}}/Q`$, the perturbative expansion (2.3) is well-defined and it describes the final states consisting of particles with relative transverse momentum that scales at large center-of-mass energy as $`Q`$. As $`e`$ approaches the three-particle upper limit, $`e=e_{\mathrm{max}}`$, $`A_t`$ and $`A_\rho `$ vanish while $`A_C`$ takes a finite value
$`A_t(1/3)=A_\rho (1/3)=0,`$
$`A_C(C)={\displaystyle \frac{256}{243}}\pi \sqrt{3}C_F\left[1{\displaystyle \frac{8}{3}}\left(C{\displaystyle \frac{3}{4}}\right)+𝒪\left((C3/4)^2\right)\right].`$ (2.4)
For $`e=(t,\rho ,C)\mathrm{\Lambda }_{_{\mathrm{QCD}}}/Q`$ the final states consist of two narrow jets with invariant mass $`M_{R,L}^2\mathrm{\Lambda }_{_{\mathrm{QCD}}}Q`$. Examining (2.3) one finds that $`A_e`$ diverges in the end-point region $`e0`$ as
$$A_e(e)=\frac{4C_F}{e}\left[\mathrm{ln}\frac{e_0}{e}\frac{3}{4}\right]+𝒪(\mathrm{ln}e)$$
(2.5)
with $`t_0=\rho _0=1`$ and $`C_0=6`$. Similar Sudakov-type corrections appear to higher orders, $`\alpha _s^N\mathrm{ln}^{2Nn}e/e`$ with $`n0`$, and need to be resummed . They originate from the effects of collinear splitting of quarks and gluons inside two narrow energetic jets and their interaction with surrounding cloud of soft gluons. The underlying QCD dynamics depends on two infrared scales, $`Qe`$ and $`Q^2e`$, such that $`1/Q1/(Qe^{1/2})1/(Qe)`$. The smallest scale $`Qe`$ sets up the typical energy carried by soft gluons, while the scale $`Q\sqrt{e}`$ defines the transverse momenta of the jets, $`k_{}^2=Q^2e`$. Applying the standard IR renormalon analysis and examining sensitivity of perturbative expressions with respect to emission of particles on each of these scales, that is soft gluons with energy $`Qe`$ and collinear particles with the transverse momentum $`Q^2e`$, one finds that nonperturbative corrections to the differential distribution appear suppressed by powers of both scales. Then, in the end-point region, $`e=𝒪(\mathrm{\Lambda }_{_{\mathrm{QCD}}}/Q)`$, we may use the fact that $`Qe=𝒪(\mathrm{\Lambda }_{_{\mathrm{QCD}}})`$ and expand the differential distribution in powers of larger scale $`Q^2e`$. Keeping only the leading term of the expansion one gets
$$\frac{1}{\sigma _{\mathrm{tot}}}\frac{d\sigma }{de}=\sigma _0(\alpha _s(Q),\mathrm{ln}e,\frac{1}{Qe})+𝒪\left(\frac{1}{Q^2e}\right),$$
(2.6)
where $`\sigma _0`$ resums perturbative corrections in $`\alpha _s(Q)`$ as well as power corrections on the smallest scale $`Qe`$
$$\sigma _0=\frac{d\sigma _{_{\mathrm{PT}}}}{de}+\underset{k=1}{\overset{\mathrm{}}{}}\frac{\lambda _k}{(Qe)^k}\mathrm{\Sigma }_k(\alpha _s(Q),\mathrm{ln}e).$$
(2.7)
Here, $`\mathrm{\Sigma }_k`$ are dimensionless perturbative coefficient functions and $`\lambda _k`$ are some nonperturbative scales, depending, in general, on the choice of the event shape variable. Using (2.7) we notice that the power corrections have a different form for $`e\mathrm{\Lambda }_{_{\mathrm{QCD}}}/Q`$ and $`e\mathrm{\Lambda }_{_{\mathrm{QCD}}}/Q`$.
For $`e`$ away from the end-point region, $`e\mathrm{\Lambda }_{_{\mathrm{QCD}}}/Q`$, one may keep in (2.7) only the first term
$$\frac{1}{\sigma _{\mathrm{tot}}}\frac{d\sigma }{de}=\frac{d\sigma _{_{\mathrm{PT}}}}{de}+\frac{\lambda _1}{Qe}\mathrm{\Sigma }_1(\alpha _s(Q),\mathrm{ln}e)+𝒪\left(\frac{1}{(Qe)^2}\right).$$
(2.8)
The coefficient function $`\mathrm{\Sigma }_1`$ can be found using the well-known property that the leading $`1/Q`$power correction to the differential distribution (2.8) is generated by a shift of perturbative spectrum, $`ee\lambda _1/Q`$. This leads to
$$\mathrm{\Sigma }_1(\alpha _s(Q),\mathrm{ln}e)=\frac{d}{d\mathrm{ln}e}\left[\frac{d\sigma _{_{\mathrm{PT}}}}{de}\right].$$
(2.9)
Then, it follows from (2.8) that for $`e\mathrm{\Lambda }_{_{\mathrm{QCD}}}/Q`$ the leading power corrections to the differential distributions have a rather simple structure: they are parameterized by a single nonperturbative scale $`\lambda _1`$. The same scale determines $`1/Q`$power correction to the mean value $`e`$. The QCD predictions (2.8) are in a good agreement with the experimental data and the value of $`\lambda _1`$ has been fitted for different shape variables . It is worthwhile to note that in the performed analysis of power corrections to the differential distributions the fitting range of event shape variables was restricted to the region $`e\mathrm{\Lambda }_{_{\mathrm{QCD}}}/Q`$. At the same time, as we shall argue below, it is in the region $`e\mathrm{\Lambda }_{_{\mathrm{QCD}}}/Q`$ where a novel QCD regime is realized and the structure of hadronization corrections is drastically changed.
For $`e\mathrm{\Lambda }_{_{\mathrm{QCD}}}/Q`$ one finds that all terms in (2.7) become equally important and, therefore, need to be resummed to all orders in $`1/(Qe)`$. The resummation is based on the remarkable factorization properties of the differential distributions. As was shown in , the nonperturbative corrections to the leading asymptotic term $`\sigma _0`$ are factorized out into nonperturbative distribution function, the so-called shape function. The general factorized expression for differential distribution looks like
$$\frac{1}{\sigma _{\mathrm{tot}}}\frac{d\sigma }{de}=_0^{eQ}𝑑\epsilon f_e(\epsilon )\frac{d\sigma _{_{\mathrm{PT}}}(e\frac{\epsilon }{Q})}{de}+𝒪\left(\frac{1}{Q^2e}\right).$$
(2.10)
Its explicit form depends on the choice of the event shape variable $`e=(t,\rho ,C)`$ and will be given in Sect. 4 (see Eqs. (4.1), (4.2) and (4.3)). For $`e\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$ one can expand the r.h.s. of (2.10) in powers of $`1/Q`$ to reproduce the expansion (2.7) with
$$\lambda _n=𝑑\epsilon \epsilon ^nf(\epsilon ),\mathrm{\Sigma }_n(\alpha _s(Q),\mathrm{ln}e)=\frac{(e)^n}{n!}\frac{d^n}{de^n}\left[\frac{d\sigma _{_{\mathrm{PT}}}}{de}\right].$$
(2.11)
Expression (2.10) has a simple physical interpretation – nonperturbative corrections increase invariant masses of jets and effectively shift perturbative spectrum towards larger values of the shape variables with the weight given by nonperturbative distribution $`f_e(\epsilon )`$.
## 3. Factorization and Shape functions
Factorization relations (2.10) take a simple form for the radiator functions $`R(e)`$ defined as
$$R(e)=_0^e𝑑e^{}\frac{1}{\sigma _{\mathrm{tot}}}\frac{d\sigma }{de^{}}\theta (ee(N)).$$
(3.1)
Here, $`\mathrm{}`$ denotes averaging over all possible final states in $`\mathrm{e}^+\mathrm{e}^{}`$annihilation with the weight given by the differential distribution $`1/\sigma _{\mathrm{tot}}d\sigma /de`$ and $`e(N)`$ denotes the value of the event shape variable $`e`$ for a given final state $`|N`$. Calculating $`R(e)`$ in perturbation theory one finds
$$R_{_{\mathrm{PT}}}(e)=1\frac{\alpha _s(Q)}{2\pi }_e^{e_{\mathrm{max}}}𝑑e^{}A_e(e^{})+𝒪(\alpha _s^2(Q)).$$
(3.2)
Close to the two-jet region, $`e0`$, perturbative expressions for $`R(e)`$ involve Sudakov logs $`\alpha _s^N\mathrm{ln}^{2Nn}e`$ with $`n0`$. In the case of event shapes under consideration, $`e=(t,\rho ,C)`$, these corrections can be systematically resummed to next-to-leading logarithmic (NLL) order and matched into exact two-loop perturbative expressions (3.2) . To this accuracy the weights $`e(N)`$ can be expressed in terms of the total invariant masses $`M_R^2`$ and $`M_L^2`$ of two jets flowing into the right and left hemispheres, respectively. Moreover, the $`t`$ and $`C`$parameters depend only on the sum of two masses and the corresponding perturbative radiation functions can be expressed to the NLL approximation as
$`R_t^{_{\mathrm{PT}}}(e)`$ $`=`$ $`\theta (et(N))_{_{\mathrm{PT}}}=\theta \left(e{\displaystyle \frac{M_R^2+M_L^2}{Q^2}}\right)_{_{\mathrm{PT}}}`$ (3.3)
$`R_C^{_{\mathrm{PT}}}(e)`$ $`=`$ $`\theta (eC(N))_{_{\mathrm{PT}}}=\theta \left(e6{\displaystyle \frac{M_R^2+M_L^2}{Q^2}}\right)_{_{\mathrm{PT}}}=R_t^{_{\mathrm{PT}}}(e/6),`$
where the subscript PT indicates that the final states in $`\text{e}^+\text{e}^{}`$annihilation are generating by perturbative branching of outgoing quark and antiquark. The radiator function for the $`\rho `$parameter depends separately on the masses of two jets. Taking into account that perturbative evolution of two jets is independent on each other to the NLL approximation one gets
$$R_\rho ^{_{\mathrm{PT}}}(e)=\theta (e\rho (N))_{_{\mathrm{PT}}}=\theta \left(e\frac{M_R^2}{Q^2}\right)_{_{\mathrm{PT}}}\theta \left(e\frac{M_L^2}{Q^2}\right)_{_{\mathrm{PT}}}.$$
(3.4)
The perturbative expressions (3.3) and (3.4) are valid in the two-jet kinematical region $`\mathrm{\Lambda }_{_{\mathrm{QCD}}}/Qe<e_{\mathrm{max}}`$ except the end-point region $`e\mathrm{\Lambda }_{_{\mathrm{QCD}}}/Q`$, in which the energy of emitted soft particles scales as $`k_{}eQ\mathrm{\Lambda }_{_{\mathrm{QCD}}}`$ and perturbation theory is expected to fail.
Calculating the radiator functions $`R_e(e)`$ one has to combine together perturbative and nonperturbative corrections. In the case of inclusive distributions, like deep inelastic structure functions and Drell-Yan distributions, this can be achieved by applying the factorisation theorems. They allow to separate short-distance dynamics into perturbatively calculable coefficient functions and absorb large-distance corrections into universal nonperturbative distributions. Specific feature of the differential event-shape distributions is that they are not inclusive quantities but rather weighted cross-sections and, as a consequence, the standard methods are not applicable in this case.
It turns out that IR factorization still holds for the leading term $`\sigma _0`$ in the expansion of the event-distributions (2.6) in the end-point region $`e\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$. Its origin has a simple physical interpretation. In end-point region, the final state in $`\text{e}^+\text{e}^{}`$annihilation consists of two narrow jets surrounding by a cloud of soft gluons. Nonperturbative corrections $`1/(Q^2e)`$ and $`1/(Qe)`$ are associated with emission of collinear particles with the transverse momenta $`k_{}^2Q^2t`$ and soft particles on the energy scale $`k_{}Qe`$, respectively. Neglecting power corrections to (2.6) on a larger scale, $`1/(Q^2e)`$, we may restrict analysis to soft particles only. Since soft particle cannot resolve the internal structure of narrow jets of transverse size $`k_{}^2Q^2e`$, we may effectively replace two jets by a pair of energetic quark and antiquark moving back-to-back with the energy $`Q/2`$. The internal dynamics of two jets is governed by perturbative branching of quark and antiquark while effects of their interaction with soft gluons can be factorized out into the eikonal phase $`W_+W_{}^{}`$ with $`W_+`$ and $`W_{}`$ being the eikonal phases of quark and antiquark, respectively. They are given by Wilson lines $`W_\pm =P\mathrm{exp}(i_0^{\mathrm{}}𝑑sn_\pm A(sn_\pm ))`$ in which soft gluon field $`A_\mu (x)`$ is integrated along the light-like directions $`n_\pm `$ defined by the momenta of two outgoing jets. In the end-point region, collinear and soft particles provide additive contributions to the shape variables $`e=(t,\rho ,C)`$. As a consequence, the radiator functions are given in all three cases by a convolution of perturbative radiators $`R_{_{\mathrm{PT}}}`$ and the same universal nonperturbative distribution $`f(\epsilon _R,\epsilon _L)`$ describing the energy flow into the right and left hemispheres in the final state, $`\epsilon _R`$ and $`\epsilon _L`$, respectively, created by nonperturbative soft gluon radiation. The nonperturbative distribution $`f(\epsilon _R,\epsilon _L)`$ is defined as follows
$$f(\epsilon _R,\epsilon _L)=\underset{N}{}|0|W_+W_{}^{}|N|^2\delta (\epsilon __R(k__Rn_+))\delta (\epsilon __L(k__Ln_{})).$$
(3.5)
Here, sum goes over all possible soft gluon final states $`|N`$ with $`k_R`$ and $`k_L`$ being the total momentum of soft particles moving into right and left hemispheres, respectively. The quantities $`(k__Rn_+)`$ and $`(k__Ln_{})`$ define the projection of the soft gluon momenta onto the directions of two jets, $`n_\pm ^\mu =(1,\mathrm{𝟎}_{},\pm 1)`$, propagating into the same hemisphere.
Finally, the factorized expressions for the radiator function for the $`t`$ and $`C`$variables look like
$`R_t(e)`$ $`=`$ $`{\displaystyle _0^{eQ}}𝑑\epsilon f_t(\epsilon )R_t^{_{\mathrm{PT}}}\left(e{\displaystyle \frac{\epsilon }{Q}}\right)`$ (3.6)
$`R_C(e)`$ $`=`$ $`{\displaystyle _0^{\frac{2}{3\pi }eQ}}𝑑\epsilon f_t(\epsilon )R_C^{_{\mathrm{PT}}}\left(e{\displaystyle \frac{3\pi }{2}}{\displaystyle \frac{\epsilon }{Q}}\right)`$ (3.7)
with nonperturbative distribution $`f_t(\epsilon )`$ defined as
$$f_t(\epsilon )=𝑑\epsilon _R𝑑\epsilon _Lf(\epsilon _R,\epsilon _L)\delta (\epsilon \epsilon _R\epsilon _L)=_0^\epsilon 𝑑\epsilon ^{}f(\epsilon \epsilon ^{},\epsilon ^{}).$$
(3.8)
In the case of the $`\rho `$variable,
$$R_\rho (e)=_0^{eQ}𝑑\epsilon _R_0^{eQ}𝑑\epsilon _Lf(\epsilon _R,\epsilon _L)R_J^{_{\mathrm{PT}}}\left(e\frac{\epsilon _R}{Q}\right)R_J^{_{\mathrm{PT}}}\left(e\frac{\epsilon _L}{Q}\right)$$
(3.9)
with $`R_\rho ^{_{\mathrm{PT}}}(e)=[R_J^{_{\mathrm{PT}}}(e)]^2`$ and $`R_J^{_{\mathrm{PT}}}(e)=\theta \left(eM_R^2/Q^2\right)_{_{\mathrm{PT}}}`$ being a single jet radiator function. We would like to stress that Eqs. (3.6), (3.7) and (3.9) hold in the region $`\mathrm{\Lambda }_{\mathrm{QCD}}^2/Q^2<e<e_{\mathrm{max}}`$. They resum all power corrections of the form $`1/(Qe)^n`$ and are valid up to corrections $`1/(Q^2e)`$. According to (3.6), (3.7) and (3.9), the power corrections have a different form for $`\rho `$ and $`e=(t,C)`$ variables. In the latter case, the radiator function depends on an overall energy flowing into both hemispheres and described by the integrated distribution (3.8).
Nonperturbative corrections to the radiator functions (3.6), (3.7) and (3.9) are governed by the universal shape function $`f(\epsilon _R,\epsilon _L)`$. This function is different from the well-known inclusive QCD distributions and its operator definition was given in . Using (3.5) it is straightforward to show that $`f(\epsilon _R,\epsilon _L)`$ is a symmetric function of its arguments, it does not depend on the center-of-mass energy $`Q`$ and is normalized as
$$\frac{d}{dQ^2}f(\epsilon _R,\epsilon _L)=0,𝑑\epsilon _R𝑑\epsilon _Lf(\epsilon _R,\epsilon _L)=1,$$
(3.10)
where the last relation follows from unitarity of the eikonal phase $`W_+W_{}^{}`$. The matrix element entering (3.5) does not depend on any kinematical scale and, as a consequence, the momenta of soft gluons contributing to (3.5) are not restricted from above. To separate the region of small gluon momenta one has to introduce the factorisation scale $`\mu `$. Then, the shape function describes the contribution of gluons with $`k_{}<\mu `$, while the contribution of gluons with $`k_{}>\mu `$, is absorbed into perturbative radiator function $`R(e)`$. In this way, both nonperturbative shape function and perturbative radiator become $`\mu `$dependent while this dependence cancel in their convolution (3.6), (3.7) and (3.9). Since the $`\mu `$dependence of radiator function $`R_{_{\mathrm{PT}}}`$ can be calculated perturbatively, the above condition allows to obtain the evolution equations on the nonperturbative distributions . Clearly, there exists an ambiguity in implementing IR cut-off inside perturbative expressions. Different prescriptions correspond to different ways of regularizing IR renormalon singularities and therefore lead to the different expressions for the nonperturbative distributions. In what follows we shall impose a “hard” IR cut-off on gluon momenta inside the perturbative radiator functions entering (3.6), (3.7) and (3.9) as
$$R_{_{\mathrm{PT}}}(e)R_{_{\mathrm{PT}}}(e;\mu )=\theta \left(e\frac{\mu }{Q}\right)R_{_{\mathrm{PT}}}^{^{\mathrm{NLL}}}(e)+\theta \left(\frac{\mu }{Q}e\right)R_{_{\mathrm{PT}}}^{^{\mathrm{NLL}}}(\mu /Q).$$
(3.11)
Throughout the paper we shall substitute $`R_{_{\mathrm{PT}}}^{^{\mathrm{NLL}}}(e)`$ by its perturbative expression resummed to the NLL accuracy and matched into two-loop explicit expressions within the modified $`\mathrm{ln}R`$matching scheme . Thus defined radiator function (3.11) depends on two scales, $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ and IR cut-off $`\mu `$, that we choose as
$$\mathrm{\Lambda }_{\mathrm{QCD}}=\mu =0.25\mathrm{GeV}.$$
(3.12)
Within the prescription (3.11), the “regularized” perturbative spectrum $`d\sigma _{_{\mathrm{PT}}}(e;\mu )de=dR_{_{\mathrm{PT}}}(e;\mu )/de`$ coincides with the $`\mathrm{ln}R`$matched perturbative distribution $`dR_{_{\mathrm{PT}}}^{^{\mathrm{NLL}}}(e)/de`$ for $`\mu /Q<e<e_{\mathrm{max}}`$ and it vanishes inside the nonperturbative “window” $`0<e<\mu /Q`$. Choosing the value of $`\mu `$ in (3.12) one has to be sure that the end-point of the perturbative distribution, $`e=\mu /Q`$, belongs to applicability range of the NLL resummed radiator function $`R_{_{\mathrm{PT}}}^{^{\mathrm{NLL}}}(e)`$ , $`2\beta _0\alpha _s(Q^2)\mathrm{ln}e<1`$. Despite the fact that the perturbative spectrum is well defined at $`e=\mu /Q`$ we do not expect that it provides a reasonable description of the physical distribution in the end-point region. Indeed, it is in this region that nonperturbative power corrections become dominant.
## 4. Differential distributions
Differentiating the radiator functions (3.6) and (3.7) we obtain the following expressions for the differential $`t`$distribution
$$\frac{1}{\sigma _{\mathrm{tot}}}\frac{d\sigma _t}{de}=Qf(Qe;\mu )R_t^{_{\mathrm{PT}}}(0;\mu )+_0^{Qe}𝑑\epsilon f_t(\epsilon ;\mu )\frac{d\sigma _t^{_{\mathrm{PT}}}(e\epsilon /Q;\mu )}{de}$$
(4.1)
and $`C`$distribution
$$\frac{1}{\sigma _{\mathrm{tot}}}\frac{d\sigma _C}{de}=\frac{2}{3\pi }Qf(\frac{3\pi }{2}Qe;\mu )R_C^{_{\mathrm{PT}}}(0;\mu )+_0^{\frac{2}{3\pi }Qe}𝑑\epsilon f_t(\epsilon ;\mu )\frac{d\sigma _C^{_{\mathrm{PT}}}(e\frac{3\pi }{2}\frac{\epsilon }{Q};\mu )}{de}.$$
(4.2)
Here, we indicated explicitly the dependence of nonperturbative shape function and perturbative distributions on the factorization scale $`\mu `$. Two terms entering the r.h.s. of (4.1) and (4.2) have the following interpretation. Since the shape function $`f_t(\epsilon )`$ rapidly vanishes for large $`\epsilon `$, the first term contributes inside the nonperturbative window $`0e<\mu /Q`$. In this region the emission of perturbative real soft gluons is suppressed due to cut-off imposed on soft gluon momenta $`k_{}>\mu `$ and the shape of the distribution is governed entirely by nonperturbative function $`f_t(\epsilon )`$. Additional Sudakov factor $`R^{_{\mathrm{PT}}}(0;\mu )`$ takes into account the contribution of virtual soft gluons with $`\mu <k_{}<Q`$ and it rapidly vanishes as $`\mu `$ decreases. The second term in (4.1) and (4.2) defines the spectrum inside the perturbative window $`\mu /Q<e<e_{\mathrm{max}}`$. In this region, nonperturbative corrections smear the perturbative spectrum over the interval $`\mathrm{\Delta }e\mathrm{\Lambda }_{_{\mathrm{QCD}}}/Q`$.
For the heavy mass distribution one finds
$$\frac{1}{\sigma _{_{\mathrm{tot}}}}\frac{d\sigma _\rho }{de}=Qf_\rho (eQ,eQ;\mu )R_J^{_{\mathrm{PT}}}(0;\mu )+_0^{eQ}𝑑\epsilon f_\rho (\epsilon ,eQ;\mu )\frac{d\sigma _J^{_{\mathrm{PT}}}(e\epsilon /Q;\mu )}{de},$$
(4.3)
where $`d\sigma _J^{_{\mathrm{PT}}}/de`$ is single jet distribution resummed to the NLL order and defined by the radiator function (3.9), $`d\sigma _J^{_{\mathrm{PT}}}/de=dR_J^{_{\mathrm{PT}}}(e)/de`$. The heavy mass nonperturbative distribution is given by
$$f_\rho (\epsilon ,eQ)=2_0^{eQ}𝑑\epsilon ^{}f(\epsilon ,\epsilon ^{})R_J^{_{\mathrm{PT}}}\left(e\frac{\epsilon ^{}}{Q}\right).$$
(4.4)
Comparing (4.3) with (4.1) and (4.2) we notice that the factorized expressions for the differential $`t`$, $`C`$ and $`\rho `$distributions have a similar form but the structure of power corrections is different in the case of the heavy mass. In distinction with (3.8), the heavy mass nonperturbative function $`f_\rho `$ depends on the shape variable, $`e`$, and the center-of-mass energy, $`Q`$. This dependence is controlled by perturbative radiator function and has the following interpretation. In the two-jet limit, the invariant mass of each jet is given by the sum of perturbative and nonperturbative contributions, $`M_R^2=M_{R,\mathrm{PT}}^2+\epsilon _RQ`$ and $`M_L^2=M_{L,\mathrm{PT}}^2+\epsilon _LQ`$. Perturbative radiation leads to $`M_{R,\mathrm{PT}}^2/Q^2M_{L,\mathrm{PT}}^2/Q^2𝒪(\alpha _s(Q))`$, while nonperturbative contribution scales as $`\epsilon _R\epsilon _L𝒪(\mathrm{\Lambda }_{\mathrm{QCD}})`$. In contrast with the $`t`$variable, which depends on the sum of both masses and therefore is additive with respect to perturbative and nonperturbative contributions, the $`\rho `$parameter is defined by the largest mass for which the “additivity” property is lost. Namely, comparing invariant masses flowing into two hemispheres one encounters a situation when masses of two perturbative jets are of the same order, $`M_{R/L,_{\mathrm{PT}}}^2=𝒪(Q^2)`$, while their difference is much smaller $`|M_{L,\mathrm{PT}}^2M_{R,\mathrm{PT}}^2|=𝒪(Q\mathrm{\Lambda }_{\mathrm{QCD}})`$.<sup>2</sup><sup>2</sup>2To see that this configuration is not rare it is enough to notice that it corresponds to the vicinity of peak of the perturbative distribution over the difference of the jet masses $`|M_L^2M_R^2|/Q^2`$ . In this case, nonperturbative correction to the difference of the jet masses becomes comparable with the perturbative contribution $`|M_{L,\mathrm{PT}}^2M_{R,\mathrm{PT}}^2|Q|\epsilon _L\epsilon _R|`$, and therefore it can invert the perturbative hierarchy of jet masses, $`M_{R,\mathrm{PT}}^2<M_{L,\mathrm{PT}}^2`$, into $`M_R^2>M_L^2`$, for instance. Expression (4.3) takes into account this effect through the induced $`Q`$dependence of the nonperturbative function (4.4).
According to (4.1), (4.2) and (4.3) the nonperturbative corrections to the $`t`$, $`C`$ and $`\rho `$distributions involve two different nonperturbative functions. They are related however to the same universal nonperturbative shape function (3.5) describing the energy flow into two hemispheres in the final state. We recall, that $`f(\epsilon _R,\epsilon _L;\mu )`$ depends on the cut-off $`\mu `$ imposed on the maximal momenta of soft particles but it is independent on the center-of-mass energy $`Q`$. By the definition, $`f(\epsilon _R,\epsilon _L;\mu )`$ distinguishes between particles propagating into right and left hemispheres in the final state and therefore it is not completely inclusive with respect to partonic final states. Namely, it takes into account that quarks and gluons produced at short distances $`1/Q`$ and moving into one of the hemispheres will eventually decay at large distances $`1/\mathrm{\Lambda }_{\mathrm{QCD}}`$ and their remnants could flow into opposite hemispheres.<sup>3</sup><sup>3</sup>3Similar effect has been studied using the IR renormalon approach in . This implies that, firstly, in contrast with the well-known inclusive QCD distributions, the shape function $`f(\epsilon _R,\epsilon _L;\mu )`$ is not related to the short distance QCD dynamics and, in particular, its moments can not be related to hadronic matrix elements of local composite operators. Indeed, according to the operator definition proposed in , the shape functions are defined in terms of the so-called “maximally nonlocal” QCD operators . Secondly, non-inclusive corrections to the shape function describe a “cross-talk” between two hemispheres in the final state leading to correlations between $`\epsilon _R`$ and $`\epsilon _L`$. As a consequence, the shape function is not factorizable into the product of functions depending on the energy flowing into separate hemispheres
$$f(\epsilon _R,\epsilon _L)=f_{\mathrm{incl}}(\epsilon _R)f_{\mathrm{incl}}(\epsilon _L)+\delta f_{\mathrm{non}\mathrm{incl}}(\epsilon _R,\epsilon _L).$$
(4.5)
One should notice that similar property holds for perturbative Sudakov resummed radiator function (3.4). However, one finds that there the factorization holds to the NLL accuracy, Eq. (3.4), and non-inclusive corrections first appear at the NNLL level $`\alpha _s^2(\alpha _sL)^N`$.
In what follows we shall rely on a particular ansatz for the shape function $`f(\epsilon _R,\epsilon _L)`$ which agrees with general properties of nonperturbative QCD distributions and has been used in previous studies of power corrections to the thrust distributions . Namely, one expects that for small values of $`\epsilon _R`$ and $`\epsilon _L`$ the shape function should vanish as a power of the energy. Similarly, $`f(\epsilon _R,\epsilon _L)`$ should rapidly vanish as $`\epsilon _R`$ or $`\epsilon _L`$ becomes large. Taking into account these properties together with (4.5) one chooses the following expression
$$f(\epsilon _R,\epsilon _L)=\frac{𝒩(a,b)}{\mathrm{\Lambda }^2}\left(\frac{\epsilon _R\epsilon _L}{\mathrm{\Lambda }^2}\right)^{a1}\mathrm{exp}\left(\frac{\epsilon _R^2+\epsilon _L^2+2b\epsilon _R\epsilon _L}{\mathrm{\Lambda }^2}\right).$$
(4.6)
It depends on two dimensionless parameters $`a`$ and $`b`$ and the scale $`\mathrm{\Lambda }`$. The factor $`𝒩(a,b)`$ is fixed by normalization condition (3.10).
The free parameters, $`a`$, $`b`$ and $`\mathrm{\Lambda }`$, have the following meaning. The exponent $`a`$ determines how fast the shape function vanishes at the origin. The scale $`\mathrm{\Lambda }`$ sets up the typical energy of soft radiation. The parameter $`b`$ controls the non-inclusive contribution to the shape function and its possible values are restricted as $`b>1`$ in order for the shape function (4.6) to be normalizable. Non-inclusive corrections vanish at $`b=0`$, $`\delta f_{\mathrm{non}\mathrm{incl}}=0`$ in (4.5). For $`b1`$, the shape function enhances the regions of the phase space $`\epsilon _R\epsilon _L`$ and $`\epsilon _R\epsilon _L`$, in which most of the energy flows into one of the hemispheres. For $`b1`$ the energies are of the same order, $`\epsilon _R\epsilon _L`$, and strongly correlated to each other. We expect that non-inclusive corrections to the shape function should be important and the configurations in which energy flows mostly into one of the hemispheres to be suppressed. This suggests that the possible values of the $`b`$parameter should lie within the interval $`1<b<0`$.
The parameters $`a`$, $`b`$ and $`\mathrm{\Lambda }`$ depend on the factorization scale $`\mu `$ and are independent on the center-of-mass energy $`Q`$ as well as the choice of the shape variable $`e=(t,\rho ,C)`$. This allows to fit their values by comparing the event shape distributions, Eqs. (4.1), (4.2) and (4.3), with the most precise available experimental data at $`Q=M_Z`$. Following this procedure we found that the fit to the heavy jet mass distribution is more sensitive to the choice of the parameters (especially to the non-inclusiveness parameter $`b`$) then the one to the thrust and the $`C`$parameter. Then, fitting the heavy jet mass distribution at $`Q=M_Z`$ as shown in Fig. 1a we obtain
$$a=2,b=0.4,\mathrm{\Lambda }=0.55\mathrm{GeV}.$$
(4.7)
Using these values we compare the QCD predictions for the $`C`$parameter distribution at $`Q=M_Z`$ with and without nonperturbative corrections included as shown in Fig. 1b. Similar plot for the thrust distribution can be found in . We observe that the differential distributions (4.3) and (4.2) combined with the shape function, Eqs. (4.6) and (4.7), correctly describe the data throughout the interval $`0<e<e_{\mathrm{max}}`$ including the end-point region $`e=𝒪(\mathrm{\Lambda }_{\mathrm{QCD}}/Q)`$. In addition, the $`\rho `$parameter distribution turns out to be very sensitive to the choice of the $`b`$parameter. The fact that its value, (4.7), is relatively large indicates that non-inclusive corrections to the shape function (4.5) are important indeed.
Having determined the parameters of the shape function, Eq. (4.7), at the reference energy scale $`Q=M_Z`$, we can now apply the factorized expressions for the differential distributions, (4.1), (4.2) and (4.3) with the same ansatz for the shape function (4.6) to obtain the QCD predictions at different energy and compare them with the data. The combined plot for the $`\rho `$ and $`C`$parameter distributions over the center-of-mass energy interval $`35\mathrm{GeV}Q189\mathrm{GeV}`$ is shown in Fig. 2 a and b, respectively. Similar plot for the thrust distribution can be found in . We observe that the theoretical curves reproduce the data over the whole interval of the shape variables including the end-point region.
## 5. Moments of the event shapes
Recently, the experimental data for the first few moments of various event shape distributions became available . Their analysis indicates a presence of large hadronization corrections whose form deviates from IR renormalon models describing the nonperturbative corrections to the distributions as the shift of perturbative spectrum.
Let us apply the obtained expressions for differential distributions to calculate the first two moments of the $`t`$, $`C`$ and $`\rho `$distributions defined as
$$e^n=_0^{e_{\mathrm{max}}}𝑑ee^n\frac{1}{\sigma _{\mathrm{tot}}}\frac{d\sigma }{de},(n=1,2).$$
(5.1)
Here, integration goes only over the part of the available phase space, $`0<e<e_{\mathrm{max}}`$, corresponding to the three-particle final states, and it does not take into account the contribution of multi-jet final states, $`e>e_{\mathrm{max}}`$. Quantitative description of hadronization corrections to such final states is not available yet. Putting an upper limit on the value of the shape variable in (5.1) allows us to avoid the latter contribution and to replace the differential distribution $`d\sigma /de`$ in (5.1) by the obtained expressions (4.1), (4.2) and (4.3) which are valid for $`0<e<e_{\mathrm{max}}`$.
Using general expression (2.10) one calculates the mean value of the event shape as
$$e=e_{_{\mathrm{PT}}}+\frac{\epsilon }{Q}\left[1e_{\mathrm{max}}\frac{d\sigma _{_{\mathrm{PT}}}(e_{\mathrm{max}})}{de}\right]+𝒪\left(\frac{1}{Q^2}\right),$$
(5.2)
where $`\mathrm{}_{_{\mathrm{PT}}}=_0^{e_{\mathrm{max}}}𝑑e(\mathrm{})𝑑\sigma _{_{\mathrm{PT}}}/𝑑e`$ denotes averaging with respect to perturbative distribution and the scale $`\epsilon `$ is defined as the first moment of the shape function, $`\epsilon =𝑑\epsilon \epsilon f(\epsilon )`$. It is important to remember that the factorized expressions for the differential distributions (4.1), (4.2) and (4.3) are valid up to $`𝒪(1/(Q^2e))`$corrections which may modify the mean value $`e`$ by $`𝒪(1/Q^2)`$terms. The additional factor in front of $`\epsilon /Q`$ takes into account that close to the edge of the three-particle phase space, $`ee_{\mathrm{max}}`$, the perturbative distribution (2.3) could take nonzero values
$$e_{\mathrm{max}}\frac{d\sigma _{_{\mathrm{PT}}}(e_{\mathrm{max}})}{de}=\frac{\alpha _s(Q)}{2\pi }e_{\mathrm{max}}A_e(e_{\mathrm{max}})+𝒪(\alpha _s^2).$$
(5.3)
This can be checked using the explicit expressions for perturbative distributions (2.4). We find that $`d\sigma _{_{\mathrm{PT}}}/de`$ vanishes to one-loop order as $`ee_{\mathrm{max}}`$ for the $`t`$ and $`\rho `$variables while for the $`C`$parameter it approaches a finite value. Finally, calculating the mean values $`\epsilon `$ with respect to nonperturbative distributions $`f_t(\epsilon )`$ and $`f_\rho (\epsilon ,eQ)`$ defined in (3.8) and (4.4), respectively, we obtain
$`t`$ $`=`$ $`t_{_{\mathrm{PT}}}+{\displaystyle \frac{\lambda _1}{Q}}+𝒪(1/Q^2),`$ (5.4)
$`\rho `$ $`=`$ $`\rho _{_{\mathrm{PT}}}+{\displaystyle \frac{\lambda _1}{2Q}}+𝒪(1/Q^2).`$ (5.5)
Similarly, for the mean value of the $`C`$parameter we get
$$C=C_{_{\mathrm{PT}}}+\frac{3\pi }{2}\frac{\lambda _1}{Q}\left[1\frac{\alpha _s(Q)}{2\pi }5.73+𝒪(\alpha _s^2)\right]+𝒪(1/Q^2).$$
(5.6)
Here, large perturbative coefficient originates from (2.4) and it reduces a magnitude of the nonperturbative scale $`\lambda _1`$ by $`11\%`$ at $`Q=M_Z`$. Relations (5.4) and (5.5) coincide with IR renormalon model predictions , while (5.6) differs by perturbative $`\alpha _s(Q)`$dependent “boundary” term. Nonperturbative $`Q`$independent scale $`\lambda _1`$ is given by
$$\lambda _1=𝑑\epsilon _R𝑑\epsilon _L(\epsilon _R+\epsilon _L)f(\epsilon _R,\epsilon _L)=𝑑\epsilon \epsilon f_t(\epsilon ).$$
(5.7)
Substituting expression for the shape function, Eq. (4.6), one finds $`\lambda _1=\mathrm{\Lambda }\phi (a,b)`$ with $`\phi `$ given by $`{}_{2}{}^{}F_{1}^{}`$hypergeometric series. Using the values of the parameters (4.7) we find
$$\lambda _1=1.22\mathrm{GeV}.$$
(5.8)
We would like to recall that this value depends on the factorization scale $`\mu `$, Eq. (3.12), and its $`\mu `$dependence is described by QCD evolution equation . Obviously, the value of $`\lambda _1`$, and as a consequence $`1/Q`$corrections to the mean values (5.4), (5.5) and (5.6) are less sensitive to the choice of the parameters $`a`$, $`b`$ and $`\mathrm{\Lambda }`$ as compared with nonperturbative corrections to the corresponding differential distributions.
The comparison of the QCD predictions, (5.4), (5.5) and (5.6), with the data over the energy interval $`35\mathrm{GeV}Q189\mathrm{GeV}`$ is shown in Fig. 3. One should notice that (5.4), (5.5) and (5.6) describe the contribution of the two-jet configurations while experimental data take into account all possible final states. A good agreement observed in Fig. 3 indicates that the contribution to the mean values of the final states with three and more jets as well as $`𝒪(1/(Q^2e))`$ subleading corrections to the distributions are subdominant. Indeed, the dominant contribution to the moments (5.1) comes from the vicinity of peak of the differential distribution $`e^nd\sigma /de`$. Using existing experimental data one can show that for $`n=1`$ and $`n=2`$ the position of the peak is located in the two-jet kinematical region while for higher $`n`$ it moves towards larger $`e`$ for which the integral (5.1) becomes very sensitive to the choice of the upper integration limit $`e_{\mathrm{max}}`$. We shall use this observation calculating the second moment of the event shape distributions.
Let us apply (5.1) to calculate $`e^2`$. Neglecting $`1/(Q^2e)`$corrections to the distribution (2.10) we find after some algebra the following general expression
$`e^2`$ $`=`$ $`e^2_{_{\mathrm{PT}}}+{\displaystyle \frac{\epsilon }{Q}}\left[2e_{_{\mathrm{PT}}}e_{\mathrm{max}}^2{\displaystyle \frac{d\sigma _{_{\mathrm{PT}}}(e_{\mathrm{max}})}{de}}\right]+`$ (5.9)
$`+`$ $`{\displaystyle \frac{\epsilon ^2}{Q^2}}\left[1e_{\mathrm{max}}{\displaystyle \frac{d\sigma _{_{\mathrm{PT}}}(e_{\mathrm{max}})}{de}}+{\displaystyle \frac{1}{2}}e_{\mathrm{max}}^2{\displaystyle \frac{d^2\sigma _{_{\mathrm{PT}}}(e_{\mathrm{max}})}{de^2}}\right]+𝒪(1/Q^3).`$
Similar to the mean value, (5.2), the boundary terms vanish for the $`t`$ and $`\rho `$variables while for the $`C`$parameter they provide a sizeable contribution. Using the explicit expression for the shape functions, (3.8) and (4.4), we calculate the scales $`\epsilon ^2`$ and take into account the boundary terms (2.4) to obtain<sup>4</sup><sup>4</sup>4We are grateful to O. Biebel and S. Kluth for providing us $`𝒪(\alpha _s^2)`$expressions for the moments of the event shapes.
$`t^2`$ $`=`$ $`t^2_{_{\mathrm{PT}}}+2{\displaystyle \frac{\lambda _1}{Q}}t_{_{\mathrm{PT}}}+{\displaystyle \frac{\lambda _2}{Q^2}}`$
$`\rho ^2`$ $`=`$ $`\rho ^2_{_{\mathrm{PT}}}+{\displaystyle \frac{\lambda _1}{Q}}\rho _{_{\mathrm{PT}}}+{\displaystyle \frac{\lambda _2+\delta \lambda _2(Q)}{4Q^2}}`$ (5.10)
$`C^2`$ $`=`$ $`C^2_{_{\mathrm{PT}}}+{\displaystyle \frac{3\pi }{2}}{\displaystyle \frac{\lambda _1}{Q}}\left[2C_{_{\mathrm{PT}}}{\displaystyle \frac{\alpha _s(Q)}{2\pi }}4.30\right]+{\displaystyle \frac{9\pi ^2}{4}}{\displaystyle \frac{\lambda _2}{Q^2}}\left[1{\displaystyle \frac{\alpha _s(Q)}{2\pi }}11.46\right].`$
Here, the scale $`\lambda _1`$ was defined in (5.7) and new scales $`\lambda _2`$ and $`\delta \lambda _2`$ are given by
$$\lambda _2=(\epsilon _R+\epsilon _L)^2,\delta \lambda _2(Q)=\left(\epsilon _R\epsilon _L\right)^2\left\{1+4_0^{\rho _{\mathrm{max}}}𝑑\rho ^{}\rho ^{}\left(\frac{d\sigma _J^{\mathrm{PT}}}{d\rho ^{}}\right)^2\right\},$$
(5.11)
where average is taken with respect to the shape function $`f(\epsilon _R,\epsilon _L)`$. The $`Q`$dependence of the scale $`\delta \lambda _2`$ is attributed to perturbative prefactor depending on the single jet distribution, $`d\sigma _J^{\mathrm{PT}}/d\rho `$, defined in (4.3). Its origin was explained in Sect. 4. We find that the value of this factor varies from $`2.19`$ at $`Q=10\mathrm{Gev}`$ to $`1.85`$ at $`Q=100\mathrm{Gev}`$. It is important to notice that $`\left(\epsilon _R\epsilon _L\right)^2`$ vanishes if one does not take into account non-inclusive corrections to the shape function (4.5), $`\delta f_{\mathrm{non}\mathrm{incl}}=0`$. Using (4.6) and (4.7) one gets
$$\lambda _2=1.70\mathrm{GeV}^2,\left(\epsilon _R\epsilon _L\right)^2=0.14\mathrm{GeV}^2.$$
(5.12)
It follows from (5.10) that the boundary terms generate a sizable perturbative corrections to the second moment of the $`C`$parameter distribution and diminish the magnitude of scales parameterizing $`1/Q`$power corrections. Moreover, one finds from (5.4), (5.5), (5.6) and (5.10) that the variance of the distribution, $`e^2e^2`$, does not receive $`1/Q`$power corrections for the $`t`$ and $`\rho `$variables while for the $`C`$parameter the boundary terms produce a negative $`1/Q`$correction
$$C^2C^2=C^2_{_{\mathrm{PT}}}C_{_{\mathrm{PT}}}^23.23\frac{\lambda _1}{Q}\alpha _s(Q)+𝒪(1/Q^2).$$
(5.13)
The comparison of the QCD predictions (5.10) with the experimental data is shown in Fig. 4. We would like to recall that the obtained expressions for the moments do not take into account the contribution of multi-jet final states configurations and assume a smallness of $`1/(Q^2e)`$corrections to the distributions (2.6). It is interesting to note that the last assumption is supported by the recent analysis of the power corrections to the first two moments of the thrust distribution in the single dressed gluon approximation . This analysis is complimentary to our studies since it does not resum leading power corrections in the two-jet region and takes into account the contribution coming from the region $`ee_{\mathrm{max}}`$.
## 6. Conclusions
In this paper we have studied the power corrections to the thrust, $`t`$, heavy jet mass, $`\rho `$, and $`C`$parameter distributions in the two-jet kinematical region. Our analysis was based on the observation that perturbative and nonperturbative effects can be separated in the differential event shape distributions into calculable Sudakov resummed distribution, $`d\sigma _{\mathrm{PT}}/de`$, and nonperturbative shape function, $`f(\epsilon _R,\epsilon _L)`$, respectively. Each of them depends separately on the factorization scale $`\mu `$ but this dependence cancels in their product. The shape function describes the energy flow into two hemispheres in the final state. It does not depend on the center-of-mass energy $`Q`$ as well as on the choice of the event shape variable $`e=t,\rho `$ and $`C`$.
We demonstrated that away from the end-point region, $`e\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$, nonperturbative corrections to the distributions have a simple form (2.8) with the leading $`1/Q`$power correction parameterized by a single scale given by the first moment of the shape function. In this region, to which all performed experimental analysis have been restricted so far, our predictions for the thrust and heavy mass distributions and their mean values coincide with those of IR renormalon based models while for the $`C`$parameters we find an additional sizeable perturbative contribution modifying the magnitude of the $`1/Q`$power correction (5.6).
In the end-point region, $`e\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$, the obtained factorized expressions for the distributions take into account power corrections of the form $`1/(Qe)^n`$ for arbitrary $`n`$. They are controlled by the shape function through (2.11) and are sensitive to the choice of this function. Comparing the QCD predictions with the data we have chosen the simplest ansatz for the shape function (4.6) which is consistent with general properties of nonperturbative distributions and includes nonzero correlations between energy flows into different hemispheres. Examining the dependence of the distributions on the corresponding parameter of the shape function we have observed that these correlations play an important rôle and are not negligible.
### Acknowledgments
We would like to thank E. Gardi and G. Sterman for very interesting discussions. We are grateful to O. Biebel, G. Salam and B. Webber for useful correspondence. This work was supported in part by the EU network “Training and Mobility of Researchers”, contract FMRX–CT98–0194 (G.K.) and the BFA fellowship – Bourse de coopération Franco-Algérienne (S.T.).
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# HNCO in massive galactic dense cores Based on the observations collected at the European Southern Observatory, La Silla, Chile and on observations with the Heinrich-Hertz-Telescope (HHT). The HHT is operated by the Submillimeter Telescope Observatory on behalf of Steward Observatory and the MPI für Radioastronomie. Tables 1, 2, 5, 6 are also available in electronic form and Tables 7–14 are only available in electronic form at the CDS via anonymous ftp to cdsarc.u-strasbg.fr (130.79.128.5) or via http://cdsweb.u-strasbg.fr/Abstract.html
## 1 Introduction
Systematic studies of dense molecular cores in regions of high mass star formation (HMSF) are of great importance for our general understanding of star formation. In comparison with low mass star formation regions, so far only a few rather arbitrarily selected cores associated with HMSF have been investigated in some detail.
In recent years we performed extensive surveys of dense cores in regions of high mass star formation, mainly in CS (Zinchenko et al. 1995, 1998). We used water masers as signposts of high mass star formation. Both outer and inner Galaxy were covered by these surveys ($`l120\mathrm{°}210\mathrm{°}`$ and $`260\mathrm{°}308\mathrm{°}`$). The innermost part of the Galaxy ($`l308\mathrm{°}360\mathrm{°}`$) was observed in a similar way by Juvela (1996). In addition, sources associated with water masers were surveyed in thermal SiO (Harju et al. (1998)) which is supposed to be a good indicator of shocks in molecular clouds. From these observations we derived basic physical parameters of the cores and constructed their statistical distributions (Zinchenko (1995), Zinchenko et al. (1998)). In order to investigate a range of core densities, observations of lines with different excitation conditions are needed. One of the interesting candidates is the HNCO (isocyanic acid) molecule.
HNCO was first detected by Snyder & Buhl (1972) in Sgr B2. Subsequent studies have concentrated mostly on the Galactic center region where the HNCO emission was found to be particularly strong (e.g., Churchwell et al. 1986, Wilson et al. 1996, Lindqvist et al. 1995, Kuan & Snyder 1996, Dahmen et al. 1997, Sato et al. 1997). A survey of HNCO emission throughout the Galaxy was made by Jackson et al. (1984) in the $`J_{\mathrm{K}_1\mathrm{K}_1}=5_{05}4_{04}`$ and $`4_{04}3_{03}`$ transitions with the 11 m NRAO telescope. Seven (from 18) clouds including Orion KL were detected at rather low levels of intensity (typically $`0.2`$ K on a $`T_\mathrm{A}^{}`$ scale). Churchwell et al. (1986) obtained strict upper limits on HNCO $`1_{01}0_{00}`$ and $`2_{02}1_{01}`$ emission towards about 20 galactic sources with the 36.6 m Haystack antenna.
HNCO is a slightly asymmetric rotor. Its levels may be designated as $`J_{\mathrm{K}_1\mathrm{K}_1}`$ where $`J`$ is the total angular momentum and $`K_1`$, $`K_1`$ are quantum numbers corresponding to the projection of $`J`$ on the symmetry axis for the limiting cases of prolate and oblate symmetric top, respectively (e.g. Townes & Schawlow 1975). The structure of the HNCO energy levels can be represented as a set of “ladders” with different $`K_1`$ values, like for a symmetric top. However, due to the asymmetry of the molecule radiative transitions between different $`K_1`$ ladders ($`b`$-type transitions) are allowed and, moreover, they are very fast. The corresponding component of the dipole moment is similar to its component for transitions inside the $`K_1`$ ladders ($`a`$-type transitions). Churchwell et al. (1986) found that as a result the HNCO excitation is governed mostly by radiative rather than collisional processes (at least in Sgr B2).
On the basis of their estimates of source parameters Jackson et al. (1984) concluded that HNCO is a potentially valuable probe of the densest regions ($`n10^6`$ cm<sup>-3</sup>) of molecular clouds. It was shown also that HNCO is rather sensitive to far infrared (FIR) radiation fields due to the fact that the lowest levels of the $`K_1=0,`$ 1 and $`K_1=1,`$ 2 ladders are separated by energies corresponding to FIR wavelengths (330 $`\mu `$m and 110 $`\mu `$m, respectively).
From this consideration it is clear that multitransitional data are needed to understand HNCO excitation and to derive the source properties. Bearing this in mind we undertook a survey of HNCO emission in various rotational lines, also trying to detect emission from higher excited $`K`$ladders ($`K_1>0`$). Five cores were mapped in HNCO to estimate the extent of the emission.
Several other species were observed simultaneously with HNCO. The most prominent are C<sup>18</sup>O and SO. In the following we thus also compare HNCO with C<sup>18</sup>O.
## 2 Observations
### 2.1 Source list
For this study we observed those dense cores showing particularly strong CS emission ($`T_{\mathrm{mb}}>3`$ K) in the surveys of Zinchenko et al. (1995, 1998) and Juvela (1996). Several strong SiO ($`v=0`$) sources detected by Harju et al. (1998) are also included in our sample. Sources observed at the SEST and at Onsala are presented in Tables 1, 2. Sources also observed at Effelsberg or at the HHT are marked in both tables.
We designate most sources according to their galactic coordinates. Exceptions are Orion KL and Sgr A. For Sgr A we use the position observed by Jackson et al. (1984) for comparison (known as the M-0.13-0.08 cloud, see Lindqvist et al. 1995). Common identifications with some well known objects are given in the last column.
### 2.2 Observational procedures
The most important parameters of our SEST-15m, OSO-20m, Effelsberg 100-m and HHT measurements are summarized in Tables 3, 4. Further details are given below for each instrument.
#### 2.2.1 SEST observations
The observations were performed with SIS receivers in a single-sideband (SSB) mode using dual beam switching with a beam throw of $`12`$′. At 220 GHz we used 2 acousto-optical spectrometers in parallel: (1) a 2000 channel high-resolution spectrometer (HRS) with 86 MHz bandwidth, 43 kHz channel separation and 80 kHz resolution and (2) a 1440 channel low-resolution spectrometer (LR1) with a 1000 MHz total bandwidth, 0.7 MHz channel separation and 1.4 MHz spectral resolution. The LR1 band was centered on the HNCO $`10_{0,10}9_{0,9}`$ transition. However, it covered some other HNCO transitions too (see Table 4) as well as C<sup>18</sup>O (2–1), SO ($`6_55_4`$) and other lines (Fig. 1 shows a typical spectrum).
The 110 and 154 GHz observations were performed simultaneously; the spectra were recorded by the HRS which band was split into two equal parts. The 220 GHz HRS spectra were smoothed to 170 kHz resolution and the 110 and 154 mm spectra were smoothed to 86 kHz resolution. Pointing was checked periodically by observations of nearby SiO masers; the pointing accuracy was $`5`$″.
The standard chopper-wheel technique was used for the calibration. We express the results in units of main beam brightness temperature ($`T_{\mathrm{mb}}`$) assuming the main beam efficiencies ($`\eta _{\mathrm{mb}}`$) as given in Table 3. The temperature scale was checked by observations of Orion KL.
In most sources only one position was observed, corresponding typically to the peak of the CS emission. In addition, G 270.26+0.83 and G 301.12$``$0.20 were mapped with a spacing of 10″.
#### 2.2.2 Onsala observations
At Onsala, the 110 GHz observing procedure was very similar to that at the SEST. The observations were also performed in a dual beam switching mode with a beam throw of 11$`\stackrel{}{.}`$5. The front-end was a SIS receiver tuned to SSB operation. As backend we used 2 filter spectrometers in parallel: a 256 channel filterbank with 250 kHz resolution and a 512 channel filterbank with 1 MHz resolution. The calibration procedure was the same as at the SEST. The pointing accuracy checked by observations of nearby SiO masers was $`5`$″. The strongest HNCO source from the Onsala sample, W51M, was mapped with 40″ spacing.
#### 2.2.3 Effelsberg observations
The 22 GHz observations in Effelsberg were performed with a K-band maser amplifier using position switching. The offset positions were displaced by 10′–15′ symmetrically in azimuth. Pointing was checked periodically by observations of nearby continuum sources; the pointing accuracy was $`10`$″. The integration time per position was a few hours.
The main beam temperature scale was checked by observations of nearby continuum calibration sources, NGC 7027 and W3(OH); for Sgr A we used Sgr B2. The fluxes for the first two sources were taken from Ott et al. (1994). The Sgr B2 flux at 1.3 cm was taken from Martín-Pintado et al. (1990).
#### 2.2.4 HHT observations
To observe the HNCO $`J`$ = 21–20 lines at 461 GHz we have used the Heinrich Hertz Telescope (HHT) on Mt. Graham (Baars & Martin 1996) during Feb. 1999 with a beamwidth of 18<sup>′′</sup>. Spectra were taken employing an SIS receiver with backends consisting of two acousto optical spectrometers with 2048 channels each, channel spacing $``$480 and $``$120 kHz, frequency resolution $``$930 and 230 kHz, and total bandwidths of $``$1 GHz and 250 MHz, respectively. Receiver temperatures were $``$150 K, system temperatures were $``$1000 K on a $`T_\mathrm{A}^{}`$ scale. The receiver was sensitive to both sidebands. Any imbalance in the gains of the lower and upper sideband would thus lead to calibration errors. To account for this, we have observed the CO $`J`$ = 4–3 line of Orion-KL with the same receiver tuning setup and obtain $`T_\mathrm{A}^{}`$ $``$ 70 K, in good agreement with Schulz et al. (1995).
HNCO $`J=1615`$ (351.63346 GHz) and $`J=1514`$ (329.66454 GHz) line emission was observed with a dual channel SIS receiver in early April 1999 at the HHT. The beamwidth was 22″, receiver temperatures were 135 K; system temperatures were $``$700 K on a $`T_\mathrm{A}^{}`$ scale. The receivers were also sensitive to both sidebands. We have used published spectra from Orion-KL and IRC+10216 as calibrators (Groesbeck et al. (1994); Schilke et al. (1997)).
All results displayed are given in units of main beam brightness temperature ($`T_{\mathrm{mb}}`$). This is related to $`T_\mathrm{A}^{}`$ via $`T_{\mathrm{mb}}`$ = $`T_\mathrm{A}^{}`$ ($`F_{\mathrm{eff}}`$/$`B_{\mathrm{eff}}`$) (cf. Downes 1989). The main beam efficiency, $`B_{\mathrm{eff}}`$, was 0.38 at 461 GHz and 0.5 at 330 and 352 GHz as obtained by measurements of Saturn. The forward hemisphere efficiency, $`F_{\mathrm{eff}}`$, is 0.75 at 461 GHz and 0.9 at 330 and 352 GHz (D. Muders, priv. comm.). The HHT is with an rms surface deviation of $``$20$`\mu `$m (i.e. $`\lambda `$/30 at 461 GHz) quite accurate. Thus emission from the sidelobes should not be a problem.
Pointing was obtained toward Jupiter (continuum pointing) and toward Orion-KL and R Cas (line pointing) with maximum deviations of order 5<sup>′′</sup>. Observations were carried out in a position switching mode with the off-position $``$1000<sup>′′</sup> offset from the source position.
### 2.3 Data reduction and analysis
We have reduced the data and produced maps using the GAG (Groupe d’Astrophysique de Grenoble) software package. The measured spectra were fitted by one or more gaussian components.
## 3 Results
### 3.1 One-point observations
HNCO was detected in 36 SEST sources (from 56 observed) and in 22 OSO sources (from 27). Because of one source belonging to both samples, the total number of detected objects is 57. In many cases $`K_1>0`$ transitions were detected too. The gaussian line parameters are presented in Tables 5–14 (Tables 7–14 are available only electronically). It is worth noting that a single-gaussian fit is clearly insufficient in many cases because the lines have broad wings and other non-gaussian features. Therefore, the values in the tables give only a rough representation of the line profiles (the integrated intensities were obtained by integrating over the lines in most cases).
Table 5 summarizes the 220 GHz SEST results for HNCO $`10_{0,10}9_{0,9}`$ and C<sup>18</sup>O. The Onsala $`5_{05}4_{04}`$ and C<sup>18</sup>O results are presented in Table 6. The 220 GHz results for the $`K_1=2`$, 3 ladders are given in Tables 7, 8. The 110 and 154 GHz SEST data are displayed in Table 9. The Onsala 88 GHz data are summarized in Table 10. The Effelsberg data are presented in Table 11. Tables 12–14 contain the HHT data. We fitted the Effelsberg spectra with 3-component gaussians with fixed separations corresponding to the hyperfine structure of the $`1_{01}0_{00}`$ transition.
Examples of measured spectra are given in Figs. 2, 3. Fig. 2 shows spectra of a few sources covering $`K_1=0`$, 2 and 3 transitions at 220 GHz. Fig. 3 presents HNCO spectra in the HNCO $`K_1=0`$ transitions at different wavelengths for several sources.
The HNCO line profile in Orion KL can be decomposed into at least two components which likely correspond to the so-called classical “Hot Core” and “Plateau” outflow components (see, e.g., Harris et al. 1995). The ratio between these components is practically the same for the $`5_{05}4_{04}`$, $`7_{07}6_{06}`$ and $`10_{0,10}9_{0,9}`$ lines: $`60`$% of the emission originates from the “Plateau” outflow source. The other lines do not allow such decomposition due to their weakness or blending with other spectral features.
An inspection of Table 5 shows that the derived C<sup>18</sup>O velocities are systematically lower (more negative) than the HNCO ones. The difference is $`1`$ km/s on the average. This can be an instrumental effect: the C<sup>18</sup>O line was located far away from the center of the spectrometer band and a possible non-linearity in the frequency response could lead to the apparent displacement of the line on the velocity axis. This remark is applicable also for the higher $`K_1`$ HNCO lines.
### 3.2 Maps
In order to estimate source sizes and their spatial association with YSO and infrared (IR) sources we mapped 2 southern sources in the $`10_{0,10}9_{0,9}`$ HNCO line and Orion KL, W49N and W51M in the $`15_{0,15}14_{0,14}`$, $`16_{0,16}15_{0,15}`$ and $`21_{0,21}20_{0,20}`$ lines. W51M was mapped also in the $`5_{05}4_{04}`$ line. Three of these maps are presented in Fig. 4.
The sources remain spatially unresolved. E.g. for G 301.12–0.20 we obtain a FWHM $`29`$″ in right ascension (from the strip scan across the map) which is very close to the beam size at this frequency (24″).
### 3.3 Detection of the $`K_1=5`$ HNCO transition
The highest $`K_1`$ HNCO transition reported so far was $`K_1=4`$ (the $`10_49_4`$ line) in Orion (Sutton et al. 1985). This line is located on the shoulder of the strong C<sup>18</sup>O $`J=21`$ line. In Fig. 5 we show parts of our Orion 220 GHz low resolution spectrum and 461 GHz spectra with $`K_1=2`$, 3, 4 and even 5 features (the $`K_1=1`$ transition is outside our band). The rest frequencies are assumed to be equal to those given in the JPL catalogue for the strongest components of the corresponding transitions (for $`K_1=2`$ at 220 GHz we took the mean of the frequencies of the two strongest components).
There is a weak bump in the redshifted C<sup>18</sup>O $`J=21`$ wing which can be attributed to HNCO $`10_49_4`$. Due to the uncertainty in fitting the C<sup>18</sup>O line profile the intensity of the HNCO feature cannot be reliably determined but it is lower than reported by Sutton et al. (1985). Our best estimate for the integrated intensity is $`T_{\mathrm{mb}}𝑑v0.7`$ K km/s, but a reliable error cannot be given.
There is also a feature at the $`K_1=5`$ frequency in the 220 GHz spectrum. It is located in the wing of a C<sub>2</sub>H<sub>3</sub>CN line. The integrated intensity is $`T_{\mathrm{mb}}𝑑v=0.27\pm 0.06`$ K km/s. The identification of this feature with HNCO seems to be reliable. The only other candidate is the C<sub>2</sub>H<sub>5</sub>OH $`14_{21,2}13_{11,2}`$ line at 219391.81 MHz. However, there is no sign of other ethanol lines in our spectrum so we reject this alternative. In the 461 GHz spectrum the $`K_1=5`$ feature is clearly detected. Its integrated intensity is $`T_{\mathrm{mb}}𝑑v=1.4\pm 0.2`$ K km/s.
### 3.4 Hyperfine splitting, HN<sup>13</sup>CO and optical depths
The HNCO lines are split into several hyperfine components mainly due to the <sup>14</sup>N spin. This splitting is clearly seen in the $`1_{01}0_{00}`$ transition (Fig. 3) at 22 GHz. Earlier HNCO hyperfine structure in the $`1_{01}0_{00}`$ line was only observed in the dark cloud TMC-1 (Brown (1981)) where possible deviations from the optically thin LTE (Local Thermodynamic Equilibrium) intensity ratios (3:5:1) were found. In our spectra the hyperfine ratios are consistent with the optically thin LTE values. Taking into account the measurement uncertainties, an upper limit on the optical depth in this transition for the sources detected in Effelsberg is $`\tau 1`$.
To the best of our knowledge no isotopomer of HNCO except the main one has been detected in space yet. This detection would be important for estimates of HNCO optical depths which are believed to be small (e.g. Jackson et al. 1984, Churchwell et al. 1986). The frequency separations between the HN<sup>13</sup>CO and the main isotopomer lines are rather small corresponding to a few km/s, so in sources with broad lines like Orion A or Sgr A the HN<sup>13</sup>CO lines will be blended. However, there are some strong HNCO sources in our sample with narrower lines which show features attributable to HN<sup>13</sup>CO. The most reliable one is seen in the G 301.12–0.20 $`10_{0,10}9_{0,9}`$ spectrum (Fig. 6). A weak feature on the blue shoulder of the main isotope line is very close in frequency to the expected location of the HN<sup>13</sup>CO line .
For comparison we show in addition to HNCO also the C<sup>34</sup>S spectrum. It is noteworthy that there is no bump in this spectrum corresponding to the discussed feature in HNCO.
The line we identify with HN<sup>13</sup>CO is shifted by $`0.65\pm 0.21`$ MHz from the expected HN<sup>13</sup>CO transition frequency. This $`3\sigma `$ shift, if it is significant, cannot be explained by instrumental effects like in the case of our C<sup>18</sup>O data because the feature is very close to the main isotope line. The shift greatly exceeds the uncertainty of the transition frequency derived from the laboratory data (Winnewisser et al. 1978) which is 25 kHz. This makes the identification questionable. Detection of other HN<sup>13</sup>CO lines would be important in this respect. There is no corresponding feature in the $`7_{07}6_{06}`$ HNCO spectrum (the $`5_{05}4_{04}`$ spectrum is too noisy). This could mean that the optical depth in this transition is significantly lower. Indeed, at sufficiently high temperatures ($`>30`$ K) it can be about 2 times lower than in the $`10_{0,10}9_{0,9}`$ transition according to Eq. (2) (see the discussion in Sect. 4.2).
If our identification of the discussed line with HN<sup>13</sup>CO is correct we can estimate the optical depth assuming the same excitation as for the main isotopomer. For G 301.12–0.20 we obtain $`\tau (\text{HNCO})15`$ if we assume the terrestrial <sup>12</sup>C/<sup>13</sup>C isotope ratio (<sup>12</sup>C/<sup>13</sup>C = 89) and $`7`$ for <sup>12</sup>C/<sup>13</sup>C = 40. A high optical depth in the $`10_{0,10}9_{0,9}`$ HNCO line does not contradict our conclusion of low optical depth in the $`1_{01}0_{00}`$ transition because the line strengths for these transitions are different (see discussion in Sect. 4.2). Therefore, the optical depth in some lines of the main isotopomer might be rather high. This contradicts the usual assumption of low optical depth in all HNCO lines (e.g., Jackson et al. 1984, Churchwell et al. 1986) and could imply serious consequences for the analysis of HNCO excitation and abundances.
## 4 Discussion
### 4.1 Comparison with other HNCO data
Most of our HNCO sources are new detections. Only few were included in the surveys of Jackson et al. (1984) and Churchwell et al. (1986). A direct comparison with the intensities measured by Jackson et al. is impossible due to different temperature scales. Common detected sources are Orion KL and W51. Their upper limit for W3(OH) does not contradict our value if we take into account the difference in the temperature scales. The upper limits for the $`1_{01}0_{00}`$ transition obtained by Churchwell et al. do not contradict our results taking into account the differences in the beam sizes and efficiencies.
As mentioned above, towards Orion KL several HNCO lines were observed at 220 GHz by Sutton et al. (1985). Their results agree in general with our measurements though there is a discrepancy concerning the intensity of the $`K_1=4`$ transition (Sect. 3.3).
It is worth noting that while at 22 GHz and at 110 GHz (as obtained by Jackson et al. 1984) the brightest source of HNCO emission is the Galactic center, at 220 GHz the situation changes and Orion becomes the brightest source with several other sources approaching Sgr A in intensity. Apparently this is caused by differences in excitation.
### 4.2 Rotational diagrams
As a first step in the excitation analysis we construct traditional rotational diagrams for our sources. For a recent discussion of this method see e.g. Goldsmith & Langer (1999). This means a plot of the column density ($`N_\mathrm{u}`$) per statistical weight ($`g_\mathrm{u}`$) of a number of molecular energy levels, as a function of their energy above the ground state ($`E_\mathrm{u}`$). In local thermodynamic equilibrium (LTE), this will just be a Boltzmann distribution, so a plot of $`\mathrm{ln}(N_\mathrm{u}/g_\mathrm{u})`$ versus $`E_\mathrm{u}/k`$ will yield a straight line with a slope of $`1/T_\mathrm{R}`$. The temperature inferred is often called the “rotational temperature”.
Actually from the measurements we do not obtain directly the column densities. The measured quantity is the line intensity. In an optically thin case for $`T_{\mathrm{ex}}>>T_{\mathrm{bg}}`$ ($`T_{\mathrm{ex}}`$ is the excitation temperature of the transition and $`T_{\mathrm{bg}}`$ is the background temperature)
$$\mathrm{log}\left[\frac{3k(W/f_\mathrm{b})}{8\pi ^3\nu \mu _\mathrm{x}^2S}\right]=\mathrm{log}\left(\frac{N}{Q}\right)\frac{E_\mathrm{u}}{kT_\mathrm{R}}\mathrm{log}e$$
(1)
where $`W`$ is the integrated line intensity, $`f_\mathrm{b}`$ is the beam dilution factor, $`S`$ is the line strength, $`\mu _\mathrm{x}`$ is the appropriate component of the dipole moment, $`N`$ is the total column density and $`Q(T_\mathrm{R})`$ is the partition function.
The quantity on the left hand side of Eq. (1) can be derived from the molecular data. Plotting it versus $`E_\mathrm{u}`$ we can find the rotational temperature (from the slope) and the total column density (from the intercept).
Some problems can arise from an uncertainty in the beam filling factor. As shown in Fig. 4 the sources are probably unresolved. Assuming that the source size is the same for all HNCO transitions in a given source and that the source size is small with respect to the beam, we reduced all data to the same beam size, the SEST HPBW at 220 GHz, i.e. 24″.
For Orion the highest observed transition lies $`1300`$ K above the ground level. For other sources we managed to observe transitions up to $`450`$ K above the ground state. Examples of the rotational diagrams are presented in Figs. 7, 8.
The measured integrated intensities are represented by filled squares ($`f_\mathrm{b}=1`$). The corrected results are plotted by open squares in Figs. 7, 8. One can see that they much better correspond to each other than the uncorrected values.
The rotational diagram for Orion is presented in Fig. 8. The rotational temperature from this plot is $`T_{\mathrm{rot}}25`$ K for the lowest transitions and $`T_{\mathrm{rot}}530`$ K for the highest transitions. The latter one is a very high value even for Orion KL. But in principle the diagram shows a range of rotational temperatures. We represent it by 3 components as shown in Table 15. A separate fit to the $`K_1=0`$ transitions gives $`T_{\mathrm{rot}}80`$ K (although this fit is not very satisfactory).
The rotational temperatures and column densities derived from rotational diagrams are summarized in Table 15. In this analysis we assume that the sources are optically thin in the observed transitions. This contradicts the tentative detection of HN<sup>13</sup>CO in G 301.12–0.20. The effects of high optical depth on rotational diagrams have been analyzed recently by Goldsmith & Langer (1999). In optically thick case the column density in the upper level of the transition ($`N_\mathrm{u}`$) is underestimated by the factor of $`\tau /(1e^\tau )`$ and, therefore, corresponding points in the population diagram lie lower than they should. In general, for linear molecules it produces a curvature resembling that seen in the diagrams for Orion and some other sources. It is caused by the fact that the optical depth exhibits a peak for transitions with the excitation energy $`E_\mathrm{u}kT`$ (Goldsmith & Langer 1999). However, for nonlinear molecules the optical depth effect rather leads to a “scatter” in the population diagram, because transitions with significantly different optical depth can have similar excitation energies.
There is a strong argument against high optical depth at least for transitions with $`E_\mathrm{u}200400`$ K in Orion. In this range transitions with similar energies of the upper state but with very different frequencies (belonging to different $`K_1`$ ladders) were observed. It is easy to estimate the expected ratio of peak optical depths in the lines which is
$$\frac{\tau _1}{\tau _2}=\frac{S_1}{S_2}\frac{\mathrm{exp}\left(\frac{h\nu _1}{kT}\right)1}{\mathrm{exp}\left(\frac{h\nu _2}{kT}\right)1}\mathrm{exp}\left(\frac{E_\mathrm{u}^2E_\mathrm{u}^1}{kT}\right)$$
(2)
For $`|E_\mathrm{u}^1E_\mathrm{u}^2|<<kT`$ the exponential factor is close to unity.
In our data there are pairs of transitions with similar upper state energies. The $`21_{0,21}20_{0,20}`$ and $`10_{2,9/8}9_{2,8/7}`$ transitions have similar $`E_\mathrm{u}200`$ K. However, the first one has higher line strength and higher transition frequency; therefore, according to Eq. (2) it should have higher optical depth than the second one. Then, it should be stronger influenced by possible optical depth effects and the corresponding point in Fig. 8 should lie lower than the point corresponding to the $`10_{2,9/8}9_{2,8/7}`$ transitions. However, this is not a case. Actually, the points are very close to each other and perhaps slightly shifted in the opposite sense. The same is true for the $`21_{2,20/19}20_{2,19/18}`$ and $`10_{3,8/7}9_{3,7/6}`$ transitions with $`E_\mathrm{u}400`$ K. We conclude that the optical depth for Orion in these transitions should be low. Perhaps in some other transitions or in other sources optical depths are as high as indicated by our tentative HN<sup>13</sup>CO detection. There is however no reason to apply optical depth corrections to the bulk of our sources.
Transitions with low $`E_\mathrm{u}/k`$ values are fitted by rather low temperature models, $`T_\mathrm{R}1030`$ K. Transitions between higher excited states are related to higher rotational temperatures up to $`T_{\mathrm{rot}}500`$ K. In Table 15 we also present estimates of the HNCO relative abundances. The hydrogen column densities have been calculated from the C<sup>18</sup>O data under the assumptions of LTE and a C<sup>18</sup>O relative abundance of $`\mathrm{1.7\hspace{0.17em}10}^7`$ (Frerking et al. 1982). Typical HNCO abundances are $`10^9`$. Sgr A does not look very exceptional here. The relative HNCO abundance in Sgr A is about the same as in Orion but the rotational temperature is much lower. In contrast to many other sources there is no high excitation temperature component in Sgr A, indicating that the dense gas is probably cool. This agrees with results from Hüttemeister et al. (1998) based on SiO and C<sup>18</sup>O. The opposite scenario, a hot highly subthermally excited low density gas component ($`n`$(H<sub>2</sub>) $`10^4`$ cm<sup>-3</sup>) as observed by Hüttemeister et al. (1993) in ammonia toward Sgr B2 is less likely, due to the correlations between HNCO and SiO that will be outlined in Sects. 4.4 and 4.6.
It is important to emphasize that our estimates give lower limits to the relative abundance $`X`$(HNCO) = $`N`$(HNCO)/$`N`$(H<sub>2</sub>) for at least two reasons. First, the HNCO sources are much more compact than their C<sup>18</sup>O counterparts and tend to be spatially unresolved. Our estimates give beam averaged values and “real” abundances in regions of HNCO line formation should be significantly higher. Second, if the HNCO optical depth is high we would underestimate its column densities.
Next, we have to mention that all these estimates refer to the bulk of the cores. In the high velocity gas the HNCO abundances are apparently much higher.
One might think that better estimates of HNCO abundances can be obtained from comparison with the dust emission rather than with C<sup>18</sup>O. As shown, HNCO probably arises in “warm” environments and in the dust emission we see preferentially a high temperature medium while in C<sup>18</sup>O the reverse is true. However, interferometric observations in Orion (Blake et al. (1996)) show that HNCO $`K_1=2`$ and dust distributions do not entirely coincide. At the same time, as shown in Sect. 4.5, there is a tight correlation between the FIR emission at 100 $`\mu `$m and C<sup>18</sup>O(2–1) integrated line intensity. Therefore, no large differences between estimates of HNCO abundances by both methods can be expected. There are detailed studies of dust emission towards some of our sources with comparable angular resolution. E.g. Henning et al. (2000) show that total gas column densities derived from dust and from C<sup>18</sup>O(2–1) in G301.12–0.20 coincide within a factor of 3.
In Fig. 9 we plot the HNCO abundances versus the HNCO line widths. There is a trend of increasing the HNCO abundance with increasing HNCO line width. This shows that the HNCO production can be related to dynamical activity in the sources.
Table 15 and Fig. 9 indicate that abundances derived for the sources which belong to the inner and to the outer Galaxy, respectively, are about the same. Therefore, there is no significant galactic gradient in HNCO abundance.
### 4.3 Physical conditions in regions of HNCO emission
Now we shall try to understand the physical conditions in regions of HNCO emission detected by us. An important question to start with is which excitation mechanism dominates, radiative or collisional? And which gas parameters are implied by each of them? To answer these questions properly would require a numerical model taking both into account. Useful conclusions can, however, also be obtained by semi-qualitative consideration presented below. We concentrate here on Orion KL as the best studied source.
At first, we need an estimate for the size of the HNCO emission region. Our map presented in Fig. 4 gives an upper limit of $`10\mathrm{}`$ for the $`21_{0,21}20_{0,20}`$ transition. Interferometric results (Blake et al. (1996)) give a size of $`2\mathrm{}\times 4\mathrm{}`$ for the $`K_1=2`$ transition at 220 GHz. This can be probably considered as an upper limit also for higher $`K_1`$ ladders. On the other hand we can obtain a lower limit on the source size from the comparison of the brightness and excitation temperatures. For $`T_{\mathrm{ex}}500`$ K (as follows from the population diagram) we obtain that the lower limit on the beam filling factor for the $`K_1=5`$ transitions in Orion is $`\mathrm{2\hspace{0.17em}10}^3`$. Therefore, the effective size of the emitting region is $`1`$″ or $`0.002`$ pc, i.e. $`\mathrm{7\hspace{0.17em}10}^{15}`$ cm.
Let us consider the physical requirements in the case of collisional excitation. The critical densities defined as $`n_\mathrm{c}=A_{\mathrm{ul}}/C_{\mathrm{ul}}`$ ($`A_{\mathrm{ul}}`$ is the spontaneous decay rate and $`C_{\mathrm{ul}}`$ is the collisional de-excitation rate; Scoville et al. 1980) are $`10^6`$ cm<sup>-3</sup> for the $`10_{0,10}9_{0,9}`$ transition and $`10^7`$ cm<sup>-3</sup> for the $`21_{0,21}20_{0,20}`$ transition ( the collisional rates are $`10^{10}`$ s<sup>-1</sup>cm<sup>-3</sup> as obtained from Sheldon Green’s program available on Internet – http://www.giss.nasa.gov/data/mcrates/). Much higher densities are needed for excitation of the transitions in the $`K_1>0`$ ladders. This is caused by fast $`b`$-type transitions between different $`K_1`$ ladders. E.g. the spontaneous emission rate from the $`K_1=5`$ ladder to the $`K_1=4`$ ladder is $`5`$ s<sup>-1</sup>. This implies a critical density of $`10^{11}`$ cm<sup>-3</sup>. The gas kinetic temperature should be $`500`$ K.
Such conditions cannot be excluded. Walker et al. (1994) derived from observations of vibrationally excited CS $`n10^{11}10^{12}`$ cm<sup>-3</sup> and $`T1000`$ K in a region $`10^{15}`$ cm from the stellar core toward IRAS 16293–2422. The question is whether the required amount of such gas is consistent with the observations.
Taking into account the lower limit on the source size the mass of the hot dense gas ($`n10^{11}`$ cm<sup>-3</sup>, $`T500`$ K) would be $`100`$ M. Estimates of the hot core mass from dust continuum measurements give values of $`540`$ M (Masson & Mundy 1988, Wright et al. 1992). Taking into account the uncertainties in our estimations we cannot entirely exclude the possibility of collisional excitation even for the $`K_1=5`$ ladder but this appears to be an unlikely scenario.
For the lower $`K_1`$ ladders the density requirements can be significantly relaxed. E.g. for the $`b`$-type transitions from the $`K_1=3`$ to the $`K_1=2`$ ladder the spontaneous decay rate is $`1`$ s<sup>-1</sup> and the critical density is $`10^{10}`$ cm<sup>-3</sup>.
The transitions in the $`K_1=0`$ ladder, of course, will be also excited in this hot dense gas. However, the emission in these lines will be dominated by a more extended lower density component.
Now let us turn to radiative excitation. It requires sufficient photons at the wavelengths corresponding to the $`b`$-type transitions between different $`K_1`$ ladders, from $`300`$ to $`30`$ $`\mu `$m. If the dilution factor is close to unity we need an optical depth $`\tau 1`$ and a radiation temperature $`T_\mathrm{R}500`$ K at least at 30 $`\mu `$m. As an upper limit to the source size we can take the mean interferometric value of $`3`$″. However, what will be the IR flux and luminosity of such a source? For the flux at 30 $`\mu `$m we obtain $`F(30\mu \mathrm{m})\mathrm{3\hspace{0.17em}10}^4(\theta _\mathrm{s}/1\mathrm{})^2`$ Jy. The observational value is $`\mathrm{5\hspace{0.17em}10}^4`$ Jy (van Dishoeck et al. 1998). Therefore, the angular source size should be $`\theta _\mathrm{s}1\stackrel{}{.}5`$ and the linear size $`10^{16}`$ cm. This practically coincides with the lower limit on the source size derived from the beam dilution (see above). Taking the dust absorption coefficient of $`k_\mathrm{m}10^2`$ cm<sup>2</sup>/g (Ossenkopf & Henning 1994) we conclude that the gas density in this region should be $`n\mathrm{3\hspace{0.17em}10}^7`$ cm<sup>-3</sup>. In this case we have no problem to reconcile the mass estimates with the available data.
However, at longer wavelengths the IR pumping from such a source might be not sufficient. Say, for $`\tau \lambda ^2`$ the optical depth at 300 $`\mu `$m will be only $`0.01`$. Therefore, we need even higher gas column and volume densities and/or larger source sizes at longer wavelengths. The latter implies the presence of a temperature gradient in the source which is natural for an internally heated object. The lower $`K_1`$ ladders are apparently excited by radiation with a lower effective temperature.
To conclude, it is much easier to explain the excitation of the higher $`K_1`$ ladders by the radiative process. The source size in Orion should be $`1\mathrm{}2\mathrm{}`$ which agrees with the interferometric image in the $`K_1=2`$ transition at 1.3 mm (Blake et al. 1996).
The emission in the $`K_1=0`$ ladder should be more extended. For Orion again from a comparison between the brightness and excitation temperatures the source size should be $`6\mathrm{}`$. Such a large source size for the $`K_1=0`$ transitions implies that the radiative excitation via $`K_1>0`$ ladders will become inefficient. Therefore, for the $`K_1=0`$ ladder collisional excitation may dominate which implies gas densities $`n10^610^7`$ cm<sup>-3</sup>. This scenario is supported by several sources where the HNCO emission peak is significantly displaced from any known IR source. The most obvious example is G 270.26+0.83 (Fig. 4). This implies either the presence of a very dense prestellar core or a highly obscured young stellar object at this location.
### 4.4 Comparison with C<sup>18</sup>O, CS and SiO data
An obvious step ahead to understand the properties of interstellar HNCO emission is to compare our results with data from other better studied species. The most reliable comparison can be done with our C<sup>18</sup>O data which were observed simultaneously with HNCO.
Fig. 10 shows a noticeable correlation between the HNCO and C<sup>18</sup>O integrated line intensities. However, it is produced apparently by the correlation between the line widths since the correlation between HNCO and C<sup>18</sup>O peak line temperatures is rather weak.
The plot of $`\mathrm{\Delta }V(\text{HNCO})`$ versus $`\mathrm{\Delta }V(\text{C}\text{18}\text{O})`$ looks rather interesting. Concerning the 220 GHz transitions for the narrowest C<sup>18</sup>O lines the HNCO line width is smaller than that of C<sup>18</sup>O. With increasing C<sup>18</sup>O linewidth, however, the HNCO lines broaden faster and become broader than the C<sup>18</sup>O lines. An exception is Sgr A (not shown in the plot) but its C<sup>18</sup>O spectrum is strongly distorted by emission from the reference position.
A similar comparison with the CS(2–1) data from Zinchenko et al. (1995, 1998) and Juvela (1996) (not shown here) shows even lower correlations between the line parameters than in the case of C<sup>18</sup>O. However, in this case the beam sizes for CS and HNCO are different and even the central positions not always coincide.
In contrast, much better correlations exist between the HNCO and SiO line parameters (the latter ones are taken from Harju et al. 1998). Good correlations exist for both integrated and peak intensities. The correlation between the line widths is somewhat worse but one should take into account that the SiO line widths were derived from the second moments of the line profiles while the HNCO widths represent results of the gaussian fits. Anyway, the correlation does exist and the SiO lines are almost always broader than the HNCO lines.
A more detailed comparison with other species should include the line profiles. For Orion, such a comparison is displayed in Fig. 11. It shows that HNCO lines possess an extra wing emission which is less pronounced than in SiO. A similar picture is seen in some other sources.
This comparison shows that HNCO is closely related to SiO which is thought to be produced primarily in shocks and other energetic processes. The comparison with the presumably optically thin C<sup>18</sup>O(2–1) line shows that the HNCO/CO abundance ratio is apparently enhanced in high velocity gas although to a lower degree than for SiO. Since the CO abundance is usually assumed to be constant in bipolar flows (e.g., Cabrit & Bertout (1992); Shepherd & Churchwell E. (1996)) we see that HNCO abundances are enhanced relative to hydrogen, too.
It is interesting to note that the interferometric data for Orion (Blake et al. (1996)) show that the spatial distributions of SiO and HNCO are rather different. However, this does not exclude a common production mechanism. E.g. these species can be formed at different stages in the postshock gas.
### 4.5 Comparison with IR data
The correlation between HNCO integrated line intensities and FIR flux, e.g. at 100 $`\mu `$m taken from IRAS data (Fig. 12), looks rather similar to the relationship between HNCO and C<sup>18</sup>O (Fig. 10). This is natural because there is a rather tight correlation between the 100 $`\mu `$m flux and the C<sup>18</sup>O integrated line intensity (Fig. 13). Such a good correlation shows that C<sup>18</sup>O relative abundances are rather constant and justifies the usage of the HNCO/C<sup>18</sup>O ratio for estimation of HNCO abundances.
### 4.6 HNCO chemistry
In the early work of Iglesias (1977) HNCO was suggested to form via ion-molecule reactions. The sequence leading to HNCO via electron recombination of H<sub>2</sub>NCO<sup>+</sup> is initiated by the formation of NCO<sup>+</sup> (either by a reaction between CN and O$`{}_{}{}^{+}{}_{2}{}^{}`$ or between He<sup>+</sup> and NCO; see also Brown (1981)). The predicted HNCO abundances from this reaction scheme are low. The steady state fractional abundance is of the order $`10^{10}`$ for a model with $`n_{\mathrm{H}_2}10^4`$ cm<sup>-3</sup> (Iglesias (1977)), and still lower for higher densities, because the fractional ion abundances are roughly inversely proportional to the square root of the gas density.
The abundances derived from ion-molecule chemistry are in contradiction with the observations, especially when HNCO is believed to trace high density gas. Recently, a new neutral gas-phase pathway has been suggested by Turner et al. (1999) for translucent clouds: $`\mathrm{CN}+\mathrm{O}_2\mathrm{NCO}+\mathrm{O}`$ followed by $`\mathrm{NCO}+\mathrm{H}_2\mathrm{HNCO}+\mathrm{H}`$. The importance of these reactions can however be questioned, since 1) the abundance of O<sub>2</sub> in the interstellar space is poorly known; and 2) the second reaction probably has an activation barrier of about 1000 K (Turner et al. (1999)).
Chemistry models predict high fractional O<sub>2</sub> abundances (up to $`10^5`$) at late stages of chemical evolution in dense cores and in postshock gas (e.g Caselli et al. (1993); Bergin et al. (1998)). However, the upper limits derived from observations towards several GMC cores (most recently by the SWAS satellite; Melnick et al. (1999)) are about $`10^6`$, which indicates that the oxygen chemistry is not well understood, yet. O<sub>2</sub> is destroyed by UV radiation and in powerful shocks (with shock velocities greater than 26 kms<sup>-1</sup>; Bergin et al. (1998)), and is therefore likely thriving in relatively quiescent dense gas or in regions associated with low velocity shocks. The same should be true for HNCO if the reaction suggested by Turner et al. (1999) is relevant.
The observed correlation between SiO and HNCO integrated line intensities indicates the prevalence of shocks in the HNCO emission regions. Shock heating can therefore provide the means of overcoming the energy barrier in the reaction between NCO and H<sub>2</sub>, and thereby intensify the HNCO production. On the other hand, the fact that the HNCO line widths are smaller than those of SiO could be understood by the destruction of O<sub>2</sub> in high velocity shocks.
In the light of the present observations the neutral reactions suggested by Turner et al. (1999) appear to provide a plausible production pathway of HNCO also in warm GMC cores. The formation of HNCO via grain surface reactions, e.g. through the desorption and subsequent fragmentation of some more complex molecule is an alternative, which to our knowledge has not yet been investigated.
## 5 Conclusions
We have presented the results of an HNCO survey of high mass star-forming cores at frequencies from 22 to 461 GHz. The main conclusions are the following:
1. HNCO is widespread in dense cores forming high mass stars. The detection rate was $`70`$%. There is no significant galactic gradient in its abundance as indicated by the fact that abundances derived for the sources which belong to the inner and to the outer Galaxy, respectively, are about the same.
2. Transitions in higher $`K_1`$ ladders, up to $`K_1=5`$, are detected. The excitation energy reaches $`1300`$ K above the ground level.
3. HN<sup>13</sup>CO is tentatively detected towards G 301.12–0.20. This implies an optical depth in the HNCO $`10_{0,10}9_{0,9}`$ line $`10`$ in this source. The optical depth in the $`1_{01}0_{00}`$ transition is $`\tau 1`$ for the sources detected in this line as inferred from the hyperfine ratios.
4. The sources are compact with sizes $`20\mathrm{}`$.
5. HNCO rotational temperatures vary from $`10`$ K to $`500`$ K. Typical relative abundances are $`10^9`$. These increase with increasing velocity dispersion.
6. The emission in the $`K_1>0`$ ladders is best explained by FIR radiative excitation. In order to provide a sufficiently large dust optical depth at FIR wavelengths taking into account the limitations on the source size, the gas density should be $`n\mathrm{3\hspace{0.17em}10}^7`$ cm<sup>-3</sup>; a temperature $`T500`$ K is needed to excite the $`K_1=5`$ emission in Orion KL. The $`K_1=0`$ transitions can be collisionally excited. The required densities are $`n10^610^7`$ cm<sup>-3</sup>.
7. HNCO correlates well with SiO and does not correlate with CS which is a typical high density probe. HNCO abundances are enhanced in high velocity gas. Probably HNCO production is related to shocks as for SiO. A plausible pathway is gas-phase neutral-neutral reactions at high ($`>1000`$ K) temperatures to overcome an activation barrier that is likely inhibiting the $`\mathrm{NCO}+\mathrm{H}_2\mathrm{HNCO}+\mathrm{H}`$ reaction in a cool interstellar medium.
###### Acknowledgements.
We are very grateful to Dr. J. Harju for his contribution to this work, to Dr. Lars E.B. Johansson for the help with the observations in Onsala, to the SEST staff, to Alexander Lapinov for calculating HNCO line strengths and to the referee, Dr. C.M. Walmsley, for the very useful detailed comments. I.Z. thanks the Helsinki University Observatory and Max-Planck-Institut für Radioastronomie for the hospitality. He was also supported in part by the DFG grant 436 RUS 113/203/0, INTAS grant 93-2168-ext, NASA grant provided via CRDF RP0-841 and grants 96-02-16472, 99-02-16556 from the Russian Foundation for Basic Research. This research has made use of the Simbad database, operated at CDS, Strasbourg, France.
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# 1 Introduction
## 1 Introduction
The search for supersymmetry (SUSY) has been going on at the leading high energy colliders, most notably at LEP and Tevatron (Run-I), for quite some time. From the negative results lower limits on various sparticle masses have been obtained . The prospect of SUSY searches at Tevatron(Run-II) and at the large hadron collider (LHC) has also been studied in great details . The sparticle mass reach of these colliders in different channels have also been estimated.
In most of the analyses it is assumed, for the sake of economy in the number of parameters, that all the scalars in the model, i.e. the squarks, the sleptons and the Higgs bosons, have a common SUSY breaking mass ($`m_0`$) at the grand unified theory (GUT) scale ($`M_G`$). Moreover the gaugino masses and the trilinear soft breaking terms are also assigned common values $`m_{\frac{1}{2}}`$ and $`A_0`$ respectively, at $`M_G`$. The parameters at the energy scale of interest ($``$ few hundred GeV) is determined by the usual renormalisation group (RG) running.
The number of free parameters may be further reduced by requiring radiative $`SU(2)\times U(1)`$ breaking at the electroweak scale. This fixes the magnitude of the Higgsino mass parameter ($`\mu `$). Thus $`m_0`$, $`m_{\frac{1}{2}}`$, $`A_0`$ along with the sign of $`\mu `$ and $`\mathrm{tan}\beta `$ (the ratio of the vacuum expectation values of the two neutral higgs bosons) define the model completely. This popular model is hereafter referred to as the conventional scenario.
The above framework motivated by N=1 supergravity is very attractive. However as there is no direct experimental information about physics at $`M_G`$, it is imprudent to restrict our attention to this model only. In this paper our goal is to re-examine the prospective SUSY signals at the Tevatron (Run-II) by relaxing some of the above assumptions. We shall, however, assume that the gaugino masses unify at $`M_G`$. This assumption is quite natural within the framework of any SUSY GUT, since it follows if the GUT symmetry is respected by the SUSY breaking mechanism at some high scale.
The assumption of a common soft breaking mass $`m_0`$ at $`M_G`$ is undoubtably more model dependent. Unlike the gauginos different scalars in a SUSY GUT may belong to different representations of the GUT group. This is especially so for the light higgs scalars and the sfermions, which almost always reside in different multiplets. Even if we assume the validity of the supergravity model, the universal parameter $`m_0`$ may well be generated at a scale substantially different from $`M_G`$, say the Plank Scale ($`M_P`$). Then the running of the scalar masses, belonging to different multiplets of the GUT group, between $`M_P`$ and $`M_G`$ may lead to non-universality at $`M_G`$ .
The following nonuniversal scenario is rather interesting from the phenomenological point of view. In this scenario the ‘right - handed’ down - type squarks ($`\stackrel{~}{d}_R,\stackrel{~}{s}_R,\stackrel{~}{b}_R`$ ), generically denoted by $`\stackrel{~}{d}_R`$, are significantly lighter than the other squarks. Then the gluino decays into three body final states mediated by virtual $`\stackrel{~}{d}_R`$ squarks will dominate. Further if the LSP is assumed to be dominated by the $`U(1)`$ gaugino ($`\stackrel{~}{B}`$), then practically all of these virtual $`\stackrel{~}{d}_R`$’s will decay into the LSP and a d-type quark. Thus the branching ratio ( BR) of direct gluino decays into the LSP will be enhanced, while the cascade decays of the gluino will be correspondingly suppressed. In the special case $`m_{\stackrel{~}{g}}>m_{\stackrel{~}{d}_R}`$, while all other squarks are heavier than the gluino, practically all the gluinos will decay into the jets + $`\overline{)}E_T`$ channel with a remarkably hard $`\overline{)}E_T`$ spectrum. On the other hand gluino decays into leptons + jets + $`\overline{)}E_T`$ arising through cascade decays will be strongly suppressed. The signal from $`\stackrel{~}{g}\stackrel{~}{g}`$, $`\stackrel{~}{g}\stackrel{~}{d}_R`$ and $`\stackrel{~}{d}_R\stackrel{~}{d}_R`$ production is likely to be observable, although the other squarks may be heavy to be of any consequence at Tevatron energies .
Theoretically relatively light $`\stackrel{~}{d}_R`$’s can be naturally motivated within a SUSY GUT framework in a variety of ways. If the GUT group is $`SU(5)`$, then the $`\stackrel{~}{d}_R`$ squarks residing in the 5 - plet may be renormalized between $`M_P`$ and $`M_G`$ such that the resulting soft breaking mass at $`M_G`$ is significantly smaller than that of the other squarks belonging to the 10 - plet . The numerical results of ref. , though in the right direction, does not exhibit a large enough mass split.
In this paper we shall illustrate the signatures of a light $`\stackrel{~}{d}_R`$ scenario through an $`SO(10)`$ SUSY GUT to be discussed below. Such models are now much more popular than the good old $`SU(5)`$ SUSY GUT in view of the recent excitement about neutrino masses and mixings generated by the SUPERK and other experiments .
We, however, emphasize that the novel collider signatures are essentially consequences of the above squark - gluino mass hierarchy at low energies and are fairly insensetive to the details of GUT scale or Planck scale physics responsible for generating it. Moreover, in view of the large uncertainties involved in GUT scale - Plank scale physics the quantitative results need not be regarded as firm predictions. Therefore, keeping in mind that either of the above mechanisms or their combination can in principle generate the required mass hierarchy, one might as well discuss the resulting phenomenology in a model independent way.
We shall now focus our attention on an SO(10) SUSY GUT containing all the quarks and leptons of a given generation in a 16 dimensional multiplet which includes the heavy right handed neutrino. In this model the non-universility at $`M_G`$ due to running between $`M_P`$ and $`M_G`$, is expected to be negligible for the first two generations of squarks and sleptons with small Yukawa couplings. In principle nonuniversal masses for the third generation sfermions with a large Yukawa coupling is also possible due to this mechanism. However we shall assume this intergeneration nonuniversility to be small compared to the nonuniversality due to D terms, which will be described below.
Running of the soft breaking masses between $`M_P`$ and $`M_G`$ may result in soft breaking masses of light higgs bosons at $`M_G`$ significantly different from that in the sfermion sector. The light higgs doublets reside in a 10 dimensional representation of SO(10) and hence are renormalised differently. Moreover they have to couple to other super heavy GUT fields in order to implement the mass-split between the coloured higgs bosons and the colour neutral ones responsible for $`SU(2)\times U(1)`$ breaking. Unfortunately the magnitude of the resulting nonuniversility is not calculable without specifying all the couplings of the higgs bosons, which are not known presently. We, therefore, do not attempt to study directly the impact of nonuniversality on higgs phenomenology in this paper. Instead we shall restrict ourselves to the signature of the squark-glunio production and decays which are only weakly dependent on the characteristics of the higgs sector. However, the effect of nonuniversality in the higgs sector will be taken into account indirectly by treating $`\mu `$ as a free parameter.
In summmary we shall work with a SO(10) scenario in which the soft breaking masses of the squarks and sleptons are equal (= $`m_0`$) at $`M_G`$. Non-universality at this scale may still arise due to D-term contributions to the above masses which appear when SO(10) breaks into a group of smaller rank. In general such contributions could be different for different members of the 16-plet . However, these non-universal terms are generation independent, so that no additional problem due to flavour changing neutral currents arise.
As a specific example we shall consider the breaking of SO(10) directly to the SM gauge group . The group $`SO(10)`$ contains $`SU(5)\times U(1)`$ as a subgroup. It is further assumed that the D-terms are linked to the breaking of this $`U(1)`$ only. The squark- slepton masses in this case are
$$m_{\stackrel{~}{u}_L}^2=m_{\stackrel{~}{u}_R}^2=m_{\stackrel{~}{e}_R}^2=m_{0}^{}{}_{}{}^{2}+0.5Dm_0^2$$
(1)
$$m_{\stackrel{~}{d}_R}^2=m_{\stackrel{~}{e}_L}^2=m_{0}^{}{}_{}{}^{2}1.5Dm_0^2,$$
(2)
where the unknown parameter D can be of either sign. The mass differences arise because of the differences in the $`U(1)`$ quantum numbers of the sparticles concerned. As can be readily seen from the above formula for D $`>`$ 0, the left handed sleptons ($`\stackrel{~}{e}_L`$) and right handed down type squark ($`\stackrel{~}{d}_R`$) are relatively light. In this paper we want to concentrate on the collider signatures of the light $`\stackrel{~}{d}_R`$. In principle the D term contributions to the light higgs masses lead to further mass splitting between the higgs bosons and the sfermions at $`M_G`$. This provides additional motivation for treating $`\mu `$ as a free parameter.
The phenomenology of the lighter $`\stackrel{~}{d}_R`$ have been studied by several authors . In this work we shall extend and complement these studies in several ways. In , the production cross-section of squark-glunio pairs and their decay branching ratios were studied. The effects of the kinematical cuts on the resulting SUSY signals, however, were not taken into account. In this paper we study the jets \+ $`\overline{)}E_T`$ as well as opposite sign dileptons + jets + $`\overline{)}E_T`$ signals by using a parton level Monte Carlo. We use as a guide line the kinematical cuts given in , but our main conclusions are essentially consequences of the spectrum under study and are fairly independent of the precise choices of these cuts.
Moreover, we shall comment on the sensitivity of the signal on $`\mu `$ and $`\mathrm{tan}\beta `$. this important point was not addressed in the earlier works. Event generators requiring large amount of computer time, though essential for precise quantitative studies, are rather expensive as tools for studying the dependence of the signal on a large number of parameters. A parton level Monte Carlo on the other hand enables us to carry out a qualitative study relatively easily.
We shall concentrate on two main issues:
a) What are the mass reaches of the upgraded Tevatron in the nonuniversal scenario and how do they compare with that in the conventional scenario ?
b) If a signal is seen at the upgraded Tevatron, can one distinguish between the models with lighter $`\stackrel{~}{d}_R`$ and the conventional scenario ?
The plan of the paper is as follows. In section II we shall discuss regions of the parameter space which are motivated by various theoretical considerations and are interesting for SUSY searches at the Tevatron. In section III we present our result for the jet + $`\overline{)}E_T`$ signal. In section IV the discrimination of different models using the jet + dilepton + $`\overline{)}E_T`$ and the clean trilepton signal is presented. Finally in section V the conclusions are summarised.
## 2 The Choice of Parameters and the Overall Strategy
As has been stated in the introduction the model under study has the following parameters: $`m_0,m_{1/2},A_0,\mu `$, $`\mathrm{tan}\beta `$ and $`D`$. In this set $`m_0`$ and $`m_{\frac{1}{2}}`$ are essentially free parameters.
The glunio mass reach via the jet + $`\overline{)}E_T`$ channel at the upgraded Tevatron has been studied by Baer et al. . Adopting the conventional scenario, their results can be classified into two generic casess: i) squarks much heavier than the gluino ($`m_0>>m_{\frac{1}{2}}`$) and ii) squarks roughly degenerate with the gluinos ($`m_0m_{\frac{1}{2}}`$; squarks much lighter than the gluinos are not allowed in the conventional scenario. Let us review the results in case i) for $`m_0`$ 500 GeV $`>>m_{\frac{1}{2}}`$. It was found that only $`m_{\frac{1}{2}}`$ 75 (100) GeV can be probed at the upgraded Tevatron provided the integrated luminosity accumulates to 2 $`fb^1`$ (25$`fb^1`$) . Unfortunately such low values of $`m_{\frac{1}{2}}`$ have already been ruled out by the direct chargino searches at LEP and direct squark - gluino searches by the D0 collaboration . Thus according to the conventional scenario direct squark gluino searches at the Tevatron in the jet + $`\overline{)}E_T`$ channel will draw a blank if the squarks indeed happen to be very heavy. This motivates us to focus our studies on choice i) in the nonuniversal scenario. In case ii) even the conventional scenario predicts observable signals at the upgraded Tevatron and we shall not consider it further.
It may be worthwhile to mention that recent analyses of the precision elctroweak observables have produced additional evidence, albeit rather mild, in favour of scenario i). In SUSY contributions to several of these observables were studied. Including the contributions from squarks, sleptons, gauginos and higgs bosons seperately, it was found that light squarks or sleptons (with all other sparticless rather heavy) just allowed by the current lower limits from direct searches, always make the fit to 22 data points (Z width and partial widths, various assymmetries etc.) worse than that of the SM. On the otherhand relatively light charginos and neutralinos with heavy sfermions ($`m_0>>m_{\frac{1}{2}}`$) improve the fit although the statistical significance of the improvement is rather modest.
Similar conclusions pertaining to the squark sector were obtained in . However, it was also noted that even for comparatively light sbottoms and small mass of one of the stops, special values of $`\stackrel{~}{t_L}\stackrel{~}{t_R}`$ mixing can make the fit as good as that in the SM.
Increasing the number of theoretical inputs the number of free parameters can be further reduced. Several authors have noted that , if Yukawa coupling unification at $`M_G`$ is demanded for the third generation within an $`SO(10)`$ frame work, then $`\mathrm{tan}\beta `$ becomes practically fixed, since only high values of $`\mathrm{tan}\beta `$ (in a narrow range around 50) lead to such unification. The well known difficulty in accomodating the radiative $`SU(2)\times U(1)`$ breaking in this scenario with universal soft breaking masses for the scalars, may be overcome by the nonuniversality induced by the $`SO(10)`$ D-terms .
We, however, note that the Yukawa coupling unification in SO(10) is a consequence of the assumption that the higgs sector is indeed minimal. In this case a single 10 dimensional higgs multiplet is assumed to contain both the higgs doublets required to generate the masses of the up and down type quarks and to trigger the radiative $`SU(2)\times U(1)`$ breaking. We, therefore, do not require full Yukawa unification for the third generation and the resulting large value of $`\mathrm{tan}\beta `$, since this crucially depends on the choice of the higgs sector.
From the phenomenological point of view the large $`\mathrm{tan}\beta `$ scenario in conjuction with the LEP lower bound $`M_{\stackrel{~}{\chi }_1^\pm }`$ 95 GeV necessarily implies that $`m_{\stackrel{~}{g}}`$ is almost at the edge of or beyond the kinematical reach of the Tevatron collider. Thus direct squark- gluino search at the upgraded Tevatron is of little consequence in this scenario, in particular if the squarks are much heavier than the gluino.
Even if more general higgs multiplets are assumed, b - $`\tau `$ Yukawa unification is a desirable feature of the theory. The conventional wisdom is that this requires values of $`\mathrm{tan}\beta `$ smaller than that in the case of full Yukawa unification. Yet the favoured values of tan $`\beta `$ are still too large to make gluinos sufficiently light to be produced copiously at Tevatron energies. Typical values required by unification are tan $`\beta >`$ 30. However in the presence of neutrino masses and , in particular, of large mixing in the lepton sector this conclusion may require revision .
In the presence of large lepton mixing, as required by the SUPERK data on atmospheric neutrinos , b - $`\tau `$ unification can be achieved for relatively low values of $`\mathrm{tan}\beta `$ which were previously disfavoured. In fact it has been shown in that for $`2`$ tan $`\beta 4`$ and suitable fermion mass matrix textures at $`M_G`$, one can obtain large mixings in the lepton sector along with an acceptable CKM matrix, desired neutrino mass patterns and b - $`\tau `$ unification. From the point of view of Tevatron phenomenology this finding is important, since gluino masses well within the striking range of the Tevatron are not necessarily excluded by the LEP lower bound on $`M_{\stackrel{~}{\chi }_1^\pm }`$ for such low values of tan $`\beta `$. We shall, therefore, restrict ourselves to the above narrow range of $`\mathrm{tan}\beta `$.
The sign of the parameter $`\mu `$ is chosen to be negative since otherwise the gluino mass range allowed by the LEP lower bound on $`M_{\stackrel{~}{\chi }_1^\pm }`$, turns out to be uninteresting for Tevatron phenomenology.
As has been discussed in the introduction, an attractive way of fixing the magnitude of $`\mu `$ is to require radiative $`SU(2)\times U(1)`$ breaking at the electroweak scale. The resulting numerical value, however, strongly depends on the choice of the higgs mass parameter at $`M_G`$. Since we wish to make our predictions largely free from the additional assumptions on the higgs sector we shall treat $`\mu `$ as a free parameter. In the context of Tevatron phenomenology this, however, does not make much of a difference since in any case magnitudes of $`\mu `$ can not be much beyond 450 GeV or so, if we require a gluino well within the striking range of the Tevatron and $`M_{\stackrel{~}{\chi }_1^\pm }>`$ 95 GeV. On the lower side $`\mu `$ is constrained by the requirement of a bino dominated LSP.
We shall denote the cross-section corresponding to the signal with n leptons + jets + $`\overline{)}E_T`$ by $`\sigma _n`$. First we shall consider the jet \+ $`\overline{)}E_T`$ signal ($`\sigma _0`$) arising from squark-gluino production.
In this work we shall reexamine the gluino mass reach at the upgraded Tevatron for large $`m_0`$ ($`>>m_{\frac{1}{2}})`$ in the non-universal scenario. We find that an interesting range of $`m_{\frac{1}{2}}`$ beyond the LEP-2 search limit can be probed. This is particularly so, if a high integrated luminosity ($``$ 30 $`fb^1`$) is available. We further study the distributions of various kinematical observables associated with the final states using conservative kinematial cuts given in and compare and contrast them with the corresponding distribution in the conventional scenario.
The size of the signal is very sensitive to the squark, glunio masses or alternatively with $`m_0`$, D, and $`m_{\frac{1}{2}}`$. The dependence on the magnitude of $`\mu `$ and $`\mathrm{tan}\beta `$ is relatively mild but non trivial. This variation was not studied systematically in earlier works . In this paper we shall check the sensitivity of our conclusions with respect to $`\mu `$ and $`\mathrm{tan}\beta `$.
If the signals for several values of D happen to be indistinguishable, we shall try to distinguish between them by considering the distributions of the final state observables and the corresponding dilepton ($`\sigma _2`$) and clean trilepton signals.
## 3 jet + $`\overline{)}E_T`$
In the conventional scenario Baer et al. have considered the jet + $`\overline{)}E_T`$ signals in great details using the ISAJET-ISASUSY Monte Carlo. They have given the kinematical cuts and the SM background corresponding to these cuts. In our parton level Monte Carlo we have adopted the cuts and the background estimates of Baer et al. Although our numerical estimates based on a simple minded approach give approximate guide lines and should not be treated as firm predictions, the main conclusions drawn are expected to be valid.
Baer et al. have given the jet + $`\overline{)}E_T`$ cross-section for several representative choices of the SUSY parameters (see fig. 3 of ref. ). Since we are interested in the $`m_{\stackrel{~}{q}}>>m_{\stackrel{~}{g}}`$ case, we have focussed our attention on the choice $`m_0`$ = 800 GeV, $`m_{\frac{1}{2}}`$ = 120 GeV, $`\mathrm{tan}\beta `$ = 2, $`A_0`$ = 0 and sign of $`\mu `$ negative. They have also prescribed the following set of kinematical cuts $`E_T(j_1)`$, $`E_T(j_2)>E_{T}^{}{}_{}{}^{c}`$ and $`\overline{)}E_T>E_{T}^{}{}_{}{}^{c}`$, where $`E_{T}^{}{}_{}{}^{c}`$ is a variable which should be chosen appropriately for each point of the parameter space to optimise the signal to background ratio. $`E_T(j_1)`$ and $`E_T(j_2)`$ are the transverse energies of the two leading jets respectively. The other cuts from are $`|\eta _j|`$3 for all jets and $`\mathrm{\Delta }R(\sqrt{\mathrm{\Delta }\eta ^2+\mathrm{\Delta }\varphi ^2})>`$ 0.7. Subject to these cuts the SM background is $``$ 2 $`pb`$. In most of the cases studied in ref. higher values of $`E_{T}^{}{}_{}{}^{c}`$ improves the statistical significance of the signal.
Using these cuts our parton level calculation gives cross-sections which approximately agree with Baer et al. for $`E_{T}^{}{}_{}{}^{c}`$ 50 GeV. For example for $`E_{T}^{}{}_{}{}^{c}`$ = 50 GeV we find $`\sigma _0`$ 35 $`fb`$ where as Baer et al. obtain $``$ 25$`fb`$. A part of the discrepancy ($``$ 10 %) may be attributed to the use of different parton density functions. Baer et al. have used CTEQ2L while we have used CTEQ4M . For higher values of $`E_{T}^{}{}_{}{}^{c}`$ ,however, parton level calculation grossly over estimate the cross section compared to ISAJET result. This is understandable because the reduction in $`p_T`$ of the parton jets due to fragmentation, final state radiation etc. which soften the jet $`p_T`$ in general, is not taken into account in parton level calculations. Being conservative we shall use $`E_{T}^{}{}_{}{}^{c}`$ = 50 GeV. For a reallistic estimate we scale our parton level cross-sections by a factor of 2/3. We however note that our conclusions regarding the search limits are likely to improve to some extent by the use of harder cuts.
We next present the sparticle spectrum for D = i)0.0, ii)0.4, and iii)0.6 using equations (1) and (2). The details are given in Table 1. Our main interest will be restricted to D = 0.6 where $`m_{\stackrel{~}{g}}>m_{\stackrel{~}{d}_R}`$. However, we shall also comment on the D = 0.4 scenario.
| | D=0.0 | D=0.4 | D=0.6 |
| --- | --- | --- | --- |
| $`\stackrel{~}{u}_L`$ | 550 | 593 | 614 |
| $`\stackrel{~}{u}_R`$ | 547 | 591 | 611 |
| $`\stackrel{~}{𝐝}_𝐑`$ | 548 | 390 | 281 |
| $`\stackrel{~}{𝐞}_𝐋`$ | 507 | 328 | 180 |
| $`\stackrel{~}{e}_R`$ | 503 | 550 | 572 |
| $`\stackrel{~}{g}`$ | 312 | 313 | 313 |
| $`\stackrel{~}{\chi _1}^0`$ | 46 | 46 | 46 |
| $`\stackrel{~}{\chi _2}^0`$ | 96 | 95 | 95 |
| $`\stackrel{~}{\chi }_1^\pm `$ | 96 | 95 | 95 |
Table 1: The mass spectrum in GeV at the weak scale for different values of D with $`m_0`$ = 500 GeV, $`m_{\frac{1}{2}}`$= 105 GeV, $`\mathrm{tan}\beta `$ = 3, $`A_0`$ = 0, $`\mu `$ = $``$340 GeV.
Further using the radiative $`SU(2)\times U(1)`$ breaking we find $`\mu `$ -340 GeV in the universal scenario (D = 0.0). In principle $`\mu `$ can be determined from radiative $`SU(2)\times U(1)`$ breaking in the non universal scenario as well, if we make additional assumptions about the higgs masses at $`M_G`$. We shall, however, refrain from making such assumptions and, as has already been mentioned in the introduction, treat $`\mu `$ as a free parameter. In order to study the impact of light $`\stackrel{~}{d}_R`$ squarks on the jets + $`\overline{)}E_T`$ signal we shall first use $`\mu `$ = $``$340 GeV even in the nonuniversal case. Later we shall comment on the sensitivity of the signal to $`\mu `$.
For Tevatron Run II ($``$ 2 $`fb^1`$), where $``$ indicates the integrated luminosity, we find that for $`m_{\frac{1}{2}}`$ 100 GeV, no observable signal is expected for D = 0.0 in agreement with Baer et al. For D=0.4 the conclusion remains more or less unchanged. For D=0.6 ($`m_{\stackrel{~}{g}}>m_{\stackrel{~}{d}_R}`$), however, a signal may be seen, provided $`m_{\frac{1}{2}}`$ is in a narrow range just beyond the LEP II limit. For example we find for $`m_{\frac{1}{2}}`$ = 110 GeV, $`\sigma _0`$ = 143 $`fb`$ which corresponds to $`\frac{\sigma }{\sqrt{B}}`$ 5.
This may be understood from the following facts. The gluino decay channels and corresponding branching ratio (BR)s are given in the Table 2.
| | D=0.0 | D=0.4 | D=0.6 |
| --- | --- | --- | --- |
| $`\stackrel{~}{g}\stackrel{~}{\chi }^\pm +jets`$ | 0.4657 | 0.3681 | 0.0006 |
| $`\stackrel{~}{g}\stackrel{~}{\chi _1}^0+jets`$ | 0.1418 | 0.3311 | 0.0001 |
| $`\stackrel{~}{g}\stackrel{~}{\chi _2}^0+jets`$ | 0.3925 | 0.3009 | 0.0005 |
| $`\stackrel{~}{g}\stackrel{~}{d}_R+jet`$ | 0.0 | 0.0 | 0.998 |
Table 2: The glunio decay channels and corresponding branching ratios for different values of D.
For D = 0 case 3-body decay of the glunio dominates because all the squarks ($`\stackrel{~}{q}_L`$, $`\stackrel{~}{q}_R`$) are heavier than the glunio. $`\stackrel{~}{\chi }^\pm `$ and $`\stackrel{~}{\chi _2}^0`$ decay through leptonic as well as hadronic modes. So BR($`\stackrel{~}{g}jets+\overline{)}E_T`$) is somewhat suppressed. But for D = 0.4 this BR increases. This is due to the light $`\stackrel{~}{d}_R`$ propagator which is less suppressed. As a result BR($`\stackrel{~}{g}\stackrel{~}{\chi _1}^0+jets`$) is enhanced significantly (0.14 $``$ 0.33).
For D=0.6, BR ($`\stackrel{~}{g}jets+\overline{)}E_T`$) increases rapidly. In this case $`m_{\stackrel{~}{g}}>m_{\stackrel{~}{d}_R}`$ and so the 2-body decay of glunio dominates. First we have the decay $`\stackrel{~}{g}\stackrel{~}{d}_R`$ d ($`\stackrel{~}{d}_R\stackrel{~}{d}_R`$, $`\stackrel{~}{s}_R`$, $`\stackrel{~}{b}_R`$) with BR $`=`$ 0.998, followed by $`\stackrel{~}{d}_R\stackrel{~}{\chi _1}^0d`$ with 100 % BR as the $`\stackrel{~}{\chi _1}^0`$ is $`\stackrel{~}{B}`$ dominated.
The observability of jets + $`\overline{)}E_T`$ signal may improve sgnificantly if higher $`(`$ 30 $`fb^1)`$ is available at the upgraded Tevatron (see Table 3).
| D=0.0 | | | D=0.4 | | | D=0.6 | | |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`m_{\frac{1}{2}}`$ | $`\sigma _0`$ | $`\frac{S}{\sqrt{B}}`$ | $`m_{\frac{1}{2}}`$ | $`\sigma _0`$ | $`\frac{S}{\sqrt{B}}`$ | $`m_{\frac{1}{2}}`$ | $`\sigma _0`$ | $`\frac{S}{\sqrt{B}}`$ |
| 105 | 53 | 6 | 105 | 62 | 7 | 105 | 203 | 25 |
| - | - | - | - | - | - | 125 | 53 | 6 |
Table 3: The $`jets+\overline{)}E_T`$ cross-sections (in fb) and statistical significances for different values of $`m_{\frac{1}{2}}`$ and D with $`m_0`$ = 500 GeV, $`\mathrm{tan}\beta `$ = 3, $`A_0`$ = 0, $`\mu `$ = $``$340 GeV.
We find that if the chargino mass is just above the LEP lower limit corresponding to $`m_{\stackrel{~}{g}}`$ 312 GeV, a signal may also be expected for D = 0.0 and 0.4. For the D = 0.0 case we obtain slightly enhanced $`\frac{S}{\sqrt{B}}`$ ratio compared to ref. since we have used $`\mathrm{tan}\beta `$ = 3. and $`(`$ 30 $`fb^1)`$. For heavier charginos no signal is anticipated.
For D = 0.6, however, a range of $`m_{\frac{1}{2}}`$ 105 GeV $`m_{\frac{1}{2}}`$ 125 GeV can be probed. This is the consequence of the production of relatively light $`\stackrel{~}{d}_R`$ squarks along with the gluino and the enhanced BR ($`\stackrel{~}{g}jets+\overline{)}E_T`$) for reasons discussed above.
We next study the variation of $`\sigma _0`$ with $`\mu `$. For D = 0.6 the results hardly changes with $`\mu `$. This is a consequence of the fact that in this case the decays $`\stackrel{~}{g}\stackrel{~}{d}_Rd`$ and $`\stackrel{~}{d}_R\stackrel{~}{\chi _1}^0d`$ dominates the signal. The branching ratio of the former strong decay is insensitive to $`\mu `$. The second decay has $``$100 % BR as long as the LSP is $`\stackrel{~}{B}`$ dominated.
For D = 0.0 and D = 0.4 the signal has some dependence on $`\mu `$ (see Table 4).
| $`\mu `$ | $`\sigma _0`$(D = 0.0) | $`\sigma _0`$(D = 0.4) |
| --- | --- | --- |
| $``$340 | 79 | 93 |
| $``$400 | 82 | 80.5 |
| $``$450 | 81 | 72 |
| $``$500 | 80 | 66 |
| $``$600 | 73 | 63.5 |
| $``$700 | 65.5 | 58 |
Table 4: The sensitivities of the $`jets+\overline{)}E_T`$ cross-section with $`\mu `$ for different values of D. All cross-sections are in fb and $`\mu `$ in GeV.
However for $`|\mu |>`$ 450 GeV, the chargino mass violates the LEP lower bound. Larger values of $`\mu `$, therefore, require enhanced $`m_{\frac{1}{2}}`$ which makes the $`m_{\stackrel{~}{g}}`$ larger and the signal at the Tevatron is supressed below the observable limit.
We next study the variation of $`\sigma _0`$ with $`\mathrm{tan}\beta `$. It is once again found that $`\sigma _0`$ for D = 0.6 is not at all sensitive to this parameter. For D = 0.4 and D = 0.0 the signal show some sensitivity. However, for $`\mathrm{tan}\beta `$ 4, the chargino mass again violates the LEP lower bound unless $`m_{\frac{1}{2}}`$ and correspondingly $`m_{\stackrel{~}{g}}`$ is increased.
Fig. 1: Missing $`E_T`$ distribution of signal for three different values of $`D`$-parameter.
From tables 3 and 4 it follows that the magnitudes of $`\sigma _0`$ is approximately the same in the following cases($`m_0`$ = 500, $`A_0`$ = 0, $`\mathrm{tan}\beta `$ = 3):
a) $`m_{\frac{1}{2}}`$ = 105 GeV, D = 0.0, $`\mu `$ =$``$340 GeV
b) $`m_{\frac{1}{2}}`$ = 105 GeV, D = 0.4, $`\mu `$ = $``$400 GeV
c) $`m_{\frac{1}{2}}`$ = 105 GeV, D = 0.6, $`\mu `$ = $``$340 GeV
We now explore the possiblity of distinguishing between these different scenarios by using the $`\overline{)}E_T`$ distribution and the $`p_T`$ distributions of the leading jet. In fig.1 we present the $`\overline{)}E_T`$ distribution with the kinematical cuts given above. We find that already in D = 0.4 case the distribution is considerably harder than that in the mSUGRA scenario at least in the interval 150 GeV $`\overline{)}E_T`$ 300 GeV. A bin by bin analysis of this distribution may disentangle the two models, if a signal is seen.
Fig. 2: Leading-jet $`P_T`$ distribution of signal for three different values of $`D`$-parameter.
Fig. 3: Mass reach of jets $`+\overline{)}E_T`$ signal in $`m_0m_{\frac{1}{2}}`$ plane.
For the D = 0.6 case the signal has a very hard $`\overline{)}E_T`$ spectrum extending far beyond the distributions in the other two cases.
The $`p_T`$ spectrum of the leading jet is presented in fig.2. It is clear that no distinction between parameter set a) and b) are possible. A much harder $`p_T`$ spectrum for the leading jet, on the other hand, can easily identify the $`m_{\stackrel{~}{g}}>m_{\stackrel{~}{d}_R}`$ scenario (D = 0.6).
We next investigate the variation of $`\sigma _0`$ with $`m_0`$. As expected $`m_{\stackrel{~}{d}_R}`$ increases with $`m_0`$ and beyond a certain range we find $`m_{\stackrel{~}{d}_R}>m_{\stackrel{~}{g}}`$. Upto $`m_0`$ 700 GeV a 5$`\sigma `$ signal can be obtained for values of $`m_{\frac{1}{2}}`$ beyond the LEP limit (see fig.3).
## 4 OS-dileptons and clean Trileptons
We have also studied the consequence of nonuniversality in the opposite sign(OS) dilepton + jets + $`\overline{)}E_T`$ channel. For D=0.6, $`m_0`$ = 500 GeV, $`m_{\frac{1}{2}}`$ = 105 GeV, $`\mathrm{tan}\beta `$ = 3, $`\mu <0`$, (i.e. $`m_{\stackrel{~}{g}}>m_{\stackrel{~}{d}_R}`$) the OS-dilepton cross-section vanishes since there is no cascade decay of the glunios. This coupled with a relatively large $`\sigma _0`$ makes this model totally distinct from the D = 0.0 case.
Baer et al. have studied OS dilepton signal in great details. We have used the following cuts from : $`E_T(j_1)`$, $`E_T(j_2)`$ 50 GeV, $`\overline{)}E_T`$ 50 GeV, $`E_T(l_1)`$, $`E_T(l_2)`$ 10 GeV, $`|\eta _l|<`$ 2.5, and the lepton isolation criterion $`R>`$ 0.3.
Using the above cuts we get for $`m_0`$ = 800 GeV, $`m_{\frac{1}{2}}`$ = 120 GeV, $`\mathrm{tan}\beta `$ = 2, $`\mu <`$ 0, $`\sigma _2`$ = 1.8 $`fb`$ where as Baer et al. have obtained $`\sigma _2`$ 2 $`fb`$. Thus the aggrement is rather well. With harder cuts the signal/background ratio improves. But as discussed above realistic e2stimates may not be obtained from a parton level Monte Carlo with such strong cuts.
For $`m_0`$ = 500 GeV, $`m_{\frac{1}{2}}`$ = 105 GeV and for $`\mathrm{tan}\beta `$ 2 Baer et al. have already analysed $`\sigma _2`$ in the conventional scenario. Their conclusion was that this signal is unobservable at the upgraded Tevatron. However, in their analyses they asumed $`|\mu |`$ to be a fixed number determined by the $`SU(2)\times U(1)`$ breaking condition, which may not be realistic due reasons discussed earlier.
We have re-examined the OS signal at the parton level for D = 0.0 case treating $`\mu `$ as a free parameter. We find that irrespective of the value of $`\mu `$ the cross section is unobservable even with $``$ = 30 $`fb^1`$ in the conventional scenario.
But for $`m_0`$=500 GeV, $`m_{\frac{1}{2}}`$ = 105 GeV , $`\mathrm{tan}\beta `$ = 3 and D = 0.4 an observable signal may be achieved for favourable values of $`\mu `$ (see Table 5).
| $`\mu `$ | $`\sigma _2`$ | BR($`\stackrel{~}{\chi _2}^0l^+l^{}\stackrel{~}{\chi _1}^0`$) | $`\frac{S}{\sqrt{B}}`$ |
| --- | --- | --- | --- |
| $``$340 | 5.1 | 0.08 | 3 |
| $``$400 | 7.4 | 0.12 | 4 |
| $``$450 | 8.6 | 0.14 | 5 |
Table 5: The sensitivities of the $`OSdilepton+jets+\overline{)}E_T`$ cross sections with $`\mu `$. Cross-sections are in fb and $`\mu `$ in GeV.
It is clear from the Table 5 that in the non-universal case OS dilepton + $`\overline{)}E_T`$ may be observable for large $`|\mu |`$, using $``$ = 30 $`fb^1`$. This happens since $`|\mu |`$ increases, $`\stackrel{~}{\chi _2}^0`$ and $`\stackrel{~}{\chi _1}^0`$ becomes more gaugino like and as BR ($`\stackrel{~}{\chi _2}^0l^+l^{}\stackrel{~}{\chi _1}^0`$) increases. Incidentally the OS dilepton signal may be a convenient tool for distinguishing the D = 0.4 scenario (parameter set in section 3) from the others which predicts unobservable dilepton signals. The most appropriate tool for distinguishing the three parameter sets a, b, c presented in the previous section is ,however, the clean trilepton signal .
We again use the cuts of ref; $`|\eta _l|<`$ 2.5, $`p_T(l_1)>`$ 20 GeV, $`p_T(l_2)>`$ 15 GeV, $`p_T(l_3)>`$ 10 GeV, $`\overline{)}E_T>`$ 25 GeV and $`|m(l\overline{l})M_Z|`$ 10 GeV. In recent times it has been emphasised that the background from $`W\gamma ^{}/Z^{}`$ is the most severe one in the channel . In order to take care of this background we have introduced an additional invariant mass cut of $`m_{l\overline{l}}>`$ 12 GeV. Using MADGRAPH we estimate the SM background to be $``$ 5 $`fb`$ subject to the above cuts. The clean 3l cross -section is presented in Table 6.
| parameter | $`\sigma _{\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi _2}^0}`$ | BR($`\stackrel{~}{\chi _2}^0l^+l^{})`$ | $`\sigma _{3l}`$ | $`\frac{S}{\sqrt{B}}`$ |
| --- | --- | --- | --- | --- |
| set | in $`pb`$ | | in $`fb`$ | |
| a | 2.294 | 0.019 | 3.0 | 7 |
| b | 2.371 | 0.117 | 16.7 | 41 |
| c | 1.187 | 0.152 | 32.3 | 79 |
Table 6: The clean trilepton cross-sections in pb for different set of parameters.
From Table 6 it follows that the three scenarios can be conveniently distinguished by the clean 3l signal.
## 5 Conclusions
Within the framework of $`N=1`$ SUGRA, it is quite possible that the $`\stackrel{~}{d}_R`$ squarks are significantly lighter than all other sfermions. We consider the case in which $`\stackrel{~}{d}_R`$ squarks have mass $`\stackrel{<}{}`$ $`m_{\stackrel{~}{g}}`$, while all other squarks are significantly heavier than the gluino, and outside the kinematical reach of the Tevatron. These light $`\stackrel{~}{d}_R`$ squarks have several distinctive features in comparison with the conventional MSUGRA scenario : 1) enhancement of $`j`$ $`+`$ $`\overline{)}E_T`$ cross-section, 2) suppression of multilepton $`+`$ $`j`$ $`+`$ $`\overline{)}E_T`$ cross-section and the cascade decays of the gluinos and 3)relatively hard missing energy spectrum.
Although in view of various uncertainties in Planck and GUT scale physics it is desirable to consider the scenario in a model independant way, we have considered a model based on $`SO(10)`$ $`D`$-terms to generate the mass spectra for the purpose of illustration.This model has only one extra parameter, namely the $`D`$-parameter, than the conventional MSUGRA model.
For $`m_0`$ $`>>`$ $`m_{\frac{1}{2}}`$, which makes most of the squarks much heavier than the gluino, it is well known that MSUGRA does not yield an observable $`j`$ $`+`$ $`\overline{)}E_T`$ signal at the upgraded Tevatron for gluino masses consistent with LEP or Tevatron Run-I lower bounds . On the other hand if $`m_{\stackrel{~}{d}_R}<m_{\stackrel{~}{g}}`$, an observable signal can be seen even with anintegrated luminosity $``$ $`2fb^1`$ for $`m_{\frac{1}{2}}`$ just above the current lower bound. As the integratde luminosity accumulates to $``$ $`30`$ $`fb^1`$ at the upgraded Tevatron, asignificant range ( $`105`$ GeV $`m_{\frac{1}{2}}125`$ GeV) can be probed. The variation of this signal with $`\mathrm{tan}\beta `$ and $`\mu `$ is insignificant for reasons discussed in the text <sup>4</sup><sup>4</sup>4 Variation of $`\mu `$ is an indirect way of taking into account the larger uncertainties in the higgs sector of the model.. Although we have carried out most of the calculations for $`m_o`$ = $`500`$ GeV, similar signals can be achieved for any $`m_o`$ $`<`$ $`700`$ GeV, even if $`m_o`$ $`>>`$ $`m_{\frac{1}{2}}`$.
We have also considered the possibilities of distinguishing between the various scenarios from the $`j+\overline{)}E_T`$ signal at the Tevatron. We illustrate this with three values of the $`D`$-parameter: 1) $`D=0`$ (conventional MSUGRA) 2) $`D=0.4`$ ($`m_{\stackrel{~}{d}_R}>m_{\stackrel{~}{g}}`$, but $`<<m_{\stackrel{~}{q}}`$) and 3) $`D=0.6`$ ($`m_{\stackrel{~}{g}}>m_{\stackrel{~}{d}_R}`$. We have chosen all other parameters such that the $`j+\overline{)}E_T`$ cross-sections are comparable in all the three cases. As discussed in th text the missing energy spetrum in scenario 2) is already somewhat harder compared to that in scenario 1). This difference may observable depending on the value of the integrated luminosity. The $`\overline{)}E_T`$ spectrum in scenario 3) is much harder compared to that in 1) and 2) and extends far beyond the end point of the corresponding distributions and can be easily distinguished.
Multilepton $`+j+\overline{)}E_T`$ signals may also help to distinguish between the three scenarios. In the scenario 3) the OS di-lepton $`+j+\overline{)}E_T`$ is absent. In scenario 1) a slight enhancement above the SM background is possible for all values of $`\mu `$, although this enhancement is not statistically significant to qualify as a genuine SUSY signal. In scenario 2) for suitable values of $`\mu `$ the OS di-lepton signal may be so large that it may be above the SM background in a staistically significant way.
Finally clean tri-lepton signal differ quite appreciably in the three cases. It is the largest in case 3 while in case 2 it is still muh larger than that in the conventional scenario.
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# Magnetogenesis and the dynamics of internal dimensions
## I Introduction
The past dynamical history of our Universe is not completely known and the theoretical understanding is often guided by consistency rather than by observational evidence. The remarkable similarity of the abundances of light elements in different galaxies leads to postulate that the Universe had to be dominated by radiation at the moment when the light elements were formed, namely for temperatures of approximately $`0.1`$ MeV . Prior to the moment of nucleosynthesis even indirect informations concerning the thermodynamical state of our Universe are lacking even if our knowledge of particle physics could give us important hints concerning the dynamics of the electroweak phase transition .
The success of big-bang nucleosynthesis (BBN) sets limits on alternative cosmological scenarios. Departures from homogeneity and isotropy of the background geometry can be successfully constrained. Bounds on the presence of matter–antimatter domains of various sizes can be derived . BBN can also set limits on the dynamical evolution of internal dimensions . Internal dimensions are an essential ingredient of theories attempting the unification of gravitational and gauge interactions in a higher dimensional background like Kaluza-Klein theories and superstring theories .
Defining, respectively, $`b_{BBN}`$ and $`b_0`$ as the size of the internal dimensions at the BBN time and at the present epoch, the maximal variation allowed to the internal scale factor from the BBN time can be expressed as $`b_{BBN}/b_01+ϵ`$ where $`|ϵ|<10^2`$ . The bounds on the variation of the internal dimensions during the matter dominated epoch are even stronger. Denoting with an over-dot the derivation with respect to the cosmic time coordinate, we have that $`|\dot{b}/b|<10^9H_0`$ where $`H_0`$ is the present value of the Hubble parameter . The fact that the time evolution of internal dimensions is so tightly constrained for temperatures lower of $`1`$ MeV does not forbid that they could have been dynamical prior to that epoch.
An apparently unrelated observational evidence characterizing the present Universe is the existence of large scale magnetic fields . Faraday rotation measurements, Zeeman splitting estimates (when available) and synchrotron emission patterns indicate that distant galaxies seem to be endowed with a magnetic field of roughly the same strength of the one of the Milky Way which is of the order of $`10^6`$ G. The similarity of the magnetic field strength in different galaxies led to the feeling that the origin of these large scale fields should somehow be connected with the cosmological consequences of the interplay between gravitational and gauge interactions . Observations of magnetic fields at even larger scales (i.e cluster , inter-cluster) are still under debate mainly because of the rather problematic estimates of the electron density in the inter-galactic medium. The existence of strong magnetic fields with coherence scale larger than the galactic one can be of crucial importance for the propagation of high-energy cosmic rays.
The measured large scale (galactic) magnetic fields are assumed to be the result of the exponential amplification (due to galactic rotation) of some primeval seed fields whose typical value lies around $`10^{16}`$$`10^{25}`$ G at the decoupling epoch. Several mechanisms have been invoked in order to explain the origin of these seeds . Large scale magnetic fields can also have relevant physical implications in other (related) areas of cosmology and especially in connection with polarization (and distortion) of the Cosmic Microwave Background (CMB) .
In the present paper it is argued that there may be a relation between the existence of large scale magnetic fields and the possible occurrence of a primordial phase of the Universe where the internal dimensions have been dynamical. In short the logic goes as follows. Suppose that prior to BBN internal dimensions were evolving in time and assume, for sake of simplicity, that after BBN the internal dimensions have been frozen to their present (constant) value. The evolution in time of a classical cosmological background can amplify a given distribution of (initially small) inhomogeneities of the metric and of (non conformally coupled) matter fields . If the background geometry is isotropic and conformally flat Abelian gauge fields cannot be amplified since their equations of motion are invariant under a (Weyl) rescaling of the metric tensor. However, if the internal dimensions change in time the Weyl invariance of Maxwell equations is naturally broken and electromagnetic (vacuum) fluctuations can be amplified.
The purpose of the present investigation is to clarify if a suitable dynamics of the internal dimensions can produce sizable seed fields which could turn on the galactic dynamo mechanism and have other indirect effects during the evolution of the Universe. The plan of the paper is then the following. In Section II we will discuss our basic equations. In Section III we will perform a model independent analysis of the amplification of magnetic fields from monotonically evolving internal dimensions. In Section IV some specific models of internal evolution will be studied. Section V contains the concluding remarks
## II Basic Equations
Consider a homogeneous and anisotropic manifold whose line element can be written as
$`ds^2=G_{\mu \nu }dx^\mu dx^\nu =a^2(\eta )[d\eta ^2\gamma _{ij}dx^idx^j]b^2(\eta )\gamma _{ab}dy^ady^b,`$ (2.1)
$`\mu ,\nu =0,\mathrm{},D1=d+n,i,j=1,\mathrm{},d,a,b=d+1,\mathrm{},d+n.`$ (2.2)
\[$`\eta `$ is the conformal time coordinate related, as usual to the cosmic time $`t=a(\eta )𝑑\eta `$ ; $`\gamma _{ij}(x)`$, $`\gamma _{ab}(y)`$ are the metric tensors of two maximally symmetric Euclidean manifolds parameterized, respectively, by the “internal” and the “external” coordinates $`\{x^i\}`$ and $`\{y^a\}`$\]. The metric of Eq. (2.2) describes the situation in which the $`d`$ external dimensions (evolving with scale factor $`a(\eta )`$) and the $`n`$ internal ones (evolving with scale factor $`b(\eta )`$) are dynamically decoupled from each other . The results of the present investigation, however, can be easily generalized to the case of $`n`$ different scale factors in the internal manifold.
Consider now a pure electromagnetic fluctuation decoupled from the sources, representing an electromagnetic wave propagating in the $`d`$-dimensional external space such that $`A_\mu A_\mu (\stackrel{}{x},\eta )`$, $`A_a=0`$. In the metric given in Eq. (2.2) the evolution equation of the gauge field fluctuations can be written as
$$\frac{1}{\sqrt{G}}_\mu \left(\sqrt{G}G^{\alpha \mu }G^{\beta \nu }F_{\alpha \beta }\right)=0,$$
(2.3)
where $`F_{\alpha \beta }=_{[\alpha }A_{\beta ]}`$ is the gauge field strength and $`G`$ is the determinant of the $`D`$ dimensional metric. Notice that if $`n=0`$ the space-time is isotropic and, therefore, the Maxwell’s equations can be reduced (by trivial rescaling) to the flat space equations. If $`n0`$ we have that the evolution equation of the electromagnetic fluctuations propagating in the external $`d`$-dimensional manifold will receive a contribution from the internal dimensions which cannot be rescaled away Notice that the electromagnetic field couples only to the internal dimensions through the determinant of the $`D`$-dimensional metric. In string theories, quite generically, the one-form fields are also coupled to the dilaton field. This case has been already analyzed in the context of string inspired cosmological scenarios . . In the radiation gauge ($`A_0=0`$ and $`_iA^i=0`$) For a discussion of gauges in curved spaces see the evolution the vector potentials can be written as
$$A_i^{\prime \prime }+nA_i^{}\stackrel{}{}^2A_i=0,=\frac{b^{}}{b}.$$
(2.4)
The vector potentials $`A_i`$ are already rescaled with respect to the (conformally flat) $`d+1`$ dimensional metric. In terms of the canonical normal modes of oscillations $`𝒜_i=b^{n/2}A_i`$ the previous equation can be written in a simpler form, namely
$$𝒜_i^{\prime \prime }V(\eta )𝒜_i\stackrel{}{}^2𝒜_i=0,V(\eta )=\frac{n^2}{4}^2+\frac{n}{2}^{}.$$
(2.5)
In order to estimate the amplification of the gauge fields induced by the evolution of the internal geometry we shall consider the background metric of Eq. (2.2) in the case of maximally symmetric subspaces $`\gamma _{ij}=\delta _{ij}`$, $`\gamma _{ab}=\delta _{ab}`$.
Suppose now that the background geometry evolves along three different epochs. During the first phase (taking place for $`]\mathrm{},\eta _1]`$) the evolution is truly multidimensional. At $`\eta =\eta _1`$ the multidimensional dynamics is continuously matched to a radiation dominated phase turning, after decoupling, into a matter dominated regime of expansion. During the radiation and matter dominated stages the internal dimensions are fixed to their (present) constant size in order not to conflict with possible bounds arising both at the BBN time and during the matter-dominated epoch. The evolution of the external dimensions does not affect the amplification of the gauge fields as it can be argued from Eq. (2.5) :in the limit $`n0`$ (i.e. conformally invariant background) Eq. (2.5) reduces to the flat space equation.
A model of background evolution can be generically written as
$`a(\eta )=a_1\left({\displaystyle \frac{\eta }{\eta _1}}\right)^\sigma ,b(\eta )=b_1\left({\displaystyle \frac{\eta }{\eta _1}}\right)^\lambda ,\eta <\eta _1,`$ (2.6)
$`a(\eta )=a_1\left({\displaystyle \frac{\eta +2\eta _1}{\eta _1}}\right),b(\eta )=b_1,\eta _1\eta \eta _2,`$ (2.7)
$`a(\eta )=a_1{\displaystyle \frac{(\eta +\eta _2+4\eta _1)^2}{4\eta _1(\eta _2+2\eta _1)}},b(\eta )=b_1,\eta >\eta _2.`$ (2.8)
In the parameterization of Eq. (2.8) the internal dimensions grow (in conformal time) for $`\lambda <0`$ and they shrink for $`\lambda >0`$ <sup>§</sup><sup>§</sup>§To assume that the internal dimensions are constant during the radiation and matter dominated epoch is not strictly necessary. If the internal dimensions have a time variation during the radiation phase we must anyway impose the BBN bounds on their variation . The tiny variation allowed by BBN implies that $`b(\eta )`$ must be effectively constant for practical purposes. . By inserting this background into Eq. (2.5) we obtain that for $`\eta <\eta _1`$
$$V(\eta )=\frac{n\lambda }{4\eta ^2}(n\lambda 2),$$
(2.9)
whereas $`V(\eta )0`$ for $`\eta >\eta _1`$. Since $`V(\eta )`$ goes to zero for $`\eta \pm \mathrm{}`$ we can define, in both limits, a Fourier expansion of $`𝒜_i`$ in terms of two distinct orthonormal sets of modes. If we promote the classical fields to quantum mechanical operators in the Heisenberg representation we can write, for $`\eta \mathrm{}`$
$$\widehat{𝒜}_i^{\mathrm{in}}(\stackrel{}{x},\eta )=\frac{d^3k}{(2\pi )^{3/2}}\underset{\alpha }{}e_i^\alpha (\stackrel{}{k})\left[a_{k,\alpha }\varphi _k(\eta )e^{i\stackrel{}{k}\stackrel{}{x}}+a_{k,\alpha }^{}\varphi _k^{}(\eta )e^{i\stackrel{}{k}\stackrel{}{x}}\right],$$
(2.10)
where the sum runs over the physical polarizations. For $`\eta +\mathrm{}`$ $`𝒜_i`$ can be expanded in a second orthonormal set of modes
$$\widehat{𝒜}_i^{\mathrm{out}}(\stackrel{}{x},\eta )=\frac{d^3k}{(2\pi )^{3/2}}\underset{\alpha }{}e_i^\alpha (\stackrel{}{k})\left[\stackrel{~}{a}_{k,\alpha }\psi _k(\eta )e^{i\stackrel{}{k}\stackrel{}{x}}+\stackrel{~}{a}_{k,\alpha }^{}\psi _k^{}(\eta )e^{i\stackrel{}{k}\stackrel{}{x}}\right].$$
(2.11)
Since both sets of modes are complete the old modes can be expressed in terms of the new ones
$$\varphi _k(\eta )=c_+(k)\psi _k(\eta )+c_{}(k)\psi _k^{}(\eta ).$$
(2.12)
If we insert this last equation back into Eq. (2.10) we get
$$\stackrel{~}{a}_k=c_+(k)a_k+c_{}(k)^{}a_k^{}.$$
(2.13)
The commutation relations within the “in” and “out” sets of orthonormal expansions are preserved if the two complex numbers $`c_+(k)`$ and $`c_{}(k)`$ (the so-called Bogoliubov coefficients) are subjected to the constraints $`|c_+(k)|^2|c_{}(k)|^2=1`$ so that the (unitary) connection between the two asymptotic vacua is parametrized, overall, by three real numbers. If the continuity of the field operators in $`\eta _1`$ is imposed $`c_\pm (k)`$ can be determined. The continuity between the two asymptotic forms of the field operators is ensured provided the old mode functions and their first derivatives are continuously matched to the new ones. The evolution equation satisfied by the mode functions during the multidimensional phase
$$\frac{d^2\varphi _k}{d\eta ^2}+\left[k^2V(\eta )\right]\varphi _k=0,$$
(2.14)
can be exactly solved
$$\varphi _k(\eta )=\frac{p}{\sqrt{2k}}\sqrt{k\eta }H_\nu ^{(2)}(k\eta ),p=\sqrt{\frac{\pi }{2}}e^{i\frac{\pi }{4}(1+2\nu )},$$
(2.15)
where $`H_\nu ^{(2)}`$ is simply the Hankel function of second kind . For $`\eta >\eta _1`$ the solution of Eq. (2.5) can be simply written as
$$\varphi _k(\eta )=c_+\psi _k(\eta )+\psi _k^{}(\eta ),\psi _k(\eta )=\frac{1}{\sqrt{2k}}e^{ik(\eta +\eta _1)}.$$
(2.16)
The frequency mixing coefficient $`c_{}(k)`$ determines the spectral number of produced particles in a given mode of the field which turns out to be $`n_k|c_{}(k)|^2`$. By continuously matching, in $`\eta =\eta _1`$, the solutions expressed by Eqs. (2.15) and (2.16) the explicit form of the mixing coefficients can be obtained
$$c_{}=\frac{p}{2}\{H_\nu ^{(2)}(k\eta _1)[\sqrt{k\eta _1}\frac{i}{\sqrt{k\eta _1}}(\nu +\frac{1}{2})]\pm i\sqrt{k\eta _1}H_{\nu +1}^{(2)}(k\eta _1)\}.$$
(2.17)
The mixing coefficients determined in the approximation of a sudden change of the background geometry (taking place in $`\eta _1`$) lead to an ultraviolet divergence in the number and in the energy density of produced particles . In fact for modes of comoving frequency $`k^2`$ larger than the height of the potential barrier $`V(\eta )`$ \[see Eq. (2.9)\] the sudden approximation is not adequate. The mixing coefficients, in the regime $`k>\sqrt{V}`$, should be computed using a smooth function interpolating between the two regimes. This standard analysis leads to a number of particle which is exponentially suppressed as $`\mathrm{exp}[qk/k_1]`$ where $`k_1\eta _1^1`$ and $`q`$ is a number of order one. Thus, the frequency mixing of modes which never hit the potential barrier will be approximately neglected since, for these modes, $`c_+(k)1`$ and $`c_{}(k)0`$. The modes experiencing the amplification are the ones for which $`k<\sqrt{V}`$. Since the maximal height of the barrier is set by the time scale of the transition (i.e. $`V(\eta _1)\eta _1^2`$) the results of Eq. (2.17) can be safely estimated in the small argument limit namely for $`|k\eta _1|<1`$ of the corresponding Hankel functions.
Using the explicit form of the field operators the two-point correlation function of the magnetic field fluctuations can be computed
$$𝒢_{ij}(\stackrel{}{r},\eta )=0_k|\widehat{B}_i(\stackrel{}{x},\eta )\widehat{B}_j(\stackrel{}{x}+\stackrel{}{r},\eta )|0_k=d^3k𝒢_{ij}(k)e^{i\stackrel{}{k}\stackrel{}{r}},$$
(2.18)
with
$$𝒢_{ij}(k)=\frac{𝒦_{ij}}{(2\pi )^3}𝒞(\nu )k|k\eta _1|^{12\nu },\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}\nu =|n\lambda 1|,$$
(2.19)
and where
$`𝒦_{ij}={\displaystyle \underset{\alpha }{}}e_i^\alpha (k)e_j^\alpha (k)=\left(\delta _{ij}{\displaystyle \frac{k_ik_j}{k^2}}\right),`$ (2.20)
$`𝒞(\nu )={\displaystyle \frac{2^{2\nu 3}}{\pi }}\mathrm{\Gamma }^2(\nu )\left({\displaystyle \frac{1}{2}}\nu \right)^2.`$ (2.21)
The magnetic energy density can be obtained by tracing over the physical polarizations with the result that
$$\rho _B(r)=\rho _B(k)\frac{\mathrm{sin}kr}{kr}\frac{dk}{k},$$
(2.22)
where,
$$\rho _B(k)=\frac{𝒞(\nu )}{\pi ^2}k^4|k\eta _1|^{12\nu },$$
(2.23)
is the logarithmic energy spectrum of the magnetic fluctuations expressed in terms of the comoving wavenumber $`k`$ which equals the physical momentum of the wave at $`\eta =\eta _1`$ \[we set $`a(\eta _1)a_1=1`$\]. At later times the momentum is given by $`\omega (\eta )=k/a(\eta )`$ and the wavelength is $`2\pi a(\eta )/k`$.
For $`\eta <\eta _1`$ Eq. (2.5) is written in vacuum (no conductivity is present). As soon as the Universe gets into the radiation dominated phase the conductivity $`\sigma _c`$ suddenly jumps to a finite value of the order of $`T_1`$, namely the temperature right after $`\eta _1`$. Thus, the electric fields will be rapidly dissipated whereas the magnetic fields will not be damped, namely
$$_i(\eta _1+\mathrm{\Delta }\eta )_i(\eta _1),_i(\eta _1+\mathrm{\Delta }\eta )e^{\sigma _ca_1\mathrm{\Delta }\eta }_i(\eta _1).$$
(2.24)
Indeed, when the conductivity jumps to finite value, the magnetic fields survive whereas the electric fields are dissipated in a time $`\mathrm{\Delta }\eta `$ which is proportional to $`T_1^1`$ . This physical situation is different from the one occurring in a tokamak where the system is non-relativistic since the mass of the charge carriers is normally much larger than the temperature of the plasma.
In the absence of vorticity in the bulk velocity field of the plasma (i.e. if $`\stackrel{}{}\times \stackrel{}{v}=0`$), a given mode $`_i(\omega )=B_i(\omega )a^2(\eta )`$ of the magnetic field will approximately evolve as
$$B_i(\omega ,\eta )a^2(\eta )B_i(\omega ,\eta _1)a^2(\eta _1)e^{\frac{\omega ^2}{\sigma }\eta }.$$
(2.25)
where $`\sigma =\sigma _ca(\eta )`$ is the (curved-space) form of the conductivity. Provided that $`\omega <\omega _\sigma \sqrt{\sigma /\eta }`$ the magnetic flux is conserved and the magnetic energy density red-shifts as $`[a(\eta )]^4`$. Magnetic fields whose typical momentum is much smaller than $`\omega _\sigma `$ ( sometimes called magnetic diffusivity scale) correspond to the large scales relevant for the dynamo action.
If some primordial vorticity is present (i.e. $`\stackrel{}{}\times \stackrel{}{v}0`$) the evolution of the magnetic fields cannot be disentangled from the evolution of the bulk velocity field . The parity breaking (helical) velocity field could be the result of some dynamically developed turbulence , of some specific initial conditions , or of the gravitational collapse . Indeed, the galaxy is formed with some typical rotation period. The rotation of the galaxy will switch on a further term in the evolution equations of the magnetic fields: the so-called dynamo term describing the amplification caused in the magnetic flux by the action of a topologically non-trivial bulk velocity field (i.e. $`\stackrel{}{v}\stackrel{}{}\times \stackrel{}{v}0`$ ) . In the presence of the dynamo action the seed field possibly generated by the variation of the internal manifold will be exponentially amplified. The number of e-folds characterizing the amplification is of the order of the number of rotations performed by the galaxy since its origin. In order to allow a successful dynamo action (and a successful genesis of the large scale magnetic field) the original seed field (exponentially amplified later on) should have a specific value. The problem is to check if the seed fields generated from the evolution of the interval dimensions will be strong enough to be amplified to the presently observed value. This analysis will be one of the purposes of the following section.
## III Magnetogenesis constraints
The amplified gauge fields should be compared with various constraints coming from general cosmological considerations (critical energy density, possible anisotropies induced in the cosmic microwave background) and from magnetogenesis. If the produced fields are strong enough to seed the galactic dynamo mechanism they should also be compatible with the physics described by the standard model of cosmological evolution.
Suppose that the time evolution of the internal dimensions is monotonic in conformal time. In this case the internal dimensions will either grow or shrink. If the internal dimensions shrink (to Planckian or quasi-Planckian size) they are very small today (of the order of $`10^{33}`$$`10^{32}`$ cm). If they expand they could lead, today, to a smaller (effective) four-dimensional Planck mass and to a larger value of $`b_1`$ which could lie between $`10^{33}`$ cm and $`10^4`$ cm.
Some theoretical models which could describe the multidimensional evolution are discussed in Section IV. Here a general analysis will be given. The discussion is divided into two parts. In the first part $`n`$ internal dimensions growing in conformal time will be considered. This cooresponds, within the parametrization of Eq. (2.8), to the case $`\lambda <0`$. In the second part internal $`n`$ internal dimensions shrinking in conformal time will be considered. This corresponds to $`\lambda >0`$ in Eq. (2.8).
### A Internal dimensions growing in conformal time
According to Eq. (2.8) the growth of the internal dimensions stops in $`\eta =\eta _1`$ when the radiation dominated phase suddenly begins. The curvature at which the transition to the radiation dominated phase occur is given by $`H_1=T_1^2/[M_4(b_1M_4)^{n/2}]`$ where $`T_1`$ is the temperature at $`\eta _1`$,$`M_4`$ is the four dimensional Planck scale and $`H_1`$ is the Hubble factor $`H=\dot{a}/a`$ in $`\eta _1`$. The maximal temperature of the Universe cannot be larger, in this context, than $`M_4`$ so $`T_1M_4`$ will be assumed. The present (physical) frequencies relevant for the problem can be obtained using the background evolution during the radiation and matter dominated epochs. After $`\eta =\eta _1`$, the evolution of the geometry is adiabatic. Thus, the scale corresponding today to $`1/\eta _1`$ is given by $`\omega _110^{11}(M_4/M_P)`$ Hz (where $`M_P1.22\times 10^{19}`$ GeV). The decoupling frequency (corresponding to the transition from radiation to matter) is $`\omega _{\mathrm{dec}}10^{16}`$ Hz.
As discussed in the previous Section, thanks to magnetic flux conservation, the ratio between the logarithmic energy spectrum and the radiation energy density is approximately constant and, for $`T_1M_4`$, it is given by
$$r(\omega )=\left(\frac{M_4}{M_P}\right)^4\frac{𝒞(\nu )}{\pi ^2}\left(\frac{\omega }{\omega _1}\right)^{32\nu },$$
(2.1)
where, for simplicity, we set $`\lambda =\beta `$ with $`\beta >0`$ so that $`2\nu =|n\lambda 1|n\beta +1`$. Notice that the effective number of relativistic species at $`\eta _1`$ has been included in the definition of $`M_4`$. Eq. (2.1) can be rewritten as
$$r(\omega )=\frac{1}{\rho _\gamma }\frac{d\rho _B}{d\mathrm{ln}\omega }=\left(\frac{M_4}{M_P}\right)^{2+n\beta }\frac{2^{n\beta 4}}{\pi ^3}\left(\frac{\omega }{10^{11}\mathrm{Hz}}\right)^{2n\beta }\mathrm{\Gamma }^2\left(\frac{n\beta +1}{2}\right)|n\beta |^2,$$
(2.2)
where $`\rho _\gamma `$ is the energy density in radiation. If $`n\beta <2`$, the logarithmic energy spectrum grows in frequency. The critical energy density constraint is then implemented by requiring that $`r(\omega _1)<1`$, namely by requiring that the magnetic energy density (evolving according to flux conservation) is smaller than the radiation density. This requirement automatically guarantees, because of the growth in frequency of the logarithmic energy spectrum, that $`r(\omega )<1`$ also for all $`\omega <\omega _1`$, and in particular, for $`\omega _{\mathrm{dec}}`$.
If $`n\beta >2`$ the spectrum decreases in frequency. In order to insure the compatibility with present large scale observations at low frequencies the spectrum should not induce too much anisotropy in cosmic microwave background (CMB). This condition is enforced if $`r(\omega _{\mathrm{dec}})<10^{10}`$. If $`r(\omega _{\mathrm{dec}})<10^{10}`$ and $`r(\omega )`$ decreases in frequency, then, $`r(\omega )<1`$ for any $`\omega _{\mathrm{dec}}<\omega <\omega _1`$.
Thus, compatibility with the observed features of our present Universe demands
$`r(\omega _1)<1,\mathrm{for}n\beta <2,`$ (2.3)
$`r(\omega _{\mathrm{dec}})<10^{10},\mathrm{for}n\beta >2.`$ (2.4)
Eqs. (2.4) imply, respectively Recall that $`\mathrm{ln}`$ denotes the Neperian logarithm and $`\mathrm{log}`$ denotes the logarithm in ten basis.,
$`4\mathrm{log}\left({\displaystyle \frac{M_4}{M_P}}\right)<\mathrm{log}\left[{\displaystyle \frac{\pi ^3}{(n\beta )^2}}\right](n\beta 4)\mathrm{log}22\mathrm{log}[\mathrm{\Gamma }\left({\displaystyle \frac{n\beta +1}{2}}\right)],n\beta <2`$ (2.5)
$`(n\beta +2)\mathrm{log}\left({\displaystyle \frac{M_4}{M_P}}\right)<8.55+27(2n\beta )2\mathrm{log}n\beta `$ (2.6)
$`2\mathrm{log}\left[\mathrm{\Gamma }\left({\displaystyle \frac{n\beta +1}{2}}\right)\right]+(4n\beta )\mathrm{log}2,n\beta >2.`$ (2.7)
Magnetogenesis demands that strong seed fields should be produced. Such a demand translates into the following equations
$$r(\omega _G)3.7\times 10^{33},\omega _G10^{14}\mathrm{Hz}$$
(2.8)
where $`\omega _G^11`$ Mpc. The condition expressed in Eq. (2.8) mildly depends upon the cosmological parameters. Eq. (2.8) assumes $`h_0=0.65`$, $`\mathrm{\Omega }_\mathrm{m}=0.3`$ and $`\mathrm{\Omega }_\lambda =0.7`$ and it corresponds to a seed field of $`2.5\times 10^{16}`$ G at the decoupling epoch. The seed field generated through our mechanism will be exponentially amplified by the galactic rotation. The amplification induced by the differential rotation of the primeval galaxy through the dynamo mechanism goes roughly as $`e^{\mathrm{\Gamma }t}`$ where $`\mathrm{\Gamma }`$ is the dynamo amplification rate and $`t`$ is the galactic age (roughly $`10`$ Gyrs) . In Eq. (2.8) it is assumed that the dynamo amplification rate is of the order of $`\mathrm{\Gamma }^10.5`$ Gyr. If the dynamo amplification rate is (more or less artificially) increased the initial seed can be even smaller that the one assumed in Eq. (2.8). For instance if we take $`\mathrm{\Gamma }^10.3`$ the seed field is required to satisfy $`r(\omega _G)>6.07\times 10^{50}`$ . According to Eq. (2.8), Eq. (2.2) implies
$`(2+n\beta )\mathrm{log}\left({\displaystyle \frac{M_4}{M_P}}\right)`$ (2.9)
$`30.92+25(2n\beta )2\mathrm{log}\left[\mathrm{\Gamma }\left({\displaystyle \frac{n\beta +1}{2}}\right)\right]2\mathrm{log}n\beta (n\beta 4)\mathrm{log}2,`$ (2.10)
The constraints imposed by Eqs. (2.5)–(2.7) and by Eq. (2.10) entail a region in the parameter space of the model where the cosmological constraints can be safely satisfied and the magnetogenesis requirements are met. The parameter space (defined by $`n\beta `$ and by $`M_4/M_P`$) is illustrated in Fig. 1. The shaded area corresponds to the allowed region.
In the shaded area of Fig. 1 (left plot) magnetic fields are produced with growing frequency spectrum. In the shaded area of the right plot of Fig. 1 magnetic fields are produced with decreasing frequency spectrum. The theoretical bias would point, a priori, towards increasing frequency spectra. In fact in this case the two-point functions are decreasing at large distance scales. However, in the present analysis all the possibilities should be borne in mind. The full (thin) line appearing in both plots of Fig. 1 denotes a less conservative magnetogenesis requirement and it corresponds to larger a dynamo amplification rate. The thin curve can be obtained by requiring, from Eq. (2.2), $`r(\omega _G)>6.07\times 10^{50}`$ which corresponds to $`\mathrm{\Gamma }^10.5`$ Gyr. According to Fig. 1 reasonable seeds are produced provided $`n\beta `$ is larger than $`0.6`$ . The four dimensional Planck mass should be small in units of the fundamental Planck mass. The parameter space seems to select typical values of $`M_4`$ between $`10^{13}`$ Tev and $`10^3`$ TeV.
These conclusions were reached in the case $`T_1M_4`$. If $`T_1M_4`$ similar conclusions can be reached by following the same steps. The conclusion is that the intercepts between the thick and dot-dashed curves are slightly shifted to the left. To give some numerical value, assume, for instance, $`T_110^4M_4`$. Then we would get that the intercept gets to $`n\beta 0.8`$.
In order to avoid confusion we would like to stress that in spite of the fact that, in our case, $`M_4/M_P1`$ the physical size of the internal dimensions is always much smaller than $`0.1`$ mm corresponding to a $`M_4`$ TeV. Therefore, further constraints coming from submillimiter tests of the Newton law are not directly relevant to our present discussion.
To complete the analysis it is worth mentioning that magnetic fields can be directly constrained from BBN. These bounds are qualitatively different from the ones previously quoted and coming, alternatively, from homogeneity and isotropy of the background geometry at the BBN time. As elaborated in slightly different frameworks through the years , magnetic fields possibly present at the BBN epoch could have a twofold effect. On one hand they could enhance the rate of reactions (with an effect proportional to $`\alpha \rho _B`$ where $`\alpha 1/137`$) and, on the other hand they could artificially increase the expansion rate (with an effect proportional to $`\rho _B`$). It turns out that the latter effect is probably the most relevant. In order to prevent the Universe from expanding too fast at the BBN epoch $`\rho _B<0.27\rho _\nu `$ where $`\rho _\nu `$ is the energy density contributed by the standard three neutrinos for $`T<1`$ MeV. This bound can be imposed also to the spectra discussed in the present Section. Taking into account that the horizon at the BBN epoch (assuming a typical temperature of $`0.1`$ MeV) corresponds to a present frequency of $`\omega _{\mathrm{NS}}2.28\times 10^{12}`$ Hz, the BBN bound implies in the case of the spectra of Eq. (2.2)
$`(2+n\beta )\mathrm{log}\left({\displaystyle \frac{M_4}{M_P}}\right)<0.19+(4n\beta )\mathrm{log}2+22.6(2n\beta )`$ (2.11)
$`2\mathrm{log}\left[\mathrm{\Gamma }\left({\displaystyle \frac{n\beta +1}{2}}\right)\right]2\mathrm{log}n\beta `$ (2.12)
This constraint is qualitatively different from the ones already discussed and that is why it is important to mention it. However, from the quantitative point of view it is roughly equivalent to the critical density constraint. If $`n\beta >2`$ this bound is always satisfied. If $`n\beta <2`$ this bounds is a bit more constraining than the critical density limit (by a tiny amount) in the region $`1.4<n\beta <2`$. Within the numerical accuracy of the present analysis it is not crucial.
Another interesting moment in the life of the Universe where the generated magnetic fields could have an impact is the electroweak epoch . The present frequency corresponding at the electroweak horizon is of the order of $`10^5`$ Hz. Fields coherent on that scale can be certainly produced within the present mechanism but they are not constrained by the electroweak physics.
### B Internal dimensions shrinking in conformal time
If $`\lambda >0`$ (internal dimensions shrinking in conformal time), then $`2\nu =|n\lambda 1|`$. Therefore, the logarithmic energy spectrum has a twofold form
$`r(\omega )=\left({\displaystyle \frac{H_1}{M_P}}\right)^2\left({\displaystyle \frac{\omega }{\omega _1}}\right)^{4n\lambda }{\displaystyle \frac{2^{n\lambda 6}}{\pi ^3}}(2n\lambda )^2\mathrm{\Gamma }^2\left({\displaystyle \frac{n\lambda 1}{2}}\right),\mathrm{for}n\lambda >1`$ (2.13)
$`r(\omega )=\left({\displaystyle \frac{H_1}{M_P}}\right)^2\left({\displaystyle \frac{\omega }{\omega _1}}\right)^{2+n\lambda }{\displaystyle \frac{2^{2n\lambda }}{\pi ^3}}(n\lambda )^2\mathrm{\Gamma }^2\left({\displaystyle \frac{1n\lambda }{2}}\right),\mathrm{for}n\lambda <1`$ (2.14)
If $`n\lambda <4`$ the spectrum is always increasing in frequency. On the contrary for $`n\lambda >4`$ the spectrum decreases and it is peaked in the infra-red. In order to implement the critical density bound for $`n\lambda <4`$, $`r(\omega _1)<1`$ should be demanded. For $`n\lambda >4`$, $`r(\omega _{\mathrm{dec}})<10^{10}`$ should be required in order to guarantee the compatibility of our toy model with the CMB anisotropies. Then, using Eq. (2.14), the critical density bound demands
$`2\mathrm{log}\left({\displaystyle \frac{H_1}{M_P}}\right)(2+n\lambda )\mathrm{log}2+\mathrm{log}\pi ^32\mathrm{log}n\lambda 2\mathrm{log}\left[\mathrm{\Gamma }\left({\displaystyle \frac{1n\lambda }{2}}\right)\right],n\lambda <1.`$ (2.15)
$`2\mathrm{log}\left({\displaystyle \frac{H_1}{M_P}}\right)(6n\lambda )+\mathrm{log}\pi ^3`$ (2.16)
$`2\mathrm{log}(n\lambda 2)2\mathrm{log}\left[\mathrm{\Gamma }\left({\displaystyle \frac{n\lambda 1}{2}}\right)\right],\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\; 1}<n\lambda <4.`$ (2.17)
In this case the constraints are expressed using curvature (rather than temperature ) scales. This is because $`H_11/(a_1\eta _1)`$ cannot be much smaller than $`M_P`$ In this language the (present) maximal frequency turns out to be $`\omega _1=10^{11}\sqrt{H_1/M_P}`$ Hz. Using Eq. (2.14) in the case $`n\lambda >4`$ the CMB bound becomes
$`n\lambda \mathrm{log}\left({\displaystyle \frac{H_1}{M_P}}\right)`$ (2.18)
$`17.01+54(4n\lambda )2(n\lambda 6)\mathrm{log}22\mathrm{log}(2n\lambda )^24\mathrm{log}\left[\mathrm{\Gamma }\left({\displaystyle \frac{n\lambda 1}{2}}\right)\right]`$ (2.19)
Finally the magnetogenesis requirements \[already introduced in Eq. (2.8)\] impose
$`(2n\lambda )\mathrm{log}\left({\displaystyle \frac{H_1}{M_P}}\right)61.88+50(2+n\lambda )`$ (2.20)
$`+2(2+n\lambda )\mathrm{log}24\mathrm{log}n\lambda 4\mathrm{log}\left[\mathrm{\Gamma }\left({\displaystyle \frac{1n\lambda }{2}}\right)\right],n\lambda <1`$ (2.21)
$`n\lambda \mathrm{log}\left({\displaystyle \frac{H_1}{M_P}}\right)61.88+50(4n\lambda )`$ (2.22)
$`2(n\lambda 6)\mathrm{log}24\mathrm{log}(n\lambda 2)4\mathrm{log}\left[\mathrm{\Gamma }\left({\displaystyle \frac{n\lambda 1}{2}}\right)\right],n\lambda >1`$ (2.23)
Fig. 2 shows that, once the cosmological bounds are imposed, magnetogenesis is possible only of the parameter space implying that $`n\lambda >1`$. If $`n\lambda <1`$ magnetogenesis is excluded unless the produced magnetic fields over-close the Universe. This possibility must be rejected.
## IV Tailoring the model of internal evolution
In the previous Sections it has been established that large scale magnetic fields can be produced either if $`n\beta 0.6`$ (if the internal dimensions grow in conformal time) or if $`n\lambda >2.7`$ (if the internal dimensions shrink in conformal time) . Some toy model of internal evolution will now be analyzed. The purpose of this analysis is to show that the physical requirements coming from magnetogenesis and from other cosmological considerations can be indeed realized in specific models of dynamical evolution of a $`D`$-dimensional background geometry.
### A General considerations
Consider the $`D`$-dimensional action for a scalar field, minimally coupled to gravity:
$$S=S_g+S_m=d^Dx\sqrt{g}R+d^Dx\sqrt{g}\left[\frac{1}{2}g^{\alpha \beta }_\alpha \phi _\beta \phi V(\phi )\right],$$
(3.1)
where we work, for simplicity, in Planck units. We shall consider a homogeneous, Bianchi-type I metric background, whose spatial part is the product of two conformally flat manifolds as introduced in Eq. (2.2). For the background given in Eq. (2.2), the equations of motion obtained by varying the action with respect to $`g_{\mu \nu }`$ and $`\phi `$,
$$R_\mu ^\nu \frac{1}{2}\delta _\mu ^\nu R=\frac{1}{2}\left[_\mu \phi ^\nu \phi \frac{1}{2}\delta _\mu ^\nu g^{\alpha \beta }_\alpha \phi _\beta \phi +\delta _\mu ^\nu V(\phi )\right],g^{\alpha \beta }_\alpha _\beta \phi +\frac{V}{\phi }=0,$$
(3.2)
reduce simply to
$`d(d1)^2+n(n1)^2+2nd=\left({\displaystyle \frac{\phi ^2}{2}}+a^2V\right),`$ (3.3)
$`2(d1)^{}+(d1)(d2)^2+2n^{}+n(n+1)^2+2n(d2)`$ $`=`$ $`\left(a^2V{\displaystyle \frac{\phi ^2}{2}}\right),`$ (3.4)
$`2(n1)^{}+2d^{}+d(d1)^2+n(n1)^2+2(d1)(n1)`$ $`=`$ $`\left(a^2V{\displaystyle \frac{\phi ^2}{2}}\right),`$ (3.5)
$`\phi ^{\prime \prime }+\left[(d1)+n\right]\phi ^{}+{\displaystyle \frac{V}{\phi }}=0,`$ (3.6)
where $`=(\mathrm{ln}a)^{}`$, $`=(\mathrm{ln}b)^{}`$.
These equations are not all independent, and the scalar field equation, for instance, can be obtained from the other Einstein equations. By summing and subtracting the above equations one obtains
$`a^2V`$ $`=`$ $`\left[(d1)+n\right]^{}+\left[(d1)+n\right]^2,`$ (3.7)
$`{\displaystyle \frac{1}{2}}\phi ^2`$ $`=`$ $`\left[(d1)+n\right]^{}+(d1)^2n^2+2n,`$ (3.8)
$`^{}^{}`$ $`=`$ $`()\left[(d1)+n\right].`$ (3.9)
Consider now the case where the external and the internal scale factors can be parametrized as $`a(\eta )(\eta )^\sigma `$ and $`b(\eta )(\eta )^\lambda `$. If $`V=0`$ and $`\phi =0`$ we have that the well known solution of this sysmtem is indeed given by a power law behaviour of the scale factors provided the exponents satisfy, according to Eqs. (3.7)–(3.9)
$`(d1)\sigma +n\lambda =1,`$ (3.10)
$`(d1)\sigma ^2+n\lambda ^2=1+2\lambda .`$ (3.11)
These two conditions can be also easily expressed in cosmic time . In this case we have to recall that $`a(\eta )d\eta =dt`$. Therefore if we parameterize the solutions in cosmic time as
$$a(t)(t)^ϵ,b(t)(t)^\zeta $$
(3.12)
we will have that the exponents in cosmic and conformal time are simply related as $`ϵ=\sigma /(\sigma +1)`$ and $`\zeta =\lambda /(\sigma +1)`$. Using these two last relations into Eqs. (3.10)–(3.11) the two standard Kasner conditions are recovered, namely $`dϵ+n\zeta =1`$ and $`dϵ^2+n\zeta ^2=1`$. The solution of Eqs. (3.10)–(3.11) gives, for each dimension $`n`$ of the internal space, two conjugate solutions which we write in the case $`d=3`$:
$`\sigma _\pm ={\displaystyle \frac{1}{2}}\left[\sqrt{{\displaystyle \frac{3n}{n+2}}}1\right],`$ (3.13)
$`\lambda _\pm =\pm \sqrt{{\displaystyle \frac{3}{n(n+2)}}}.`$ (3.14)
The related cosmic time exponents can be easily obtained
$`ϵ_\pm ={\displaystyle \frac{1\sqrt{\frac{n+2}{3n}}}{13\sqrt{\frac{n+2}{3n}}}},`$ (3.15)
$`\zeta _\pm ={\displaystyle \frac{2}{n(13\sqrt{\frac{n+2}{3n}})}}.`$ (3.16)
In the parameterization of Eq. (3.12) the solution labeled by $`+`$ in the exponents of the internal manifold denote a growing solution whereas the solution denoted by $``$ denote a shrinking solution. In order to produce large scale magnetic fields we should require that, when the internal dimensions expand,
$$|n\lambda _{}|>0.6.$$
(3.17)
If we now insert the explicit expression for $`\lambda _{}`$ we have that Eq. (3.17) implies
$$\sqrt{\frac{3n}{n+2}}0.6$$
(3.18)
which is satisfied for $`n=1`$ and it always satisfied for any $`n>1`$. Since $`\lambda _{}`$ is negative for any $`n`$ $`b(\eta )(\eta )^\lambda _{}`$ will be growing for any $`n`$. It is useful to look also at the cosmic time picture. In cosmic time $`b(t)(t)^\zeta _{}`$. Since $`\dot{b}>0`$ for any $`n`$ the solution is expanding also in cosmic time Notice that the the over-dot denotes derivation with respect to cosmic time. The kinematical conditions on the expansions or contraction of a given background geometry should be stated in cosmic rather than in conformal time..
In order to produce large scale magnetic fields when the internal dimensions shrink according to the Kasner solutions we have to require
$$n\lambda _+>2.7.$$
(3.19)
If we insert the explicit expression of $`\lambda _+`$ we can see that Eq. (3.19) imposes the following inequality
$$\sqrt{\frac{3n}{n+2}}2.7$$
(3.20)
which can never be satisfied for any $`n`$. If the dynamics of the internal dimensions follows vacuum (Kasner) solutions, then, the requirement of generating large scale magnetic fields compatible with the cosmological bounds automatically selects an expanding dynamics of the internal dimensions.
Vacuum Kasner solutions can be generalized to the case of vanishing scalar potential. In this case the equations of motion can be always solved and the power-law ansatz can still be made for the scale factors. The solutions will be given by:
$`(d1)\sigma +n\lambda =1`$ (3.21)
$`\phi _{}^{}{}_{}{}^{2}=2\{1(d1)\sigma ^2n\lambda ^2+2\sigma \}{\displaystyle \frac{1}{\eta ^2}}`$ (3.22)
Also in this case we can easily show that if the solution is contracting (in the internal manifold) large scale magnetic fields cannot be produced. Indeed if $`b(\eta )`$ contracts we will have that the allowed range of $`n\lambda `$ is given by $`0<n\lambda <d`$. If we take $`d=3`$ we see that $`n\lambda <3`$. But according to our analysis $`n\lambda >2.7`$. Therefore only a tiny slice of parameter space (i.e. $`2.7<n\lambda <3`$ ) could give a positive result.
### B Explicit examples
Consider, as an example, the case $`n=6`$. In this case the internal manifold has six internal scale factors $`b(\eta )`$. As assumed in Eq. (2.2) both the internal and the external metrics will be maximally symmetric. Suppose that the internal dimensions are then expanding according to the Kasner solutions, namely $`b(\eta )(\eta )^\lambda _{}`$ and $`a(\eta )(\eta )^\sigma _{}`$. If $`n=6`$ $`\lambda _{}=1/4`$ and $`\sigma _{}=5/4`$. The the background model given in Eq. (2.8) will be in this case
$`a(\eta )=a_1\left({\displaystyle \frac{\eta }{\eta _1}}\right)^{5/4},b(\eta )=b_1\left({\displaystyle \frac{\eta }{\eta _1}}\right)^{1/4},\eta <\eta _1,`$ (3.23)
$`a(\eta )=a_1\left({\displaystyle \frac{\eta +2\eta _1}{\eta _1}}\right),b(\eta )=b_1,\eta _1\eta \eta _2,`$ (3.24)
where a sudden match to the radiation dominated epoch is assumed in $`\eta _1`$. The internal scale factors grow from $`\eta \mathrm{}`$ up to $`\eta _1`$. The external scale factor contracts. This is because of the mathematical nature of vacuum Kasner solutions. Using Eqs. (2.4) and (2.9)–(2.14) the evolution of the mode functions during the multidimensional phase will be given by
$$\varphi _k^{\prime \prime }\left[k^2\frac{21}{16\eta ^2}\right]\varphi _k=0,$$
(3.25)
whose solution is given by Eq. (2.15) with $`\nu =5/4`$. At the time $`\eta _1`$ the background will pass to radiation with a temperature $`T_1M_4`$. The produced large scale magnetic fields will evolve according to flux conservation and the ratio between the magnetic and radiation energy density is roughly constant. Then, from Eqs. (2.23)–(2.1), $`r(\omega )`$ can be explicitly written as
$$r(\omega )=\left(\frac{M_4}{M_P}\right)^4\frac{𝒞(5/4)}{\pi ^2}\left(\frac{\omega }{\omega _1}\right)^{1/2},\omega _{\mathrm{dec}}<\omega <\omega _1$$
(3.26)
where $`\omega _110^{11}(M_4/M_P)`$ Hz and $`\omega _{\mathrm{dec}}10^{16}`$ Hz are both evaluated at the present time and where $`𝒞(5/4)/\pi ^21.86\times 10^3`$ \[see Eq. (2.21)\]. The value of $`r(\omega )`$ at the scale relevant for magnetogenesis (i.e. $`\omega _G10^{14}`$ Hz) is given by
$$r(\omega _G)=1.86\times 10^{15.5}\left(\frac{M_4}{M_P}\right)^{3.5}.$$
(3.27)
As discussed in Eq. (2.8), in order to have a successful magnetogenesis $`r(\omega _G)>10^{33}`$ should be required. Imposing this requirement it turns out that from Eq. (3.27
$$\left(\frac{M_4}{M_P}\right)10^5.$$
(3.28)
If a less conservative magnetogenesis requirement is assumed (i.e. $`r(\omega _G)10^{50}`$) it can be checked that
$$\left(\frac{M_4}{M_P}\right)10^{10}$$
(3.29)
should be demanded. If $`M_410^{12}`$ TeV both Eqs. (3.28)–(3.29) are satisfied and $`\omega _110^7`$ Hz. Moreover, $`r(\omega _{\mathrm{dec}})10^{30}`$ with negligible impact on the CMB. If $`M_410^7`$ TeV large scale magnetic fields can be generated only if a higher dynamo amplification rate is assumed, so that Eq. (3.29) can be applied. In this case $`\omega _1100`$ Hz.
## V Discussion and Conclusions
In this paper the possible generation of large scale magnetic fields in the case of dynamical internal dimensions has been considered. The time dependence of internal dimensions (a generic ingredient of Kaluza-Klein models and of superstring theories) is severely constrained by our knowledge of BBN. The fact that internal dimensions could not be dynamical after the Universe was old of approximately one second does not exclude that they could have been dynamical for temperatures larger that $`1`$ MeV.
Since internal dimensions naturally break the conformal invariance of the evolution equations of Abelian gauge fields propagating in the four-dimensional world, electromagnetic fluctuations can be amplified. The typical amplitude and scale of the produced magnetic fields has been investigated under the assumption that the internal dimensions were dynamically evolving prior to the radiation dominated epoch.
In this context the problem of generating a seed field for the galactic dynamo action has been studied. The problem of magnetogenesis amounts to satisfying the numerical requirements necessary for a successful dynamo action together with all the bounds imposed by the standard model of cosmological evolution. If the conformal time evolution of the internal dimensions is monotonic and if the multidimensional phase is suddenly matched to radiation, then magnetogenesis is possible provided the dimensions belonging to the internal manifold are expanding. If the internal dimensions are contracting magnetogenesis is not possible.
The considerations reported in this paper should emerge as a part of a complete cosmological scenario. This problem has been only partially addressed in our considerations and further work is certainly needed. However, the idea of connecting large scale magnetic fields with the dynamics of internal dimensions seems to be of some interest.
## Acknowledgments
The author wishes to thank M. E. Shaposhnikov for useful discussions and valuable exchanges of ideas.
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# On the Origin of the Rotation Curves of Dark-Matter-Dominated Galaxies
## 1 Introduction
The rotation curves of dark-matter-dominated galaxies probe the mass profile of their galactic halos. The observed rotation curves of dwarf and low-surface-brightness (LSB) disk galaxies, therefore, offer a relatively direct test of theories of cosmological halo formation, free of the complicating effects of the dynamical coupling of dissipationless dark matter and dissipative baryonic matter which affect the mass profiles of baryon-dominated galaxies. Recently, attention has focused on the apparent conflict between the singular halo density profiles predicted by CDM N-body simulations and the observed rotation curves, which favor a flat-density core (cf. Moore et al. 1999). This has led to intense scrutiny of both the observations and the CDM model, primarily in one of three directions: improvement in the numerical resolving power of the CDM N-body simulations to determine better the logarithmic slope of the predicted density profiles at small radii (e.g. Moore et al. 1999), ideas for modifying the microscopic properties of the dark matter so as to retain the more successful aspects of the CDM model while flattening the halo density profiles at small radii (e.g. Davé et al. 2000, and references therein), and suggestions that the rotation curve data lack sufficient spatial resolution near the center to distinguish unambiguously between a density profile with a flat-density core and the singular profiles predicted by CDM N-body simulations (e.g. van den Bosch & Swaters 2000). Despite these uncertainties, a meaningful comparison between the data and predictions for the global properties of the rotation curves of dark-matter-dominated galaxies is possible which serves to test the CDM model and discriminate amongst different background cosmologies. To demonstrate this, we show how the mass and formation epoch of a halo can be extracted from its observed rotation curve, by application of an analytical model for halo formation which reproduces the empirical description of these rotation curves extremely well.
Burkert (1995) showed that the observed rotation curves of several dark-matter-dominated dwarf galaxies are consistent with a common density profile with a flat-density core, according to
$$\rho (r)=\frac{\rho _{0,\mathrm{B}}}{(r/r_{0,\mathrm{B}}+1)(r^2/r_{0,\mathrm{B}}^2+1)}.$$
(1)
Kravtsov et al. (1998) found that the rotation curves for a larger sample which included both dwarf and LSB galaxies are well fit by a similar universal halo profile given by
$$\rho (r)=\frac{\rho _S}{(r/r_S)^\gamma (1+(r/r_S)^\alpha )^{(\beta \gamma )/\alpha }},$$
(2)
with $`(\alpha ,\beta ,\gamma )=(2,3,0.2)`$, which is only slightly steeper than the Burkert profile in the core but shares the slope at large radii, where $`\rho r^3`$. The circular velocity profiles \[i.e. $`v(r)(GM[r]/r)^{1/2}`$\] for equations (1) and (2) are virtually indistinguishable at all radii, even near the center, where $`vr^{1\gamma /2}`$, since $`1\gamma /2=1`$ or $`0.9`$, respectively.
This contrasts with the halo profiles found by N-body simulations of the standard CDM model. The universal fitting formula for simulated halos reported by Navarro, Frenk, & White (1997; NFW) is equation (2) with $`(\alpha ,\beta ,\gamma )=(1,3,1)`$, while Moore et al. (1999) find $`(\alpha ,\beta ,\gamma )=(1.5,3,1.5)`$. However, the circular velocity profiles for these halos differ significantly from those implied by equations (1) and (2) only at very small radii where the uncertainties in the shape of the observed rotation curves currently makes discrimination difficult. In the meantime, equation (1) generally fits the data better than does the NFW profile even when the latter is also an acceptable fit.
The Burkert profile, then, continues to serve as a useful empirical description of the universal mass profiles of dark-matter-dominated galactic halos. We show here that the truncated, nonsingular isothermal sphere model we derived elsewhere (Shapiro, Iliev & Raga 1998, Paper I; Iliev & Shapiro 2000, Paper II) has a density profile for which the circular velocity profile is essentially indistinguishable from that of the Burkert profile and, as such, provides a theoretical motivation for the latter. The TIS model goes well beyond the prediction of density profile and rotation curve, however, to provide the size, mass, velocity dispersion and collapse epoch of the halo, as well. This makes possible the further interpretation of observations of dark-matter-dominated galaxies for comparison with the predictions of various cosmological models. At the same time, the TIS model can be shown to reproduce many of the average properties of the halos found in simulations of the standard CDM model, outside of the innermost region where the TIS halo has a flat-density core, unlike the CDM halos. As a result, a comparison of the analytical TIS predictions with observed properties of galaxies also offers insight into the standard CDM model. Finally, if suggestions like the self-interacting dark matter proposal of Spergel & Steinhardt (2000) are correct, that CDM might be more “collisional” as a way to eliminate the central cusp of standard CDM halos, then our TIS solution will also apply to these models, to the extent that the halo relaxation process makes the final equilibrium approximately isothermal.
In § 2, we briefly summarize the relevant properties of the TIS model and present a simple analytical formula for it with which to fit an observed rotation curve. In § 3 we demonstrate the excellent agreement between the TIS model rotation curve and that implied by equation (1) and show how this allows us to deduce the total mass and collapse epoch of a given halo directly from the parameters of its observed rotation curve. In § 4, we describe how the dependence of the average formation epoch of a halo on its mass in the CDM model results in statistical correlations amongst the parameters of the observed rotation curves. As an example, we use this approach to derive analytically the known correlation between the maximum velocity of each rotation curve and the radius at which it occurs.
## 2 The Truncated Isothermal Sphere Model
The TIS model is a particular solution of the Lane-Emden equation (suitably modified when $`\mathrm{\Lambda }0`$) which results from the collapse and virialization of a top-hat density perturbation (c.f. Paper I for Einstein-de Sitter \[EdS\] universe and Paper II for $`\mathrm{\Omega }_0<1`$ and $`\lambda _0=0`$ or $`1\mathrm{\Omega }_0`$). The size $`r_t`$ and velocity dispersion $`\sigma _V`$ are unique functions of the mass and redshift of formation of the object for a given background universe. While the Lane-Emden equation requires a straightforward numerical solution, Paper I provides a convenient analytical fitting formula,
$$\rho (r)=\rho _0\left[\frac{A}{a^2+\zeta ^2}\frac{B}{b^2+\zeta ^2}\right],$$
(3)
where $`\zeta r/r_0`$ and $`(A,a^2,B,b^2)=(21.38,9.08,19.81,14.62)`$, accurate to within 3% over the full range $`0\zeta \zeta _t30`$ for both the EdS case and most low-density models of interest (i.e. $`\mathrm{\Omega }_00.3`$). \[Note: Our definition of the core radius is $`r_{0,\mathrm{TIS}}r_{\mathrm{King}}/3`$, where $`r_{\mathrm{King}}`$ is the “King radius” defined in Binney & Tremaine (1987), p. 228.\] Equation (3) can be integrated to yield an analytical fitting formula for the TIS rotation curve, as well, given by
$$\frac{v(r)}{\sigma _V}=\left\{AB+\frac{1}{\zeta }\left[bB\mathrm{tan}^1\left(\frac{\zeta }{b}\right)aA\mathrm{tan}^1\left(\frac{\zeta }{a}\right)\right]\right\}^{1/2},$$
where $`\sigma _V=(4\pi G\rho _0r_0^2)^{1/2}`$. This fit has a fractional error less than 1% over the full range of radii $`0rr_t`$ for all matter-dominated background models. With a nonzero cosmological constant, the circular velocity becomes $`v(r)=(GM[r]/r)^{1/2}[12\rho _\lambda /\rho (r)]^{1/2}`$, where $`\rho _\lambda `$ is the constant vacuum energy density associated with the cosmological constant. For cases of current interest for a flat universe with $`\mathrm{\Lambda }0`$ (i.e. $`\mathrm{\Omega }_00.3`$), equation (2) can still be used, however, since it departs from the exact solution only slightly in the outer halo (i.e. at $`rr_t`$) for halos collapsing even as late as $`z=0`$; the fractional error is less than 6% at all radii and less than 1% for $`r(2/3)r_t`$.
The TIS model quantitatively reproduces the average structural properties of halos found in CDM simulations to good accuracy, suggesting that it is a useful analytical approximation for halos which form from more realistic initial conditions. An exception to this agreement is the very inner profile, where the TIS has a uniform-density core instead of a central cusp. Our TIS predictions agree to astonishingly high accuracy (i.e. to of order 1%) with the cluster mass-radius and radius-temperature relationships and integrated mass profiles derived from detailed CDM simulations of X-ray cluster formation by Evrard, Metzler, and Navarro (1996). Apparently, these simulation results are not sensitive to our disagreement in the core. A direct comparison of the TIS and NFW mass profiles reveals a very close agreement (fractional deviation of less than $`10\%`$) at all radii outside of a few TIS core radii (i.e. about one King radius), for NFW concentration parameters $`4c_{\mathrm{NFW}}7`$.
## 3 Application to Galaxy Rotation Curves
The rotation curve for equation (1) is given by
$$\frac{v_\mathrm{B}(r)}{v_{,\mathrm{B}}}=\left\{\frac{\mathrm{ln}\left[(\zeta _B+1)^2(\zeta _B^2+1)\right]2\mathrm{tan}^1(\zeta _B)}{\zeta _B}\right\}^{1/2},$$
(4)
where $`\zeta _Br/r_{0,\mathrm{B}}`$ and $`v_{,\mathrm{B}}(\pi G\rho _{0,\mathrm{B}}r_{0,\mathrm{B}}^2)^{1/2}`$. To compare our TIS directly with equation (1), we solve for the ratios $`\rho _{0,\mathrm{TIS}}/\rho _{0,\mathrm{B}}`$ and $`r_{0,\mathrm{TIS}}/r_{0,\mathrm{B}}`$ which minimize the $`\chi ^2`$ of the fit of equation (4) to equation (2) over radii from 0 to $`r_t`$. The result in Figure On the Origin of the Rotation Curves of Dark-Matter-Dominated Galaxies, with $`\rho _{0,\mathrm{TIS}}/\rho _{0,\mathrm{B}}=0.7790`$ and $`r_{0,\mathrm{TIS}}/r_{0,\mathrm{B}}=0.3276`$, shows extremely good agreement, with a fractional deviation below 10% from $`0.01r_t`$ to $`r_t`$ and below 4% over the range $`r>0.03r_tr_{0,\mathrm{TIS}}`$. For the TIS, the maximum circular velocity and its location are $`v_{\mathrm{max},\mathrm{TIS}}=1.5867\sigma _{\mathrm{V},\mathrm{TIS}}`$ at $`r_{\mathrm{max},\mathrm{TIS}}=8.99r_{0,\mathrm{TIS}}`$ for all matter-dominated cosmologies. For a flat universe with $`\mathrm{\Lambda }0`$, ($`\mathrm{\Omega }_00.3`$), these numbers depend weakly upon collapse redshift, but are reduced by no more than 0.2% and 1.7%, respectively, for $`z_{\mathrm{coll}}=0`$, and by even less for earlier collapse. For the Burkert profile, $`v_{\mathrm{max},\mathrm{B}}=1.2143v_{,\mathrm{B}}`$ at $`r_{\mathrm{max},\mathrm{B}}=3.2446r_{0,\mathrm{B}}`$. Hence, our best fit finds that $`v_{\mathrm{max}}`$ and $`r_{\mathrm{max}}`$ for the TIS and Burkert profiles are extremely close, with $`v_{\mathrm{max},\mathrm{TIS}}/v_{\mathrm{max},\mathrm{B}}=1.02`$ and $`r_{\mathrm{max},\mathrm{TIS}}/r_{\mathrm{max},\mathrm{B}}=0.91`$.
In short, our TIS model provides a solid, theoretical underpinning for the empirical fitting formula of Burkert, and, by extension, a self-consistent theoretical explanation of the observed galaxy rotation curves it was invented to fit. In addition, the TIS model allows us to calculate the total mass $`M_0`$, collapse epoch $`z_{\mathrm{coll}}`$, and other parameters of each observed dark-matter-dominated halo from its rotation curve, as follows. Assuming $`v_{\mathrm{max}}`$ and $`r_{\mathrm{max}}`$ are provided by observation<sup>1</sup><sup>1</sup>1 In practice, observations are generally restricted to the inner parts of galaxy rotation curves, often not extending to radii as large as $`r_{\mathrm{max}}`$. In that case, $`v_{\mathrm{max}}`$ and $`r_{\mathrm{max}}`$ are inferred by fitting the observed rotation curve to the TIS rotation curve at other radii., the mass of the galaxy can be shown to be
$$\frac{M_0}{h^1M_{}}=6.329\times 10^{10}\left(\frac{r_{\mathrm{max}}}{10h^1\mathrm{kpc}}\right)\left(\frac{v_{\mathrm{max}}}{100\mathrm{km}\mathrm{s}^1}\right)^2$$
(5)
for any matter-dominated universe. When $`\mathrm{\Lambda }0`$, $`M_0r_{\mathrm{max}}v_{\mathrm{max}}^2`$, but the coefficient depends on $`z_{\mathrm{coll}}`$ and background cosmology (Paper II). For any flat universe ($`\mathrm{\Lambda }0`$) of current interest, equation (5) is still a very good approximation. For $`\mathrm{\Omega }_0=1\lambda _00.3`$, equation (5) underestimates the mass by less than 6.3%, 2.8%, and 1.2% for $`z_{\mathrm{coll}}=0,0.5`$, and 1, respectively.
In terms of $`r_{\mathrm{max}}`$ and $`v_{\mathrm{max}}`$, $`z_{\mathrm{coll}}`$ is given to good accuracy by the implicit equation
$$F(\mathrm{\Omega }_0,\lambda _0,z_{\mathrm{coll}})=2.284\left(\frac{v_{\mathrm{max}}/100\mathrm{km}\mathrm{s}^1}{r_{\mathrm{max}}/10h^1\mathrm{kpc}}\right)^{2/3}$$
(6)
where $`F\{[\mathrm{\Omega }_0/\mathrm{\Omega }(z_{\mathrm{coll}})][\mathrm{\Delta }_{\mathrm{c},\mathrm{SUS}}/18\pi ^2]\}^{1/3}(1+z_{\mathrm{coll}})`$, $`\mathrm{\Omega }(z)=[\mathrm{\Omega }_0(1+z)^3]/[(1\mathrm{\Omega }_0\lambda _0)(1+z)^2+\mathrm{\Omega }_0(1+z)^3+\lambda _0]`$, and $`\mathrm{\Delta }_{\mathrm{c},\mathrm{SUS}}`$ is the density \[in units of $`\rho _{\mathrm{crit}}(z_{\mathrm{coll}})`$\] after top-hat collapse and virialization for the standard uniform sphere approximation (SUS)<sup>2</sup><sup>2</sup>2This $`\mathrm{\Delta }_{\mathrm{c},\mathrm{SUS}}`$ is well-approximated by $`\mathrm{\Delta }_{\mathrm{c},\mathrm{SUS}}=18\pi ^2+c_1xc_2x^2`$, where $`x\mathrm{\Omega }(z_{\mathrm{coll}})1`$, and $`c_1=82(60)`$ and $`c_2=39(32)`$ for the flat (open) cases, $`\mathrm{\Omega }_0+\lambda _0=1`$ ($`\mathrm{\Omega }_0<1,\lambda _0=0`$), respectively (Bryan & Norman, 1998).. For the EdS case, $`F=(1+z_{\mathrm{coll}})`$, while for open, matter-dominated and flat cases, $`F\mathrm{\Omega }_0^{1/3}(1+z_{\mathrm{coll}})`$ at early times \[i.e. $`x0`$\]. The coefficient on the r.h.s. of equation (6) is correct for any matter-dominated cosmology and a very good approximation for flat ($`\mathrm{\Lambda }0`$) cases of current interest. For $`\mathrm{\Omega }_00.3`$, it underestimates $`(1+z_{\mathrm{coll}})`$ by less than 2.5% for $`z_{\mathrm{coll}}0`$.
## 4 Statistical Correlations Amongst the Properties of Dwarf and LSB Galaxies: The $`v_{\mathrm{max}}`$$`r_{\mathrm{max}}`$ Relation
The halos in our TIS model are fully described for a given background universe by their total mass $`M_0`$ and collapse epoch $`z_{\mathrm{coll}}`$. In hierarchical models of structure formation like CDM, these are not statistically independent parameters, however. Smaller mass halos on average collapse earlier and are denser than larger mass halos. In terms of galaxy rotation-curve parameters, this dependence should be observable, for example, as a correlation between $`v_{\mathrm{max}}`$ and $`r_{\mathrm{max}}`$. Mori & Burkert (2000) used the Burkert (1995) fits to rotation curves of dwarf galaxies to report such an observed correlation, expressed as follows:
$$v_{\mathrm{max},\mathrm{B}}=9.81\left(r_{\mathrm{max},\mathrm{B}}/1\mathrm{kpc}\right)^{2/3}\mathrm{km}\mathrm{s}^1.$$
(7)
Kravtsov et al. (1998) also found a correlation using fits to observed rotation curves. Their results for a sample of dwarf and LSB galaxies are shown in Figure On the Origin of the Rotation Curves of Dark-Matter-Dominated Galaxies, along with that in equation (7). Kravtsov et al. (1998) further showed that the results of their CDM N-body simulations agreed with these data points.<sup>3</sup><sup>3</sup>3This result by Kravtsov et al. (1998) is not sensitive to the question of whether their simulations adequately resolved the halo density profiles at very small radii. This suggests that if we can relate the mass and collapse epoch of our TIS model halos in a statistical way within the context of the CDM model, a comparison of our predicted $`v_{\mathrm{max}}`$-$`r_{\mathrm{max}}`$ relation with this observed one will further check the relevance of our TIS model to halo formation from realistic initial conditions and, at the same time, give a theoretical explanation for both the data and the simulation results.
To predict the $`v_{\mathrm{max}}`$$`r_{\mathrm{max}}`$ correlation for a CDM universe using our TIS model, we apply the well-known Press-Schechter (PS) approximation to derive $`z_{\mathrm{coll}}(M_0)`$ \- the typical collapse epoch for a halo with a given mass. Halos of mass $`M`$ which collapse when $`\sigma (M)=\delta _{\mathrm{crit}}/\nu `$ are referred to as “$`\nu `$-$`\sigma `$” fluctuations, where $`\sigma (M)`$ is the standard deviation of the density fluctuations at $`z_{\mathrm{coll}}`$ according to linear perturbation theory, after the density field is filtered on the scale $`M`$, and $`\delta _{\mathrm{crit}}`$ is the amplitude of a top-hat perturbation according to linear theory at the epoch $`z_{\mathrm{coll}}`$ at which the exact solution predicts infinite collapse. The typical collapse epoch for halos of a given mass is that for which $`\nu =1`$, the 1-$`\sigma `$ fluctuations.
Our results for 1-$`\sigma `$ fluctuations are shown in Figure On the Origin of the Rotation Curves of Dark-Matter-Dominated Galaxies (upper panel), for different background cosmologies. The flat, untilted and the open, slightly-tilted ($`n_p=1.14`$) models are in reasonable agreement with the observed $`v_{\mathrm{max}}r_{\mathrm{max}}`$ relation, while the untilted and strongly tilted ($`n_p=1.3`$) open models are not. The empirical Burkert scaling relation is closely approximated by the currently-favored $`\mathrm{\Lambda }`$CDM model, less well by other models.
In practice, observed galaxies should exhibit a statistical spread of halo properties in accord with the expectations of the Gaussian statistics of the density fluctuations which formed them. Since observed galaxies will not all be “typical”, the scatter of the data points in Figure On the Origin of the Rotation Curves of Dark-Matter-Dominated Galaxies is natural. To probe this in our model, we calculate the masses and collapse redshifts for halos formed by $`\nu `$-$`\sigma `$ fluctuations for different values of $`\nu `$. As shown in Figure On the Origin of the Rotation Curves of Dark-Matter-Dominated Galaxies (lower panel) for $`\mathrm{\Lambda }`$CDM, the Burkert scaling relation is closest to the TIS model prediction for 1-$`\sigma `$ fluctuations, while all the observed galaxy data points except one correspond to $`0.7\nu 1.5`$. Hence, on average, the galaxies which constitute the observed $`v_{\mathrm{max}}`$-$`r_{\mathrm{max}}`$ correlation correspond to halos which formed at close to the typical collapse time expected theoretically for objects of that mass. This means that our TIS model is a self-consistent explanation for the observed $`v_{\mathrm{max}}`$-$`r_{\mathrm{max}}`$ correlation.
This success of the TIS model in explaining the observed $`v_{\mathrm{max}}`$-$`r_{\mathrm{max}}`$ relation and, by extension, the CDM simulation results of Kravtsov et al. (1998) which follow it can be understood by a completely analytical argument, as follows. We approximate the density fluctuation power spectrum as a power-law in wavenumber $`k`$, $`P(k)k^n`$. If we define a mass $`M`$ which corresponds to $`k`$ according to $`Mk^3`$, then $`\sigma (M)M^{(3+n)/6}`$ if we set $`n=n_{\mathrm{eff}}3(2y_{\mathrm{eff}}+1)`$, where $`y_{\mathrm{eff}}(d\mathrm{ln}\sigma /d\mathrm{ln}M)_{\mathrm{exact}}`$ at the relevant mass scale $`M`$. Our results for the $`\mathrm{\Lambda }`$CDM case indicate that the galaxies which make up the $`v_{\mathrm{max}}r_{\mathrm{max}}`$ data points in Figure On the Origin of the Rotation Curves of Dark-Matter-Dominated Galaxies collapsed at redshifts $`1z_{\mathrm{coll}}6`$ with masses in the range $`8\times 10^9M_0/(M_{}h^1)3\times 10^{11}`$. Hence, the precollapse fluctuation growth rate is approximately EdS, and we can let $`\mathrm{\Omega }(z_{\mathrm{coll}})=1`$ in equations (5) and (6). In that case, $`(1+z_{\mathrm{coll}})\sigma (M)M^{(3+n)/6}`$, $`r_{\mathrm{max}}M^{(5+n)/6}\mathrm{\Omega }_0^{1/3}`$ and $`v_{\mathrm{max}}M^{(1n)/12}\mathrm{\Omega }_0^{1/6}`$, which combine to yield
$$v_{\mathrm{max}}=v_{\mathrm{max},}\left(r_{\mathrm{max}}/r_{\mathrm{max},}\right)^{(1n)/[2(5+n)]},$$
(8)
where $`v_{\mathrm{max},}`$ and $`r_{\mathrm{max},}`$ are for a 1-$`\sigma `$ fluctuation of fiducial mass $`M_{}`$, as given by equations (5) and (6) with $`(1+z_{\mathrm{coll}})=(1+z_{\mathrm{rec}})\sigma (M_{},z_{\mathrm{rec}})/\delta _{\mathrm{crit}}`$, where $`\delta _{\mathrm{crit}}=1.6865`$ and $`\sigma (M_{},z_{\mathrm{rec}})`$ is the value of $`\sigma (M_{})`$ evaluated at the epoch of recombination $`z_{\mathrm{rec}}`$ \[i.e. early enough that $`(1+z)\sigma `$ is independent of $`z`$\]. Over the relevant mass range $`M=10^{10\pm 1}h^1M_{}`$, $`n_{\mathrm{eff}}2.4\pm 0.1`$ for our $`\mathrm{\Lambda }`$CDM case. For $`M_{}=10^{10}h^1M_{}`$, our COBE-normalized, flat $`\mathrm{\Lambda }`$CDM case ($`\mathrm{\Omega }_0=0.3,h=0.65`$) yields $`(1+z_{\mathrm{rec}})\sigma (M_{},z_{\mathrm{rec}})=5.563`$, so $`(1+z_{\mathrm{coll}})=3.30`$, $`v_{\mathrm{max},}=53.2\mathrm{km}\mathrm{s}^1`$ and $`r_{\mathrm{max},}=5.59h^1\mathrm{kpc}`$. With these values and $`n_{\mathrm{eff}}=2.4`$, equation (8) yields the TIS model analytical prediction $`v_{\mathrm{max}}=(13.0\mathrm{km}\mathrm{s}^1)(r_{\mathrm{max}}/1\mathrm{kpc})^{0.65}`$, remarkably close to the empirical relation in equation (7).
We thank Hugo Martel for valuable discussion, grants NASA ATP NAG5-7363 and NAG5-7821, NSF ASC-9504046, and Texas Advanced Research Program 3658-0624-1999, and an NSF International Research Fellow Award INT-0003682 to ITI.
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# 1 Introduction
## 1 Introduction
Numerous attempts have been made to explain the origin of CP violation in the context of string-derived effective supergravity models. Early on, Strominger and Witten suggested compactification might lead to explicit violation of CP from the four-dimensional point of view, since the operation of CP is orientation-changing in the six-dimensional compact space . Inspired by this prospect, detailed analyses were carried out for various compactifications of the heterotic string; for example, a Calabi-Yau manifold was examined in and $`Z_N`$ orbifolds were investigated in . In each case where explicit CP violation was investigated, it was found not to occur. Subsequently, Dine et al. argued that CP is a gauge symmetry in string theory and that explicit breaking is therefore forbidden both perturbatively and nonperturbatively . In the same work it was shown that this is certainly true for heterotic orbifolds. Based on these results, it is clear that if a heterotic orbifold is to provide a reasonable approximation to the correct underlying theory of known interactions, CP violation must occur spontaneously from string moduli or matter fields (or perhaps both) acquiring vacuum expectation values (vevs).
In a series of papers, Bailin, Kraniotis and Love (BKL) considered the possibility of supersymmetric CP violation (SCPV)<sup>1</sup><sup>1</sup>1By this, we mean CP-violating complex phases in the soft supersymmetry breaking operators of minimal extensions to the Standard Model (SM) . in heterotic orbifold models . They related SCPV to the complex phases of string moduli vevs. While SCPV is an interesting possibility, it is important to understand how the Kobayashi-Masakawa (KM) phase might occur in semi-realistic orbifold models and how generic it is. Furthermore, SCPV in the context of supergravity is phenomenologically problematic unless soft terms meet a variety of stringent criteria ; the current status of SCPV and various solutions to related phenomenological problems are reviewed in ref. . By contrast, the KM phase does not have dangerous side-effects; the smallness of CP violation is for the most part explained by the small mixing angles between heavy and light generation quarks while large electric dipole moments do not arise from this source .
BKL have considered complex vevs for the string moduli of heterotic orbifolds as a source of the KM phase . T- and U-moduli, which parameterize deformations of the six-dimensional compact space consistent with the orbifold construction, enter into the effective Yukawa couplings for twisted fields, as described below. In this article we will restrict ourselves to $`Z_3`$ and $`Z_3\times Z_3`$ orbifolds, which have no U-moduli because the complex structure is fixed. It has been demonstrated by BKL and others that is possible for nonperturbative effects to stabilize these string moduli at complex values . BKL have shown that, as a result of complex string moduli vevs, $`𝒪(1)`$ phases can arise in twisted Yukawa coupling coefficients. However, they do not construct any models and it is unclear whether or not the phases they find are physical: phases in the quark Yukawa matrices do not necessarily imply the existence of a nonzero KM phase.<sup>2</sup><sup>2</sup>2Complex phases which give rise to a nontrivial KM phase will be termed physical while those which can be eliminated by rephasing quark fields will be termed spurious. As an example, let us write the quark mass superpotential as
$$W_{qm}=\lambda _{ij}^uH_uQ_{iL}u_{jL}^c+\lambda _{ij}^dH_dQ_{iL}d_{jL}^c$$
(1.1)
and suppose Yukawa matrices of the form
$$\lambda _{ij}^u=|\lambda _{ij}^u|e^{i(\alpha _i+\beta _j^u)},\lambda _{ij}^d=|\lambda _{ij}^d|e^{i(\alpha _i+\beta _j^d)}.$$
(1.2)
The phases $`\alpha _i,\beta _j^{u,d}`$ are not physical since they can be removed by rephasing the quark fields according to
$$Q_{iL}e^{i\alpha _i}Q_{iL},u_{jL}^ce^{i\beta _j^u}u_{jL}^c,d_{jL}^ce^{i\beta _j^d}d_{jL}^c.$$
(1.3)
Naive orbifold models of quark Yukawa couplings assign the nine SM multiplets $`Q_{iL},u_{jL}^c,d_{jL}^c`$ to different sectors of the Hilbert space and rely on trilinear couplings to give all of the quarks mass; the assignments to different sectors determine the moduli-dependence of the effective Yukawa matrices, hence the complex phases. If these assignments can be brought into correspondence with phases that enter the Yukawa matrices in a way similar to (1.2), then they will not give rise to CP violation. An example of this is described in Section 2 below. Furthermore, when one reviews the tables of phases displayed in the appendix of ref. , one finds that many of them are identical or zero for a given orbifold model and T-modulus vev. It then becomes a concern whether or not this high degree of degeneracy causes the phases to “wash out” in the final analysis. An example of this is also given in Section 2.
Beyond these more obvious ways that complex phases in the Yukawa matrices may not give rise to a nontrivial KM phase, we must also be concerned that symmetry constraints imposed on the Yukawa matrices by the underlying string theory might relate the phases to each other in such a way that the CKM matrix can be made entirely real. The embedding of the orbifold action into the internal left-moving gauge degrees of freedom typically leaves a surviving gauge group significantly larger than the SM. For example, in the three generation constructions which we will discuss below the surviving gauge group $`G`$ has rank sixteen . Low-energy effective Yukawa couplings must be constructed from high-energy operators invariant under $`G`$. The high-energy operators are also subject to orbifold selection rules, which result from symmetries of the six-dimensional compact space. (A brief review of orbifold selection rules may be found in ref. .) Finally, the underlying conformal field theory has target space modular symmetries associated with the identification of equivalent string moduli backgrounds. It is conceivable (albeit highly unlikely) that orbifold selection rules, gauge invariance under $`G`$ and target-space modular invariance may conspire to make phases derived from the scalar background spurious. In the present paper we construct a toy $`Z_3`$ orbifold model which is subject to these symmetry constraints; we construct explicit quark mass matrices in order to make conclusive statements about the existence of a nonzero KM phase.
The case of complex vevs for the T-moduli is treated below; however, we place more emphasis on another origin of complex phases, which in our opinion is a much more generic and natural source of CP violation in semi-realistic heterotic orbifold models. Nonrenormalizable couplings are often important in semi-realistic heterotic orbifold models for the following reason. Wilson lines are typically included in the embedding to get a reasonable gauge group.<sup>3</sup><sup>3</sup>3 Standard GUT scenarios require large higgs representations, whereas they are absent in affine level 1 orbifold constructions . This makes it phenomenologically advantageous to obtain $`G=G_{SM}\times G_{other}`$ from the start. In many cases their inclusion leads to an anomalous $`U(1)_X`$ factor: $`\text{tr}Q_X0`$. The apparent anomaly is cancelled by the Green-Schwarz mechanism, which induces a Fayet-Illiopoulos (FI) term . Several scalar fields get vevs $`v_i𝒪(10^{2\pm 1})`$ (in units where $`m_P=1/\sqrt{8\pi G}=1`$) to restabilize the vacuum in the presence of the FI term; thus, the suppression of nonrenormalizable couplings due to vevs $`v_i`$ may be as little as $`𝒪(10^1)`$. The scalar fields which cancel the FI term break the $`U(1)_X`$ gauge symmetry at the scale $`\mathrm{\Lambda }_X10^{17}`$ GeV by the Higgs mechanism; in order to distinguish them from the higgses associated with electroweak symmetry breaking, we will for convenience and with all due apologies refer to them as Xiggses. The Xiggses are usually charged under other $`U(1)`$ factors of $`G`$ besides $`U(1)_X`$; a basis of generators can always be chosen such that these other $`U(1)`$’s are not anomalous.<sup>4</sup><sup>4</sup>4We will not consider the case where Xiggses are in nontrivial representations of the nonabelian factors of $`G`$. Vacuum stabilization near the scale $`\mathrm{\Lambda }_X`$ requires the D-terms of these other $`U(1)`$’s to vanish. To satisfy the numerous D-flatness conditions, it is generally necessary for several Xiggses to get vevs. Typically there are $`𝒪(10)`$ or more such Xiggses.
In this paper we demonstrate how the KM phase in semi-realistic heterotic orbifold models arises generically from the “Planck slop” created by the Xiggses. Nonrenormalizable couplings make significant and in some cases leading order contributions to the effective quark mass matrices. For instance, in the semi-realistic model developed by Font, Ibáñez, Quevedo and Sierra in Section 4.2 of ref. (FIQS model), the down-type quarks get their leading order mass from dimension 10 holomorphic couplings. (We count dimensions as 1 for each elementary superfield entering a coupling.) For up-type quarks, only the top and charm get masses from renormalizable couplings.<sup>5</sup><sup>5</sup>5We identify top, charm, etc., by ranking the mass at a given level of analysis; of course, the identification will be imperfect once higher-order corrections are included. The up quark must receive its mass from nonrenormalizable couplings or radiative mass terms. Although high order nonrenormalizable couplings are suppressed by large powers of the $`𝒪(10^{2\pm 1})`$ Xiggs vevs, the suppression is not as large as one might think, for two reasons. Many high order operators involving the Xiggses generically exist, with their number increasing at each higher order. Most of the operators are closely related by variations in fixed point locations and oscillator “directions”, as will be shown in detail below; however, this does not change the fact that the number of distinct operators tends to be large. The second reason why high order couplings are important is that the coupling strength tends to be much larger than $`𝒪(1)`$ and to grow as the dimension of the coupling increases, as was pointed out by Cvetič et al. . The combination of these two effects forces one to proceed to rather high order before couplings make negligible contributions to the quark mass matrices. Here, “negligible” is taken to mean less than, say, 10% contributions to the lightest quark masses. As a consequence, each effective Yukawa coupling $`\lambda _{ij}^{u,d}`$ depends on the vev of a linear combination of a large number of monomials of Xiggses. The principal point of this article is that since $`𝒪(10)`$ Xiggses get complex vevs, which appear in a large number of monomials contributing to effective quark Yukawa couplings, the Yukawa matrices are generically complex and a nonvanishing KM phase is almost inevitable. Indeed, in any orbifold model with an anomalous $`U(1)_X`$ present, it seems improbable that one would not have CP violation in this way, since it is difficult to see how nonrenormalizable couplings involving the Xiggses would not contribute to the effective quark Yukawa couplings.
Previous authors have noted the possible role of nonrenormalizable couplings in heterotic orbifold and free fermionic models for giving large mass hierachies and CP violation from a KM phase . In this respect the mechanism analyzed here is not new. However, we present much more detailed results by constructing an explicit toy model and we impose modular invariance on the nonrenormalizable couplings. The constraint of modular invariance leads to relationships between coupling coefficients which were not accounted for in earlier efforts. The model is inspired by three generation heterotic $`Z_3`$ orbifold models previously investigated in the literature . Our model is, by construction, quite similar to the FIQS model, which is string-derived: in the FIQS model, the gauge group, the spectrum of states and the allowed superpotential couplings are completely determined from the underlying string theory. The FIQS model suffers from phenomenolgical difficulties related to the quark mass matrices; these difficulties were previously pointed out in . At leading order in the FIQS model, the top and charm come from different $`SU(2)`$ doublets than the bottom. As a result, the leading order CKM matrix has some of its diagonal entries zero and some $`𝒪(1)`$ heavy-light generation mixing angles, which is clearly unacceptable. The assignments into $`SU(2)`$ doublets are determined by the H-momenta<sup>6</sup><sup>6</sup>6 H-momenta are the $`SO(10)`$ weights of bosonized NSR fermions, which appear in the vertex operators creating asymptotic states. of the untwisted states, so the identification of which doublet a given $`u_{iL}`$ or $`d_{iL}`$ sits in and how it couples at leading order is fixed. A second problem with the FIQS model is that it is difficult to give the three lightest quarks mass. We have searched for dimension $`d30`$ allowed holomorphic couplings which might do the job (under the assumptions made in the FIQS model about Xiggs and T-moduli vevs) and found that none exist. We further found that radiative masses only occur at high loop levels and would be minuscule in comparison to the experimental values. However, we assumed the leading order Kähler potential in this analysis. It is possible that higher order terms in the Kähler potential could provide a mechanism for giving the light quarks masses in agreement with experimental values.
We have attempted in many ways to evade these problems. However, we were not able to do so without creating other difficulties. We are actively searching for a superior three generation heterotic $`Z_3`$ orbifold model to study. In the meantime, we have developed a toy model which replicates the FIQS model wherever possible while avoiding its problems, in order that we might illustrate that a nontrivial KM phase generically arises from the complex vevs of Xiggses, even after orbifold selection rules and target space modular invariance have been accounted for.
Although the coupling coefficients for nonrenormalizable superpotential couplings (which are in principle obtainable from the underlying conformal field theory) are not apparently known, we propose couplings which transform in the requisite manner under the $`[SL(2,𝐙)]^3`$ diagonal subgroup of the full $`SU(3,3,𝐙)`$ modular duality group of the $`Z_3`$ orbifold . We take into account the non-trivial transformations of twisted sector fields in our construction of modular invariant couplings. In the three generation $`Z_3`$ orbifold models upon which the toy model of Section 4 is based, different species are distiguished by quantum numbers (other than the fixed point location in the third complex plane) of massless states in the underlying theory . For twisted fields $`\mathrm{\Phi }_n^i`$, where $`n`$ labels the species and $`i=1,2,3`$ labels the fixed point locations in the third complex plane, $`\mathrm{\Phi }_n^1,\mathrm{\Phi }_n^2`$ and $`\mathrm{\Phi }_n^3`$ mix amongst themselves under the $`T^31/T^3`$ duality transformation in the third complex plane . Constraining holomorphic polynomials of dimension $`d>3`$ to transform with<sup>7</sup><sup>7</sup>7Modular weight will be explained below. modular weight $`1`$ in light of these nontrivial mixings places strong constraints on the form of the superpotential terms and gives some confidence that our proposed nonrenormalizable coupling coefficients may reproduce key features of the actual couplings which would be derived using conformal field theory techniques. In the course of discussing our assumptions for the coupling coefficients of nonrenormalizable superpotential terms, we will explain why the calculation of these coefficients from the underlying string theory represents an extremely difficult problem. For now, we remark that the most intimidating aspect of such a calculation is the integration of the string correlator over the $`d3`$ vertex locations which cannot be fixed by $`SL(2,𝐂)`$ invariance, where $`d`$ is the dimension of the coupling. The string correlator is generally a very complicated function of the unfixed vertex locations.
The reliability of the effective supergravity approach is hampered by theoretical uncertainties in the Kähler potential of $`Z_3`$ orbifold models. Nonleading operators in the Kähler potential are neglected in most analyses; however, with $`𝒪(10^{2\pm 1})`$ Xiggs vevs, higher order terms in the Kähler potential give non-negligible corrections to the mass matrices of quarks: corrections from higher order terms give the quarks noncanonical, nondiagonal kinetic terms which must be rendered canonical by nonunitary field redefinitions when one goes to compute mass eigenstates and mixings. These complications cannot be ignored if one hopes to develop an accurate picture of the low-energy phenomenology predicted by a given model. Higher order terms in the Kähler potential ought to be included in order to be consistent with the high order expansion of the superpotential, both of which are necessary in order to pick up all significant contributions to the quark mass matrices. The calculation of higher order corrections to the Kähler potential is notoriously difficult because of the lack of holomorphicity. It is hoped that future work on the Kähler potential of heterotic $`Z_3`$ orbifold models will amend these deficiencies and allow for an improved analysis of the low energy phenomenology of semi-realistic models. Present ignorance regarding these aspects of the Kähler potential has forced us to make a number of oversimplifications. However, we do not expect these oversimplifications to affect our main result that a nontrivial KM phase is generic. Introducing higher order terms in the Kähler potential is not expected to eliminate the complex phases which enter into the effective quark Yukawa couplings from the vevs of Xiggses.
In Section 2 we present examples where the complex phases found by BKL are spurious. In Section 3 we introduce modular invariant coupling coefficients for nonrenormalizable superpotential couplings. In Section 4 we discuss our string-inspired toy heterotic $`Z_3`$ orbifold model. In Section 5 we make concluding remarks and suggest further investigations motivated by our results. In the Appendix we address normalization conventions for $`U(1)`$ charges in string-derived models. We explain how to account for different conventions when determining the FI term and describe the Green-Schwarz cancellation of the $`U(1)_X`$ anomaly in the linear multiplet formulation.
## 2 Counterexample
Here, we consider some assignments of quarks and higgses in $`Z_3\times Z_3`$ orbifold models and show that the complex phases found by BKL do not lead to a nontrivial KM phase for these particular examples. We certainly do not wish to imply that the phases found by these authors cannot lead to a nonzero KM phase; we only wish to point out that due to the degeneracy in phases (in this case always $`0`$ or $`\pi /3`$), they can in many cases be eliminated by rephasing quark fields. In our opinion, a more careful analysis is required in order to conclude whether or not the phases found by BKL can account for CP violation.
We will use the notation and conventions of ref. in our discussion of the twisted sectors and fixed tori of the $`Z_3\times Z_3`$ orbifold. This orbifold is constructed using twists
$$\theta =\frac{1}{3}(1,0,1),\omega =\frac{1}{3}(0,1,1).$$
(2.1)
We make use of the $`\theta `$, $`\theta \omega `$ and $`\theta \omega ^2`$ twisted sectors. The fixed tori for each of these sectors are given by
$`f_\theta `$ $`=`$ $`{\displaystyle \frac{m_1}{3}}(2e_1+\stackrel{~}{e}_1)+{\displaystyle \frac{m_3}{3}}(e_3\stackrel{~}{e}_3)+v_2,m_{1,3}=0,\pm 1,v_2K_2,`$ (2.2)
$`f_{\theta \omega }`$ $`=`$ $`{\displaystyle \frac{1}{3}}{\displaystyle \underset{i=1}{\overset{3}{}}}r_i(2e_i+\stackrel{~}{e}_i)+\mathrm{},r_i=0,\pm 1,\mathrm{}\mathrm{\Lambda },`$ (2.3)
$`f_{\theta \omega ^2}`$ $`=`$ $`{\displaystyle \frac{p_1}{3}}(2e_1+\stackrel{~}{e}_1)+{\displaystyle \frac{p_2}{3}}(e_2\stackrel{~}{e}_2)+v_3,p_{1,2}=0,\pm 1,v_3K_3,`$ (2.4)
where $`\mathrm{\Lambda }`$ is the $`[SU(3)]^3`$ root lattice and $`K_i`$ is the $`i`$th complex plane. Physical states must be simultaneous eigenstates of $`\theta `$ and $`\omega `$; they are therefore linear combinations of states whose zero modes are given by different fixed tori. Since the first complex plane is neutral under $`\omega `$, physical states in the $`\theta `$ sector can be chosen with a definite quantum number $`m_1`$. Because the form of the fixed torus (2.4) in the first complex plane is the same as in the $`\theta `$ sector (2.2), and since $`\omega `$ does not rotate in the first plane, a physical state in the $`\theta \omega ^2`$ sector can be chosen to have a definite quantum number $`p_1`$ as well. The coefficients of the trilinear Yukawa couplings are determined by the evaluation of correlation functions in the underlying string theory. The contribution from a complex $`T^1`$ in the classical partition function is the only source of complex phases in the trilinear couplings. The phases of $`T^{2,3}`$ do not matter because the twist operator contributions to the correlation function in the second and third complex planes reduce to the identity, as discussed in ref. . The phases found in ref. for the case $`T^1=\mathrm{exp}(i\pi /6)`$ are determined by the difference $`m_1p_1`$:
$$\gamma (m_1,p_1)\mathrm{arg}\underset{X_{cl}}{}e^{S_{cl}}=\{\begin{array}{cc}0,\hfill & m_1p_1=0;\hfill \\ \frac{\pi }{3},\hfill & m_1p_1=\pm 1,\pm 2.\hfill \end{array}$$
(2.5)
We next suppose assignments as follows: $`H_u,H_d`$ are in the $`\theta `$ sector with fixed tori quantum numbers $`m_1^u,m_1^d`$ resp.; $`Q_{iL}`$ are in the $`\theta \omega ^2`$ sector with fixed tori quantum numbers $`p_1^i`$ resp.; $`u_{jL}^c,d_{jL}^c`$ are in the $`\theta \omega `$ sector. The phases (2.5) enter into the Yukawa couplings (1.1) according to:
$$\lambda _{ij}^u=|\lambda _{ij}^u|e^{i\alpha _i^u},\lambda _{ij}^d=|\lambda _{ij}^d|e^{i\alpha _i^d},$$
(2.6)
where $`\alpha _i^u=\gamma (m_1^u,p_1^i)`$ and $`\alpha _i^d=\gamma (m_1^d,p_1^i)`$. If $`m_1^u=m_1^d`$ then we recover the CP conserving forms of (1.2) with $`\alpha _i=\alpha _i^u=\alpha _i^d`$ and $`\beta _j^{u,d}=0`$.
It is an interesting question whether or not CP violation really can occur, and the example just presented demonstrates that a careful case-by-case analysis is probably required in order to draw any firm conclusions. It should be noted that in the case where nonrenormalizable couplings give significant corrections to the trilinear $`(\theta ,\theta \omega ,\theta \omega ^2)`$ coupling, the above arguments likely fail.
Lastly, we would like to comment on another suggested source of complex phases. As noted above, physical states are constructed from linear combinations of states whose zero modes are given by different fixed tori. It was noted by Kobayashi and Ohtsubo that complex phases enter from the coefficients in the linear combinations and that these might be a souce of CP violation . However, this amounts to explicit CP violation since it does not require a particular scalar background. As explained in Section 1, explicit CP violation does not occur in heterotic orbifolds since CP is a gauge symmetry of the underlying theory. Therefore, the phases which enter into couplings from this source cannot contribute to the KM phase.
## 3 Modular invariance
In the toy model to be considered below, nonrenormalizable superpotential couplings will play a crucial role. In this section we present a set of assumptions for coupling coefficients of holomorphic couplings of arbitrary order; the result will be couplings which transform in the requisite manner under the $`[SL(2,𝐙)]^3`$ subgroup of the full $`SU(3,3,𝐙)`$ modular duality group of the $`Z_3`$ orbifold.
We begin by considering the simpler case of a two dimensional $`Z_3`$ orbifold, where there is a single modulus $`T`$ and there are only three fixed points. Consequently, twisted matter fields carry a single fixed point label. The twisted trilinear couplings are known in this simple case . These couplings can be expressed in terms of the Dedekind $`\eta `$ function
$$\eta (T)=e^{\pi T/12}\underset{n=1}{\overset{\mathrm{}}{}}(1e^{2\pi nT})$$
(3.1)
and the level-one $`SU(3)`$ characters
$$\chi _i(T)=\eta ^2(T)\underset{v\mathrm{\Gamma }_i}{}e^{\pi T|v|^2}.$$
(3.2)
In this expression, $`\mathrm{\Gamma }_0`$ is the $`SU(3)`$ root lattice, while $`\mathrm{\Gamma }_{1,2}`$ are shifted by $`SU(3)`$ weight vectors. Explicitly,
$$|v|^2=\left(\begin{array}{cc}n_1+\frac{i}{3},& n_2+\frac{i}{3}\end{array}\right)\left(\begin{array}{cc}2& 1\\ 1& 2\end{array}\right)\left(\begin{array}{c}n_1+i/3\\ n_2+i/3\end{array}\right)$$
(3.3)
with $`n_1`$ and $`n_2`$ integers to be summed over in (3.2) and $`i`$ is the integer labeling $`\mathrm{\Gamma }_i`$. It is easy to check that $`\chi _1=\chi _2`$. The values of these three functions at the self-dual points $`T=1,e^{i\pi /6}`$ are approximately given by Table 1.
The trilinear couplings between twisted fields $`\mathrm{\Phi }_1^{i_1},\mathrm{\Phi }_2^{i_2},\mathrm{\Phi }_3^{i_3}`$ in the superpotential are
$$\lambda [\eta (T)]^2f_{i_1i_2i_3}^T\mathrm{\Phi }_1^{i_1}\mathrm{\Phi }_2^{i_2}\mathrm{\Phi }_3^{i_3},$$
(3.4)
where $`f_{i_1i_2i_3}^T`$ is given by:
$$f_{i_1i_2i_3}^T=\{\begin{array}{c}\chi _0(T),i_1=i_2=i_3;\hfill \\ \chi _1(T),i_1i_2i_3i_1;\hfill \\ 0,\text{else}.\hfill \end{array}$$
(3.5)
The overall T-independent coupling strength $`\lambda `$ does not depend on fixed point locations and is obtained by factorization of the four-point string correlator . Note that $`\eta (T)`$ is superfluous and occurs only because of the definition of $`\chi _i(T)`$.
It can be checked that the above Yukawa couplings transform correctly under the target-space modular duality group $`SL(2,𝐙)`$. The Kähler potential for the modulus $`T`$ of the 2-d orbifold is given by $`K(T,\overline{T})=\mathrm{ln}(T+\overline{T})`$, which transforms under the $`SL(2,𝐙)`$ modular transformations
$$T\frac{aTib}{icT+d},adbc=1,a,b,c,d𝐙,$$
(3.6)
as
$$K(T,\overline{T})K(T,\overline{T})+F(T)+\overline{F}(\overline{T}),$$
(3.7)
where $`F(T)=\mathrm{ln}(icT+d)`$. We require that the quantity $`K+\mathrm{ln}|W|^2`$ remain invariant , which implies that the superpotential $`W`$ transform as
$$We^{i\gamma (a,b,c,d)}e^{F(T)}W,$$
(3.8)
where $`\gamma (a,b,c,d)`$ is a T-independent phase which does not appear in the functional which is physically meaningful, $`K+\mathrm{ln}|W|^2`$. We can take $`b=1,c=1,d=0`$ to obtain the $`T1/T`$ transformation of $`SL(2,𝐙)`$, so that $`F(T)=\mathrm{ln}(iT)`$ in this case. It has been shown that the twist fields $`\sigma _i`$ which create twisted vacua corresponding to fixed points labeled by $`i`$, and hence the twisted states $`\mathrm{\Phi }^i`$, transform under the duality transformation $`T1/T`$ as
$$\left(\begin{array}{c}\sigma _1^{}\\ \sigma _2^{}\\ \sigma _3^{}\end{array}\right)=\frac{e^{i\beta }}{\sqrt{3}}\left(\begin{array}{ccc}1& 1& 1\\ 1& \overline{\alpha }& \alpha \\ 1& \alpha & \overline{\alpha }\end{array}\right)\left(\begin{array}{c}\sigma _1\\ \sigma _2\\ \sigma _3\end{array}\right),$$
(3.9)
where $`\mathrm{exp}(3i\beta )\sqrt{\overline{T}/T}`$ and $`\alpha \mathrm{exp}(2\pi i/3)`$. Since the vertex operator which creates the twisted state $`\mathrm{\Phi }^i`$ is proportional to $`\sigma _i`$, the twisted fields must transform in a way which is proportional to the same matrix. The nonholomorphic phase $`\beta `$ must be absent in the supergravity definition of the fields. For example, we could define $`\stackrel{~}{\sigma }_i=e^{i\beta /2}\sigma _i`$ and use these to create the supergravity fields, which must transform in a holomorphic way. Aside from the mixing of different fixed points, the twisted fields tranform with a modular weight q: $`\mathrm{\Phi }\mathrm{\Phi }^{}e^{qF(T)}\mathrm{\Phi }`$. These weights are known ; a non-oscillator twisted field has modular weight $`q=2/3`$. Only non-oscillator twisted fields are allowed to enter into the trilinear twisted couplings due to the automorphism selection rule, corresponding to invariance of string correlators under automorphisms of the underlying $`SU(3)`$ lattice. This rule is explained and illustrated with examples in ref. . Since $`F(T)=\mathrm{ln}(iT)`$ for the $`T1/T`$ duality transformation, non-oscillator twisted fields must transform as
$$\left(\begin{array}{c}\mathrm{\Phi }_{}^{1}{}_{}{}^{}\\ \mathrm{\Phi }_{}^{2}{}_{}{}^{}\\ \mathrm{\Phi }_{}^{3}{}_{}{}^{}\end{array}\right)=\frac{1}{\sqrt{3}}\left(\frac{1}{iT}\right)^{2/3}\left(\begin{array}{ccc}1& 1& 1\\ 1& \overline{\alpha }& \alpha \\ 1& \alpha & \overline{\alpha }\end{array}\right)\left(\begin{array}{c}\mathrm{\Phi }^1\\ \mathrm{\Phi }^2\\ \mathrm{\Phi }^3\end{array}\right).$$
(3.10)
Under the duality transformation $`TT^{}=1/T`$ it has been shown that the $`SU(3)`$ characters transform as
$$\left(\begin{array}{c}\chi _0^{}\\ \chi _1^{}\\ \chi _2^{}\end{array}\right)=\frac{1}{\sqrt{3}}\left(\begin{array}{ccc}1& 1& 1\\ 1& \alpha & \overline{\alpha }\\ 1& \overline{\alpha }& \alpha \end{array}\right)\left(\begin{array}{c}\chi _0\\ \chi _1\\ \chi _2\end{array}\right).$$
(3.11)
It is also well-known that
$$\eta ^2(1/T)=\eta ^2(T^{})=T\eta ^2(T)=ie^{F(T)}\eta ^2(T).$$
(3.12)
We define a polynomial $`p`$ which encodes the superpotential couplings (3.4), up to a power of $`\eta (T)`$:
$`p(T;\mathrm{\Phi }_1,\mathrm{\Phi }_2,\mathrm{\Phi }_3)`$ $`=`$ $`\chi _0(T)(\mathrm{\Phi }_1^1\mathrm{\Phi }_2^1\mathrm{\Phi }_3^1+\mathrm{\Phi }_1^2\mathrm{\Phi }_2^2\mathrm{\Phi }_3^2+\mathrm{\Phi }_1^3\mathrm{\Phi }_2^3\mathrm{\Phi }_3^3)`$ (3.13)
$`+\chi _1(T)(\mathrm{\Phi }_1^1\mathrm{\Phi }_2^2\mathrm{\Phi }_3^3+\mathrm{\Phi }_1^3\mathrm{\Phi }_2^1\mathrm{\Phi }_3^2+\mathrm{\Phi }_1^2\mathrm{\Phi }_2^3\mathrm{\Phi }_3^1`$
$`+\mathrm{\Phi }_1^3\mathrm{\Phi }_2^2\mathrm{\Phi }_3^1+\mathrm{\Phi }_1^1\mathrm{\Phi }_2^3\mathrm{\Phi }_3^2+\mathrm{\Phi }_1^2\mathrm{\Phi }_2^1\mathrm{\Phi }_3^3).`$
Using the transformation properties enumerated above, it can be shown that
$$\eta ^2pi\left(\frac{1}{iT}\right)\eta ^2p=ie^{F(T)}\eta ^2p.$$
(3.14)
Thus, the functional $`\eta ^2p`$ transforms with modular weight $`1`$ up to a moduli independent phase $`i`$, as required by (3.8). Here, we draw attention to the fact that the monomials contained in (3.13) do not by themselves transform in the required way. Rather, it is the linear combination of fields in (3.13) together with $`\chi _i(T)`$ factors which is modular covariant, in the sense of (3.8).
Similar arguments hold for the axionic shift $`TT^{}=Ti`$. Indeed, $`\eta ^2(T)\mathrm{exp}(i\pi /6)\eta ^2(T)`$ and
$$\left(\begin{array}{c}\chi _0^{}\\ \chi _1^{}\\ \chi _2^{}\end{array}\right)=e^{i\pi /6}\left(\begin{array}{ccc}1& 0& 0\\ 0& \alpha & 0\\ 0& 0& \alpha \end{array}\right)\left(\begin{array}{c}\chi _0\\ \chi _1\\ \chi _2\end{array}\right),$$
(3.15)
$$\left(\begin{array}{c}\mathrm{\Phi }_{}^{1}{}_{}{}^{}\\ \mathrm{\Phi }_{}^{2}{}_{}{}^{}\\ \mathrm{\Phi }_{}^{3}{}_{}{}^{}\end{array}\right)=\left(\begin{array}{ccc}\overline{\alpha }& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right)\left(\begin{array}{c}\mathrm{\Phi }^1\\ \mathrm{\Phi }^2\\ \mathrm{\Phi }^3\end{array}\right).$$
(3.16)
Then it can be checked that $`\eta ^2p\eta ^2p`$, which transforms as it should (E.g., $`c=0,d=1`$ for the axionic shift).
A general $`SL(2,𝐙)`$ transformation (3.6) can be built up out of the two operations analyzed above. Thus, we are assured that the trilinear superpotential coupling (3.4) is modular covariant under the entire duality group. To summarize, for the general $`SL(2,𝐙)`$ transformation the polynomial (3.13) and the Dedekind $`\eta `$ function transform as
$`p(T;\mathrm{\Phi }_1,\mathrm{\Phi }_2,\mathrm{\Phi }_3)`$ $``$ $`e^{i\varphi (a,b,c,d)}e^{_{n=1}^3q_nF(T)}p(T;\mathrm{\Phi }_1,\mathrm{\Phi }_2,\mathrm{\Phi }_3),`$
$`\eta ^2(T)`$ $``$ $`e^{i\gamma (a,b,c,d)}e^{F(T)}\eta ^2(T),`$
where $`q_n`$ is the modular weight of the matter fields of species $`n`$, $`\mathrm{\Phi }_n^j`$. It can then be checked that (3.8) holds for the trilinear superpotential term coupling (3.4).
Twisted couplings of dimension $`3m>3`$ in the effective field theory remain to be calculated from the underlying conformal field theory. However, these computations appear formidable. As an example, we briefly consider the form of six dimensional twisted couplings using the methods of . The classical action $`S_{cl}`$ may determined from monodromy conditions on the underlying bosonic fields $`X(z,\overline{z}),\overline{X}(z,\overline{z})`$ of the 2-d orbifold, where $`z,\overline{z}`$ provide a parameterization of the string world-sheet; for each classical solution $`X_{cl}(z,\overline{z})`$, local monodromy conditions demand
$`_zX_{cl}(z)`$ $`=`$ $`{\displaystyle \frac{a(z_4,z_5,z_6)}{[z(z1)(zz_4)(zz_5)(zz_6)]^{2/3}}},`$
$`\overline{}_{\overline{z}}X_{cl}(\overline{z})`$ $`=`$ $`{\displaystyle \frac{b(\overline{z}_4,\overline{z}_5,\overline{z}_6)}{[\overline{z}(\overline{z}1)(\overline{z}\overline{z}_4)(\overline{z}\overline{z}_5)(\overline{z}\overline{z}_6)]^{1/3}}},`$
$`\overline{}_{\overline{z}}\overline{X}_{cl}(\overline{z})`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{a}(\overline{z}_4,\overline{z}_5,\overline{z}_6)}{[\overline{z}(\overline{z}1)(\overline{z}\overline{z}_4)(\overline{z}\overline{z}_5)(\overline{z}\overline{z}_6)]^{2/3}}},`$
$`_z\overline{X}_{cl}(z)`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{b}(z_4,z_5,z_6)}{[z(z1)(zz_4)(zz_5)(zz_6)]^{1/3}}}.`$ (3.17)
Here, $`z_4,z_5,z_6`$ are the three vertex insertion locations which cannot be fixed by $`SL(2,𝐂)`$ invariance while the first three vertices $`z_1,z_2,z_3`$ have been sent to $`0,1,\mathrm{}`$ resp. The functions $`a,b,\stackrel{~}{a},\stackrel{~}{b}`$ depend on the unfixed vertex locations and are determined for each classical solution $`X_{cl}`$ using global monodromy conditions. The classical action is given by
$$S_{cl}=\frac{1}{4\pi }d^2z(X_{cl}\overline{}\overline{X}_{cl}+\overline{}X_{cl}\overline{X}_{cl}).$$
(3.18)
Upon substitution of formulae (3.17) into (3.18), one finds it necessary to perform the integrals
$$I_r=\frac{d^2z}{|z(z1)(zz_4)(zz_5)(zz_6)|^r},r=\frac{4}{3},\frac{2}{3}.$$
(3.19)
Using techniques developed in , one can show that both of these integrals may be written in the form of sums of products of integrals along the real axis. The expressions obtained are typically complicated special functions of the unfixed vertex locations. Rather than attempting the calculation, we merely write the results symbolically as<sup>8</sup><sup>8</sup>8 At the risk of annoying the reader, we have explicitly shown the dependence on unfixed vertex locations, in order that we might stress wherein the difficulty lies.
$$I_r=\underset{i}{}c_r^iF_i(r;z_4,z_5,z_6)\overline{G}_i(r;\overline{z}_4,\overline{z}_5,\overline{z}_6).$$
(3.20)
It follows that
$$S_{cl}(z_4,z_5,z_6;\overline{z}_4,\overline{z}_5,\overline{z}_6)=\frac{1}{4\pi }\left(a\stackrel{~}{a}I_{4/3}+b\stackrel{~}{b}I_{2/3}\right)(z_4,z_5,z_6;\overline{z}_4,\overline{z}_5,\overline{z}_6).$$
(3.21)
The functions $`a,b,\stackrel{~}{a},\stackrel{~}{b}`$ are also complicated special functions of the unfixed vertex locations. It should be clear that $`S_{cl}`$ is a horrendous function. What is more, this action must be exponentiated, summed over an infinity of classical solutions $`X_{cl}`$, and multiplied by the quantum part of the partition function to obtain the correlator:
$`V_1(0,0)V_2(1,1)V_3(\mathrm{},\mathrm{})V_4(z_4,\overline{z}_4)V_5(z_5,\overline{z}_5)V_6(z_6,\overline{z}_6)`$
$`=Z(z_4,z_5,z_6;\overline{z}_4,\overline{z}_5,\overline{z}_6)`$
$`=Z_{qu}(z_4,z_5,z_6;\overline{z}_4,\overline{z}_5,\overline{z}_6){\displaystyle \underset{X_{cl}}{}}\mathrm{exp}[S_{cl}(z_4,z_5,z_6;\overline{z}_4,\overline{z}_5,\overline{z}_6)].`$
Here, the quantum part of the partition function $`Z_{qu}(z_4,z_5,z_6;\overline{z}_4,\overline{z}_5,\overline{z}_6)`$ will be some other horrific function of the unfixed vertex locations. Finally, we must extract the effective field theory coupling coefficient by integrating the unfixed vertex locations over the complex plane:
$$f_{i_1\mathrm{}i_6}d^2z_4d^2z_5d^2z_6Z(z_4,z_5,z_6;\overline{z}_4,\overline{z}_5,\overline{z}_6).$$
(3.23)
Integrating over exponentials of sums of products of special functions of several complex variables is bad enough, but one also has the $`Z_{qu}`$ prefactor and the functions $`a,b,\stackrel{~}{a},\stackrel{~}{b}`$ to deal with. Suffice it to say, the explicit calculation of coupling coefficients for higher dimensional twisted couplings certainly looks like a major undertaking.
As a result, we take a more phenomenlogical approach, using modular covariance as a guide. In this respect our effective field theory is “string-inspired” rather than “string-derived”. It is our hope that by appealing to symmetries of the underlying theory, we will capture the most important features of the bona fide couplings. Modular covariant $`3m`$-dimensional twisted couplings can be constructed by tensoring the polynomials (3.13). This leads us to the implicit definition of $`T`$-dependent twisted coupling coefficients $`f_{i_1\mathrm{}i_{3m}}^T`$ given by
$$\underset{\{i_j\}}{}f_{i_1\mathrm{}i_{3m}}^T\mathrm{\Phi }_1^{i_1}\mathrm{}\mathrm{\Phi }_{3m}^{i_{3m}}=\frac{1}{m!(3!)^m}\underset{\{n_j\}}{}^{}\underset{k=0}{\overset{m1}{}}p(T;\mathrm{\Phi }_{n_{3k+1}},\mathrm{\Phi }_{n_{3k+2}},\mathrm{\Phi }_{n_{3k+3}}),$$
(3.24)
Here, $`_{\{n_j\}}^{}`$ indicates that the $`3m`$-tuple of subscripts $`(n_1,\mathrm{},n_{3m})`$ should be summed over all permutations of $`(1,2,\mathrm{},3m)`$. The factor $`1/m!(3!)^m`$ accounts for trivial permutation symmetries. This construction treats the different species of fields $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_{3m}`$ in a symmetric way with respect to fixed point couplings. If we define the $`3m`$-dimensional twisted superpotential coupling as
$$\lambda f_{i_1\mathrm{}i_{3m}}^T[\eta (T)]^{2({\scriptscriptstyle q_n}1)}\mathrm{\Phi }_1^{i_1}\mathrm{}\mathrm{\Phi }_{3m}^{i_{3m}},$$
(3.25)
it can be checked that this will satisfy the requirement (3.8). We generically denote the overall modulus- and fixed-point-independent coupling strength by $`\lambda `$. The actual value of this strength will depend on the dimensionality of the coupling and the species of fields which enter. For $`m>1`$ in (3.25), it is likely that $`\lambda 1`$, due to the numerous world-sheet integrals which must be performed in (3.23) above, as pointed out recently by Cvetič, Everett and Wang . This aspect of nonrenormalizable couplings in string-derived models has been overlooked in much of the earlier literature, due to the temptation to estimate unknown quantities as $`𝒪(1)`$.
All of the above considerations dealt with a two-dimensional orbifold. We must generalize our results to the six-dimensional case. Also, it is necessary to say what should be done if untwisted states appear in a coupling or if some of the twisted states have nonzero oscillator numbers. Here, we address these complications only to the extent that it is necessary for the quark mass couplings in the toy model which is to be discussed in Section 4. Again, we take a phenomenological approach, using modular covariance as a guide, rather than attempting to explicitly derive effective field theory coupling coefficients from the underlying conformal field theory. The reason once again is that the calculations appear to be extremely difficult.
For a six-dimensional $`Z_3`$ orbifold the twist field operators generalize to $`\sigma _{ijk}(z,\overline{z})`$, where the triple $`ijk`$ denotes which of the 27 fixed points the twisted field $`\mathrm{\Phi }_n^{ijk}`$ created by the vertex operator sits at. The indices refer to the fixed point locations in each of the three complex planes. To the extent that the three complex planes are orthogonal to each other, which is the case if the off-diagonal T-moduli vanish, it is possible to decompose the twist operators into a tensor product of two-dimensional orbifold twist operators:
$$\sigma _{ijk}(z,\overline{z})\sigma _i^1(z,\overline{z})\sigma _j^2(z,\overline{z})\sigma _k^3(z,\overline{z}).$$
(3.26)
Then the six-dimensional orbifold couplings are a tensor product of the two-dimensional ones. Indeed, it can be checked that given the transformation properties of the two dimensional twist operators,<sup>9</sup><sup>9</sup>9Here, $`\stackrel{~}{\sigma }_i^I=e^{i\beta (T^I,\overline{T}^{\overline{I}})/2}\sigma _i^I`$ generalizes the $`\stackrel{~}{\sigma }_i`$ defined above for the two-dimensional orbifold case.
$`{\displaystyle \underset{I=1}{\overset{3}{}}}p(T^I;\stackrel{~}{\sigma }^I(z_1,\overline{z}_1)\stackrel{~}{\sigma }^I(z_2,\overline{z}_2)\stackrel{~}{\sigma }^I(z_3,\overline{z}_3))`$
$`{\displaystyle \underset{I=1}{\overset{3}{}}}e^{i\varphi (a^I,b^I,c^I,d^I)}p(T^I;\stackrel{~}{\sigma }^I(z_1,\overline{z}_1)\stackrel{~}{\sigma }^I(z_2,\overline{z}_2)\stackrel{~}{\sigma }^I(z_3,\overline{z}_3)).`$ (3.27)
This has the correct transformation, up to the factor $`_I\mathrm{exp}(_iq_i^IF(T^I))`$ which would come from the transformation properties of the matter field vertex operators not accounted for by $`\stackrel{~}{\sigma }_i^I`$. Then the generalization of (3.4) with the correct modular transformation properties is:
$$\lambda \left(\underset{I}{}[\eta (T^I)]^2f_{i_1^Ii_2^Ii_3^I}^{T^I}\right)\mathrm{\Phi }_1^{i_1^1i_1^2i_1^3}\mathrm{\Phi }_2^{i_2^1i_2^2i_2^3}\mathrm{\Phi }_3^{i_3^1i_3^2i_3^3}$$
(3.28)
It is easy to check that this holds for the higher dimensional couplings as well. We simply take the obvious products of 2-d orbifold coupling coefficients:
$$\lambda \left(\underset{I}{}[\eta (T^I)]^{2(_nq_n^I1)}f_{i_1^I\mathrm{}i_{3m}^I}^{T^I}\right)\mathrm{\Phi }_1^{i_1^1i_1^2i_1^3}\mathrm{}\mathrm{\Phi }_3^{i_{3m}^1i_{3m}^2i_{3m}^3}.$$
(3.29)
Now we consider the occurence of twisted oscillator fields and untwisted fields in higher dimensional couplings. The vertex operator for a twisted oscillator field is proportional to an excited twist operator. The excited twist operator can be written in terms of an ordinary twist operator and a factor of $`\overline{}\overline{X}^{\mathrm{}}`$, with $`\mathrm{}`$ depending on the oscillator direction . Then the classical action $`S_{cl}`$ is still computed in the presence of the same twist operators and we expect that the fixed point dependence should be the same as in the case where none of the twisted states were oscillators. Indeed, this has been found to be the case in $`Z_N`$ orbifolds where renormalizable couplings may involve oscillator fields . It has been shown that the modular weight of an $`N_L=1/3`$ oscillator state $`Y^{\mathrm{},ijk}`$ (where $`\mathrm{}`$ is the direction of the oscillator and $`ijk`$ specifies the fixed point location) is given by
$$q^I(Y^{\mathrm{},ijk})=(2/3,2/3,2/3)+\delta _{\mathrm{}}^I.$$
(3.30)
Since the vertex operator creating the twisted oscillator state is proportional to $`\stackrel{~}{\sigma }_i^1\stackrel{~}{\sigma }_j^2\stackrel{~}{\sigma }_k^3`$, we expect that the state transforms under $`T^I1/T^I`$ according to (3.10) if $`I\mathrm{}`$, while the power $`2/3`$ should be replaced by $`5/3`$ if $`I=\mathrm{}`$. Based on this assumption, the modular invariant couplings involving oscillator fields (which for the $`Z_3`$ orbifold are always higher dimensional couplings because of the automorphism selection rule) are obtained from (3.29) directly, with the oscillator nature of states reflected in the modular weights $`q_i^I`$ and a different overall strength $`\lambda `$ than would be obtained if the states were not oscillators, due to the presence of the additional operator $`\overline{}\overline{X}^{\mathrm{}}`$. Obviously, adding untwisted states to a coupling does not introduce any new twist operators. We can always choose a “picture” such that the vertex operator of the untwisted state $`U^i`$ goes like $`X^i`$. Then the change in the conformal field theory correlator will be completely in the quantum part and we expect the fixed point dependence to be unchanged. An untwisted state $`U^i`$ has modular weight $`q_i^I=\delta _i^I`$. As a result of these arguments, we conclude that the coefficients for a coupling with $`3m`$ twisted fields and $`n`$ untwisted fields can be read off from (3.29), only we must include the modular weights of the untwisted fields in the sums in the exponents of the $`\eta `$ functions, and the overall coupling strength $`\lambda `$ will be different than if no untwisted fields were in the coupling.
## 4 Toy model
Below the scale of $`U(1)_X`$ breaking, $`\mathrm{\Lambda }_X10^{17}`$ GeV, the quark and higgs spectrum is assumed to be that of the Minimal Supersymmetric Standard Model (MSSM). Extra color triplets get vector mass couplings when Xiggses get vevs. As a consequence, they get masses $`𝒪(\mathrm{\Lambda }_X)`$ and integrate out of the spectrum near the string scale. Above $`\mathrm{\Lambda }_X`$ we suppose the spectrum contains untwisted doublet quarks $`Q^i`$, untwisted $`u_L^c`$-like quarks $`u_1^i`$, twisted $`u_L^c`$-like quarks $`u_2^i`$, twisted $`d_L^c`$-like quarks $`d^i`$, untwisted $`H_u`$-like higgs doublets $`H_u^i`$ and twisted $`H_d`$-like higgs doublets $`H_d^i`$. These assignments are quite similar to those of the FIQS model. The superscript on untwisted fields corresponds to the H-momentum of the states in the underlying theory, and takes values $`i=1,2,3`$. The superscript on twisted fields corresponds to the fixed point location in the third complex plane of the six-dimensional compact space. As in the FIQS model, three linear combinations of the six $`u_L^c`$-like quarks survive in the low-energy spectrum, which we describe by mixing matrices $`X^1`$ and $`X^2`$:
$$u_1^i=X_{ij}^1u_{jL}^c+\text{heavy},u_2^i=X_{ij}^2u_{jL}^c+\text{heavy}.$$
(4.1)
The mixings to $`\mathrm{\Lambda }_X`$ scale mass eigenstates, denoted by “heavy”, are not important to our tree level analysis of low energy quark mass matrices. We assume that all extra higgses integrate out near the scale $`\mathrm{\Lambda }_X`$ due to vector couplings induced by the Xiggs vevs (as in the FIQS model), leaving one pair which we identify as the $`H_u`$ and $`H_d`$ of the MSSM:
$$H_u=H_u^1,H_d=H_d^3.$$
(4.2)
We introduce SM singlet Xiggses $`Y_1^\mathrm{}i,Y_2^\mathrm{}i`$ and $`Y_3^\mathrm{}i`$ which get $`𝒪(\mathrm{\Lambda }_X)`$ vevs and appear in the nonrenormalizable mass couplings of the quarks. The $`Y`$’s are charged under $`U(1)_X`$ and their scalar components are among those which get vevs to cancel the $`U(1)_X`$ FI term. The $`Y`$’s are twisted states with nonzero left-moving oscillator number $`N_L=1/3`$. The first superscript corresponds to the oscillator direction in the compact space, $`\mathrm{}=1,2,3`$. The second superscript corresponds to the fixed point location in the third complex plane of the compact space. Such fields also arise in the FIQS model. In what follows, we use the same symbols for superfields and scalar components of fields other than the quarks, with the meaning obvious by context. Similarly, whether we refer to a quark superfield or its fermionic component should be obvious by context.
We assume that the leading holomorphic couplings giving quarks masses are contained in the superpotential
$`W_{qm}`$ $`=`$ $`\lambda _0|ϵ_{ijk}|H_u^iQ^ju_1^k+\lambda _1\lambda _{i_1i_2i_3\mathrm{}_1\mathrm{}_2kj}Y_1^{\mathrm{}_1i_1}Y_1^{\mathrm{}_2i_2}H_u^kQ^ju_2^{i_3}`$ (4.3)
$`+\lambda _2\lambda _{i_1i_2i_3i_4i_5i_6\mathrm{}_1\mathrm{}_2\mathrm{}_3\mathrm{}_4j}Y_1^{\mathrm{}_1i_1}Y_2^{\mathrm{}_2i_2}Y_3^{\mathrm{}_3i_3}Y_3^{\mathrm{}_4i_4}H_d^{i_5}Q^jd^{i_6}.`$
The trilinear untwisted coupling is proportional to $`|ϵ_{ijk}|`$ according to the conservation of H-momentum orbifold selection rule. The fields $`Y_{1,2,3}^\mathrm{}i`$ are assumed to be charged identically under $`U(1)_X`$. We assume that the operators contained in (4.3) are each neutral under the full rank sixteen gauge group $`G`$ obtained from the orbifold embedding. Thus, the other fields are assumed to have $`U(1)_X`$ charges<sup>10</sup><sup>10</sup>10 In the three generation constructions presently under consideration, the $`U(1)_X`$ charges are independent of generation number. Thus, horizontal flavor symmetries, such as those considered in refs. , cannot be implemented in this context. such that each coupling is $`U(1)_X`$ neutral. We assume that off-diagonal T-moduli $`T^{I\overline{J}},I\overline{J},`$ have vanishing vevs, as in the FIQS model, so that the leading order kinetic terms for the matter fields are diagonal. In the FIQS model, nonvanishing off-diagonal T-moduli lead to nonvanishing F-terms which break supersymmetry at the scale $`\mathrm{\Lambda }_X`$, which is unacceptable. Though these fields may acquire vevs once supersymmetry is broken in the hidden sector, we expect the vevs to be at most of order the hidden sector supersymmetry breaking scale. As a result, off-diagonal T-moduli give negligible contributions to the kinetic terms of the quarks. We further assume that the diagonal T-moduli $`T^IT^{I\overline{I}}`$ are stabilized at one of their self-dual points $`T^I=1,e^{i\pi /6}`$ once supersymmetry is broken, consistent with models of hidden sector supersymmetry breaking by gaugino condensation . It has also been argued that the T-moduli may stabilize at other points on the unit circle . Either way, it would appear that string-derived scalar potentials for the T-moduli stabilize them to values $`|T^I|=1`$. For this reason we view models which allow $`T^I`$ as large as required to obtain hierarchies in the Yukawa couplings of twisted fields to be unmotivated.
The Kähler metric for matter fields in (0,2) $`Z_3`$ orbifolds (arbitrary Wilson lines and point group embeddings) has been determined to leading order . In the case of vanishing off-diagonal T-moduli, the metric of the untwisted fields $`Q^i`$ and $`u_1^i`$ is given by
$$K_Q_{i\overline{\mathrm{}}}=K_{u_1}_{i\overline{\mathrm{}}}=\delta _{i\overline{\mathrm{}}}T^i+\overline{T}^{\overline{ı}}^1.$$
(4.4)
We make redefinitions $`Q^iT^i+\overline{T}^{\overline{ı}}^{1/2}Q^i`$ and similarly for $`u_1^i`$. Similar arguments hold for the twisted fields $`u_2^i`$ and $`d^i`$, whose Kähler metric at leading order is
$$K_d_{i\overline{\mathrm{}}}=K_{u_2}_{i\overline{\mathrm{}}}=\delta _{i\overline{\mathrm{}}}\underset{I}{}T^I+\overline{T}^{\overline{I}}^{2/3}.$$
(4.5)
We assume that the overall coupling strengths $`\lambda _0,\lambda _1,\lambda _2`$ in (4.3) reflect these rescalings and that the quark fields entering these couplings are the rescaled ones. We also assume that the factor $`\mathrm{exp}K/2`$ has been absorbed into these coupling strengths and make use of the fact that $`D_iD_jWW_{ij}`$ is a very good approximation, in the supergravity lagrangian notation of Wess and Bagger . It can be checked that the terms which we drop are $`𝒪(m_{\stackrel{~}{G}}m_W/m_P)`$, where $`m_{\stackrel{~}{G}}`$ is the gravitino mass. When working with these rescaled fields, we may raise and lower indices with impunity since their Kähler metric is canonical in the leading order approximation made here. Once the fields $`u_1^i`$ and $`u_2^i`$ have been rescaled in this way, the mixings to mass eigenstates $`u_{iL}^c`$ and their three heavy relatives (all canonically normalized) can be made unitary. We then have as a constraint:
$$(X^1X^1)_{jk}+(X^2X^2)_{jk}=\delta _{jk}.$$
(4.6)
When one takes (4.1) and (4.2) into account, the effective Yukawa couplings for the quarks are given by
$`\lambda _{jm}^u`$ $`=`$ $`\lambda _0(\delta _{j2}X_{3m}^1+\delta _{j3}X_{2m}^1)+\lambda _1\lambda _{i_1i_2i_3\mathrm{}_1\mathrm{}_21j}Y_1^{\mathrm{}_1i_1}Y_1^{\mathrm{}_2i_2}X_{i_3m}^2,`$
$`\lambda _{jm}^d`$ $`=`$ $`\lambda _2\lambda _{i_1i_2i_3i_43m\mathrm{}_1\mathrm{}_2\mathrm{}_3\mathrm{}_4j}Y_1^{\mathrm{}_1i_1}Y_2^{\mathrm{}_2i_2}Y_3^{\mathrm{}_3i_3}Y_3^{\mathrm{}_4i_4}.`$ (4.7)
In going from (4.3) to (4.7), we have set $`k=1`$ in the second coupling of (4.3) because $`H_u=H_u^1`$ and we have fixed $`i_5=3`$ and $`i_6=m`$ in the third coupling of (4.3) since we couple to $`H_d=H_d^3`$ and $`d^m`$.
We now apply the assumptions of Section 3 to the toy model. The toy model is based on string-derived models where two “discrete” Wilson lines are included in the embedding to give a three generation model . By construction, states which differ only by their fixed point location in the third complex plane have identical gauge quantum numbers. On the other hand, states which differ by fixed point locations in the first two complex planes generally have different gauge quantum numbers under the rank 16 gauge group which survives the orbifold compactification. Typically, the embedding is arranged so that the rank 16 gauge group has the form $`SU(3)\times SU(2)\times [U(1)]^m\times G_c`$, where $`G_c`$ is a simple group which condenses in the hidden sector to break supersymmetry. The extra $`U(1)`$’s get broken down to $`U(1)_Y\times `$ (hidden $`U(1)`$’s) by the FI term associated with the anomalous $`U(1)_X`$. Fixed point locations in the first two complex planes become species labels. In what follows, the only fixed point location superscript on twisted states is that corresponding to the third complex plane. This serves as a family number for twisted states in these models. Twisted “relatives” differ only by their fixed point location in the third complex plane, so in many respects the effective Yukawas behave as if we were working with a two-dimensional orbifold. The coupling coefficients in (4.7) are given by:
$`\lambda _{i_1i_2i_3\mathrm{}_1\mathrm{}_21j}`$ $`=`$ $`\left({\displaystyle \underset{I}{}}\eta (T^I)^{2(Q_1^I1)}\right)\chi _0(T^1)\chi _0(T^2)f_{i_1i_2i_3}^{T^3}`$ (4.8)
$`\text{if}(\mathrm{}_1,\mathrm{}_2)=(\underset{¯}{1,j}),`$
$`=`$ $`0\text{else};`$
$`\lambda _{i_1i_2i_3i_43m\mathrm{}_1\mathrm{}_2\mathrm{}_3\mathrm{}_4j}`$ $`=`$ $`\left({\displaystyle \underset{I}{}}\eta (T^I)^{2(Q_2^I1)}\right)\chi _1(T^1)^2\chi _1(T^2)^2f_{i_1i_2i_3i_43m}^{T^3}`$ (4.9)
$`\text{if}(\mathrm{}_1,\mathrm{}_2,\mathrm{}_3,\mathrm{}_4)=(\underset{¯}{1,2,3,j}),`$
$`=`$ $`0\text{else};`$
$$Q_1^I=2+2\delta _1^I+2\delta _j^I,Q_2^I=5+2\delta _j^I.$$
(4.10)
The constraints on the allowed values of $`\mathrm{}_i`$ follow from the automorphism selection rule; underlining denotes that any permutation of entries is permitted. The coefficients $`f_{i_1i_2i_3}^T`$ carry the dependence on third complex plane fixed point locations of twisted fields appearing in the nonrenormalizable up-type quark mass coupling, and are given explicitly in (3.5) above. The third complex plane fixed point dependence for the down-type quark mass coupling follows from the six-dimensional twisted coupling, and is defined implicitly by (3.24). The six twist coupling coefficients $`f_{i_1i_2i_3i_4i_5i_6}^T`$ vanish by the lattice group selection rule unless $`i_1+\mathrm{}+i_6=0\text{mod}\mathrm{\hspace{0.33em}3}`$. It can be checked that the choices $`i_1,\mathrm{},i_6`$ satisfying this rule can be divided into four classes, depending on whether triples of the indices can be formed where the entries of the triples are either all the same (s) or all different (d). Members of the same class have identical values for $`f_{i_1i_2i_3i_4i_5i_6}^T`$. The nonvanishing values of $`f_{i_1i_2i_3i_4i_5i_6}^T`$ are given in Table 2, according to which of the four classes the indices belong to. A representative example $`(i_1\mathrm{}i_6)`$ for each class is given to avoid any confusion. The factor of $`\chi _1(T^1)^2\chi _1(T^2)^2`$ in (4.9) follows from an additional assumption of our model: the fixed point locations (of the six species of twisted fields in the down-type Yukawa coupling) in the first two complex planes are such that the lattice group selection rule in each of the two planes is satisfied in the (dd) way of Table 2.
The strengths $`\lambda _0,\lambda _1,\lambda _2`$, the mixing matrices $`X_{ij}^1`$, $`X_{ij}^2`$, and the background $`Y_{1,2,3}^\mathrm{}i`$ are treated as phenomenological parameters. We tune the couplings, mixings and vevs to values which yield a reasonable phenomenology. In principle, all of these quantities would be fixed by a full and complete analysis of the string-derived effective supergravity. However, some of the background fields in $`Y_{1,2,3}^\mathrm{}i`$ are D-moduli ; in order to fix these we must say how the D-moduli flat directions are lifted. In the reference just cited it was suggested how these flat directions may be lifted by nonpertubative effects in the hidden sector, via superpotential couplings of the D-moduli to hidden sector matter condensates. It should also be noted that the symmetries which give rise to the D-moduli are only valid for the classical scalar potential under the assumption of vanishing of F-terms for the D-moduli in the background. As a result, they are pseudo-Goldstone bosons and we expect that the D-moduli flat directions will also be lifted by loop corrections.<sup>11</sup><sup>11</sup>11We thank Korkut Bardakci for bringing this to our attention. Preliminary estimates show that the D-moduli get contributions to their masses of order the gravitino mass from either effect. We are curently exploring how the phases of the Xiggses may be fixed by these mechanisms and what effect this will have on the KM phase in models of the type discussed here. Our results will be presented elsewhere.
We do not intend to be exhaustive in our analysis of the phenomenology of (4.7). Rather, we would simply like to demonstrate that it is possible to obtain a quark phenomenology which is consistent with experimental data. The shortest route to this goal is to implement textures in the effective Yukawa couplings. In this way our scan over parameter space is biased toward viable models. We make use of the results of a recent analysis of viable mass matrices , though we will not impose the hermiticity constraint implemented there since we have no motivation for it in the present context. We impose the textures
$$\lambda ^u=f_t\left(\begin{array}{ccc}A_u\theta ^8& 0& C_u\theta ^4\\ 0& D_u\theta ^4& E_u\theta ^2\\ \stackrel{~}{C}_u\theta ^4& \stackrel{~}{E}_u\theta ^2& 1\end{array}\right),\lambda ^d=f_b\left(\begin{array}{ccc}0& B_d\theta ^3& 0\\ \stackrel{~}{B}_d\theta ^3& D_d\theta ^2& 0\\ 0& 0& 1\end{array}\right),$$
(4.11)
through an arrangement of the mixing matrices $`X_{ij}^1,X_{ij}^2`$ and vevs $`Y_{1,2,3}^\mathrm{}i`$. Here,
$$\theta V_{us}0.22,f_tm_t/H_u^0,f_bm_b/H_d^0,$$
(4.12)
$$A_u,C_u,\stackrel{~}{C}_u,D_u,E_u,\stackrel{~}{E}_u,B_d,\stackrel{~}{B}_d,D_d𝒪(1).$$
(4.13)
The forms (4.11) were obtained by imposing the following textures in $`X^{1,2}`$ and $`Y_{1,2,3}`$:
$$X^1=\left(\begin{array}{ccc}& & \\ r_1\theta ^4& r_2\theta ^2& r_3\\ 0& r_4\theta ^4& r_5\theta ^2\end{array}\right),X^2=\left(\begin{array}{ccc}& & \\ & & \\ s_1\theta ^4& 0& s_2\end{array}\right),$$
(4.14)
$$Y_1^\mathrm{}i=y_1\theta ^2\delta _1^{\mathrm{}}\delta _3^i,Y_2^\mathrm{}i=y_2\delta _2^{\mathrm{}}\delta _3^i,Y_3=\left(\begin{array}{ccc}y_{31}\theta ^3& 0& 0\\ y_{32}\theta ^2& \stackrel{~}{y}_{31}\theta ^3& 0\\ 0& 0& y_{33}\end{array}\right),$$
(4.15)
The elements denoted by $``$ in (4.14) are left unspecified since they do not appear in the effective Yukawa matrices. We assume that they are chosen such that (4.6) is satisfied, which is generally true provided $`|r_3|^2+|s_2|^2+|r_5|^2\theta ^41`$ because of the $`\theta ^n`$ suppressions on entries of the other columns. Up to this restriction, the quantities $`r_i,s_i,y_i,y_{ij},\stackrel{~}{y}_{31}𝒪(1)`$. We note that there is no inconsistency in having Xiggs vevs larger than the FI term, since the vevs of fields having opposite $`U(1)_X`$ charge can be played off against each other in the $`U(1)_X`$ D-term. As an example, in the FIQS model the Y-type Xiggses can be made arbitrarily large while maintaining D-flatness by simultaneously increasing some of the vevs of non-oscillator Xiggses which they denote by $`S_6^i`$. Of course at some point the nonlinear $`\sigma `$-model perturbation theory breaks down.
Given the assumptions enumerated above, the effective Yukawa matrices take the form ($`T_c^IT^I`$):
$$\lambda ^u=\lambda _0\left(\begin{array}{ccc}h_us_1\theta ^8& 0& h_us_2\theta ^4\\ 0& r_4\theta ^4& r_5\theta ^2\\ r_1\theta ^4& r_2\theta ^2& r_3\end{array}\right),$$
(4.16)
$$h_u\frac{\lambda _1}{\lambda _0}\eta (T_c^1)^{10}\eta (T_c^2)^2\eta (T_c^3)^2\chi _0(T_c^1)\chi _0(T_c^2)\chi _0(T_c^3)y_1^2,$$
(4.17)
$$\lambda ^d=h_d\left(\begin{array}{ccc}0& 2\eta (T_c^1)^4\chi _1(T_c^3)y_{31}\theta ^3& 0\\ 2\eta (T_c^2)^4\chi _1(T_c^3)\stackrel{~}{y}_{31}\theta ^3& 2\eta (T_c^2)^4\chi _1(T_c^3)y_{32}\theta ^2& 0\\ 0& 0& 5\eta (T_c^3)^4\chi _0(T_c^3)y_{33}\end{array}\right),$$
(4.18)
$$h_d2\lambda _2\left[\eta (T_c^1)\eta (T_c^2)\eta (T_c^3)\right]^8\chi _1(T_c^1)^2\chi _0(T_c^2)^2\chi _0(T_c^3)y_1y_2y_{33}\theta ^2.$$
(4.19)
The quantities $`B_d,\stackrel{~}{B}_d,D_d`$ in (4.11) can be varied independently by adjusting the ratios $`y_{31}/y_{33},\stackrel{~}{y}_{31}/y_{33},y_{32}/y_{33}`$. The ratio of heavy generation Yukawa eigenvalues $`f_b/f_t`$ can be varied independently of $`B_d,\stackrel{~}{B}_d,D_d`$ and $`\lambda ^d`$ by adjusting $`y_2y_{33}^2/y_1`$. The top quark Yukawa eigenvalue $`f_t`$ can be adjusted independently by varying $`r_3`$. However, if $`\lambda _0`$ is too small there may be a minimum $`\mathrm{tan}\beta `$ below which we cannot match experimental data, since $`|r_3|<1`$ is required by (4.6). Recall that we have absorbed a factor $`\mathrm{exp}K/2`$ into $`\lambda _0`$, as well as the effects of quark field rescalings to account for noncanonical kinetic terms. Typically, $`\mathrm{exp}K/2<1`$ when the string moduli get $`𝒪(1)`$ vevs, so this may be a worry. Without an explicit model of the superpotential couplings and Xiggs vevs which determine the mixing matrices $`X^{1,2}`$, it is not possible to say whether or not the entries of $`\lambda _u`$ can be varied independently of each other and $`\lambda _d`$; we will assume that this is true.
With the above assumptions, scanning over the Xiggs vevs and the mixing coefficients $`r_i,s_i`$ for viable models is equivalent to varying the coefficients in (4.13) independently and tuning the values of $`f_t,f_b`$ to agree with experimental data. We then rephase the quarks according to the convention
$$V_{ud}>0,V_{us}>0,V_{cb}>0,V_{ts}<0,V_{cd}<0,$$
(4.20)
to which the Wolfenstein parameterization is an approximation. As is well known, the advantage of such a parameterization is that the elements with significant complex phase are the smallest ones, $`V_{ub}`$ and $`V_{td}`$.
All of these calculations are done at the $`U(1)_X`$ breaking scale, and are therefore subject to evolution under the renormalization group. The evolution of the quark masses and mixing angles assuming the MSSM spectrum has been studied extensively; approximate analytic formulas are available, for example in refs. . We will use the approximations of to evolve the low energy data to the scale of $`U(1)_X`$ breaking, which we assume to be $`\mathrm{\Lambda }_X\mathrm{\Lambda }_s5\times 10^{17}`$ GeV, based on what occurs in the FIQS model.<sup>12</sup><sup>12</sup>12See Appendix A. The following quantities are approximately scale-independent:
$$\frac{m_c}{m_u},\frac{m_s}{m_d},|V_{ud}|,|V_{cs}|,|V_{tb}|,|V_{us}|,|V_{cd}|.$$
(4.21)
The running of the other quantities is approximately given by:
$$\frac{m_t}{m_c}|_{\mathrm{\Lambda }_X}=\frac{1}{\xi _t^3\xi _b}\frac{m_t}{m_c}|_{M_Z},\frac{m_b}{m_s}|_{\mathrm{\Lambda }_X}=\frac{1}{\xi _t\xi _b^3}\frac{m_b}{m_s}|_{M_Z},|V_{cb}|_{\mathrm{\Lambda }_X}=\xi _t\xi _b|V_{cb}|_{M_Z};$$
(4.22)
the quantities $`|V_{ub}|,|V_{td}|,|V_{ts}|`$ scale in the same manner as $`|V_{cb}|`$. The scaling functions are given by
$$\xi _{t,b}=\mathrm{exp}\left[\frac{1}{16\pi ^2}_0^{\mathrm{ln}(\mathrm{\Lambda }_X/M_Z)}𝑑\chi f_{t,b}^2(\chi )\right],$$
(4.23)
where $`\chi =\mathrm{ln}(\mu /M_Z)`$ and $`f_{t,b}(\mu )`$ are the Yukawa coupling eigenvalues of the top and bottom quarks appearing in (4.11), at the scale $`\mu `$. We assume the scale of observable sector supersymmetry breaking is 1 TeV and we set $`\mathrm{tan}\beta =5`$. For low energy data we use the values of the running quark masses at the scale $`M_Z`$ as determined in ref. and the CKM data listed in ref. . Taking into account the errors quoted in these two source, we find the following values:
$$f_t(\mathrm{\Lambda }_X)=0.74_{0.24}^{+1.65},f_b(\mathrm{\Lambda }_X)=0.028(4),$$
$$\frac{m_t}{m_c}|_{\mathrm{\Lambda }_X}=440_{100}^{+390},\frac{m_c}{m_u}|_{\mathrm{\Lambda }_X}=290(60),$$
$$\frac{m_b}{m_s}|_{\mathrm{\Lambda }_X}=38(7),\frac{m_s}{m_d}|_{\mathrm{\Lambda }_X}=20(4),$$
$$|V_{CKM}|_{\mathrm{\Lambda }_X}=\left(\begin{array}{ccc}0.9752(8)& 0.220(4)& 0.0027(12)\\ 0.220(4)& 0.9745(8)& 0.033(4)\\ 0.014(11)& 0.065(36)& 0.9992(2)\end{array}\right).$$
(4.24)
We stress that theoretical errors due to the approximations made in (4.21) and (4.22) have not been included in the estimates of uncertainty. However, for our purposes this is not an important issue since we can always make a small shift in the $`𝒪(1)`$ parameters of our toy model to account for small corrections and larger uncertainties will just mean that more points in parameter space will give viable models.
In our analysis we consider both generic and extreme possibilities in order to get a feel for how the KM phase depends on the various sources of phases in (4.16) and (4.18).
Case 1: generic mixings and Xiggs vevs
We have scanned over the magnitudes of the parameters in (4.13) with Gaussian distributions centered on values suggested by the central values in (4.24) and with spreads suggested by the estimated uncertainties. Phases have been scanned on a flat distribution over the interval $`(\pi ,\pi ]`$. We then compared mass ratios and the magnitudes of CKM elements, except $`|V_{ub}|`$, to the values in (4.24); if these results agreed with the values in (4.24), except $`|V_{ub}|`$, up to the stated uncertainties, we stored the values of $`V_{ub}(\mathrm{\Lambda }_X)`$. We then scaled the magnitude of $`V_{ub}`$ according to (4.22) but left the phase unrotated to get an estimate of $`V_{ub}(M_Z)`$. In Figure 1 we plot our results, showing only points near the acceptable region. The results are hardly surprising: if we allow the phases of the fields Xiggses to float randomly, the KM phase can take on any value we like. No magical cancellation occurs. Although regions of parameter space in this toy model do exist which have reasonable quark masses, mixings and CP violation, the model provides no understanding of why we live in one region of parameter space rather than another. All that can be said is that our model, which contains many more free parameters than the number of experimental data points which we are attempting to fit, can be made to agree with what is known about the quark sector. One promising point does emerge, however. Figure 1 shows that CP violation is generic in the toy model under consideration. To be fair, one could argue that we have gone through a lot of unnecessary work to prove the obvious: if nonrenormalizable couplings contribute significantly to the effective quark Yukawa matrices, and the Xiggses in these nonrenormalizable couplings get complex vevs, then CP violation is to be expected. However, as we discussed in Section 1, one can wonder whether the symmetry constraints of modular invariance and orbifold selection rules might render these phases spurious. We have explicitly shown that this is not the case.
Case 2: complex Xiggs vevs
Here, we make the quantities $`r_i,s_i`$ in (4.14) real and positive and keep $`T_c^I=1`$ in order to isolate the effects of the phases of the Xiggs. With these assumptions it can be seen from (4.16) and (4.18) that the $`𝒪(1)`$ coefficients (4.13) satisfy
$$\mathrm{arg}D_u=\mathrm{arg}E_u=\mathrm{arg}\stackrel{~}{C}_u=\mathrm{arg}\stackrel{~}{E}_u=0,$$
(4.25)
$$\mathrm{arg}A_u=\mathrm{arg}C_u,$$
(4.26)
with $`\mathrm{arg}C_u,\mathrm{arg}B_d,\mathrm{arg}\stackrel{~}{B}_d`$ and $`\mathrm{arg}D_d`$ independent parameters to be scanned over. As in the previous case, we scan over the $`𝒪(1)`$ magnitudes of the coefficients (4.13) using a Gaussian distribution, and plot values of $`V_{ub}(M_Z)`$ for models which satisfy all constraints in (4.24) except the one on $`|V_{ub}|_{\mathrm{\Lambda }_X}`$. The results are given in Figure 2. Comparing to Figure 1, it can be seen that whether or not the mixings $`X^{1,2}`$ are a source of phase makes little difference. Complex Xiggs vevs provide a source of a KM phase and allow us to obtain any value we like.
Case 3: complex T-moduli
As discussed above, some of the T-moduli may stabilize at $`e^{i\pi /6}`$, their other self-dual point under $`SL(2,𝐙)`$. If we keep the mixing matrices $`X^{1,2}`$ and Xiggs vevs $`Y_{1,2,3}^\mathrm{}i`$ complex, then the results are indistinguishable from those of Fig. 1. To isolate the effect of T-moduli sitting at the other self-dual point, we have constrained the coefficients $`r_i,s_i`$ in (4.14) and the Xiggs vevs $`Y_{1,2,3}^\mathrm{}i`$ to be real and positive in what follows. Next, referring to (4.11), (4.16) and (4.18), we define
$`\gamma _1`$ $``$ $`\mathrm{arg}h_u=\mathrm{arg}A_u=\mathrm{arg}C_u,`$ (4.27)
$`\gamma _2`$ $``$ $`\mathrm{arg}B_d=\mathrm{arg}{\displaystyle \frac{\eta (T_c^1)^4\chi _1(T_c^3)}{\eta (T_c^3)^4\chi _0(T_c^3)}},`$ (4.28)
$`\gamma _3`$ $``$ $`\mathrm{arg}\stackrel{~}{B}_d=\mathrm{arg}D_d=\mathrm{arg}{\displaystyle \frac{\eta (T_c^2)^4\chi _1(T_c^3)}{\eta (T_c^3)^4\chi _0(T_c^3)}},`$ (4.29)
$`\mathrm{\Gamma }`$ $``$ $`\gamma _1\gamma _2+\gamma _3.`$ (4.30)
It can be checked that the Yukawa matrices (4.16) and (4.18) can be rephased such that $`\lambda ^d`$ has all positive entries and
$$\mathrm{arg}\lambda ^u=\left(\begin{array}{ccc}\mathrm{\Gamma }& 0& \mathrm{\Gamma }\\ 0& 0& 0\\ 0& 0& 0\end{array}\right).$$
(4.31)
This is easily implemented in a scan of the parameters in (4.13) by requiring all of them to be positive except $`\mathrm{arg}A_u=\mathrm{arg}C_u=\mathrm{\Gamma }`$. Using Table 1 it is straightforward to determine $`\mathrm{\Gamma }`$. We summarize the possible values in Table 3.
The results of the scan are presented in Figure 3. Once again, these are values of $`V_{ub}(M_Z)`$ for models which satisfy all constraints in (4.24) except the one on $`|V_{ub}|_{\mathrm{\Lambda }_X}`$.
It can be seen that neither possibility is consistent with the experimentally preferred region. This result only rules out T-moduli as the sole source of CP violating phases in the toy model considered here. In another model the powers of $`\eta (T),\chi _0(T),\chi _1(T)`$ would likely enter differently, since these depend on the dimension of a given nonrenormalizable coupling. It should also be noted that whereas we have set the mixing matrices $`X^{1,2}`$ real, they typically have a nontrivial dependence on $`\mathrm{arg}T_c^I`$, since they are determined at least in part by couplings involving twisted fields. Thus, the results of Figure 3 would likely change if this part of the model were made explicit.
## 5 Conclusions
In this article we have discussed several possible sources of CP violation in semi-realistic heterotic orbifold models. We have presented examples where CP violation does not occur in spite of the presence of phases, derived from complex string moduli vevs, in renormalizable coupling coefficients. However, it was described how nonrenormalizable couplings give a significant contribution to the effective quark Yukawa matrices when an anomalous $`U(1)_X`$ is present and we argued that this generically leads to a nontrivial KM phase.
In order to make a detailed analysis of models with nonrenormalizable couplings, we introduced modular covariant nonrenormalizable superpotential couplings. It was explained why it is difficult to obtain the bona fide effective coupling coefficients from conformal field theory techniques. It was also pointed out that higher order terms in the Kähler potential should be important in cases where an anomalous $`U(1)_X`$ is present. These theoretical uncertainties represent a significant stumbling block to further progress in string-derived effective supergravity models and it is hoped that they will be resolved at some point in the near future.
The KM phase was determined explicitly in a toy model inspired by three-generation heterotic $`Z_3`$ orbifold constructions. Though target space modular invariance and orbifold selection rules greatly restrict the coupling coefficients of nonrenormalizable couplings, we found it possible to obtain viable Yukawa couplings for quarks by adjusting the vevs of Xiggses. This result highlights the necessity of understanding how D-moduli flat directions are lifted in a given model. In principle the Xiggs vevs should be determined by the mechanisms which lift these flat directions. This would eliminate our ability to tune the scalar background to our liking and would in most cases probably render the quark phenomenology inconsistent with low energy data. We are currently investigating this issue and hope to report on it in a future publication.
Acknowledgements
The author would like to thank Prof. Mary K. Gaillard for innumerable discussions and helpful comments during the development of the work contained here. I am particularly indebted to her for showing me how the Green-Schwarz cancellation of the $`U(1)_X`$ anomaly works in the linear multiplet formulation and have borrowed heavily from written communications with her in preparing the Appendix. I would also like to thank Brent Nelson for useful comments. This work was supported in part by the Director, Office of Science, Office of High Energy and Nuclear Physics, Division of High Energy Physics of the U.S. Department of Energy under Contract DE-AC03-76SF00098 and in part by the National Science Foundation under grant PHY-95-14797.
Appendix
## Appendix A Charge normalization
In the FIQS model, unconventional normalizations for the $`U(1)`$ charges have been chosen to keep the tables of charges simple and amenable to computer assisted analysis. The generator $`Q_X^I`$ acting on the $`E_8\times E_8`$ root torus is given by
$$Q_X=6(0,0,0,0,0,0,0,0;1,1,1,0,0,0,0,0).$$
(A.1)
The affine level of a $`U(1)`$ group may be defined as
$$k_Q=2\underset{I=1}{\overset{16}{}}(Q^I)^2.$$
(A.2)
With this convention, the FIQS normalization gives $`k_{Q_X}=6^3`$. To go to a normalization where the coupling constant for the $`U(1)_X`$ group will be the universal coupling at the string scale, we must rescale the generator $`Q_XQ_X^{}`$ so that $`k_{Q_X^{}}=1`$. Then
$$Q_X^{}=\frac{1}{6\sqrt{6}}Q_X.$$
(A.3)
Since the original normalization satisfied<sup>13</sup><sup>13</sup>13The nonabelian generators $`T^a`$ are normalized such that $`\text{tr}T^aT^a=1/2`$ for a fundamental representation of $`SU(N)`$.
$$\text{tr}Q_X^3=27\text{tr}Q_X=2724\text{tr}T^aT^aQ_X=272454,$$
(A.4)
it can be checked that the rescaled generator satisfies
$$24\text{tr}(T^aT^aQ_X^{})=8\text{tr}Q_{X}^{}{}_{}{}^{3}=\text{tr}Q_X^{}=36\sqrt{6},$$
(A.5)
as required by anomaly matching . Indeed, if the $`U(1)_X`$ vector superfield $`V_X`$ is shifted by $`\delta V_X=(1/2)(\mathrm{\Lambda }+\overline{\mathrm{\Lambda }})`$, then the resulting anomalous transformation of the lagrangian is
$$\delta =\frac{1}{16\pi ^2}\underset{a}{}\text{tr}(T^aT^aQ_X^{})\left[\text{Re}\lambda F^aF^a+\text{Im}\lambda F^a\stackrel{~}{F}^a\right]+\mathrm{}$$
(A.6)
where $`\lambda =\mathrm{\Lambda }|`$. We introduce our counterterm<sup>14</sup><sup>14</sup>14We work in Kähler superspace and use the linear multiplet formulation where $`V`$ is a real superfield which satisfies modified linearity conditions and contains the dilaton $`\mathrm{}`$ as its lowest component . as
$$_{GS,V_X}=\delta _XEVV_X$$
(A.7)
from which it follows that under the shift in $`V_X`$
$$\delta _{GS,V_X}=\frac{\delta _X}{2}EV(\mathrm{\Lambda }+\overline{\mathrm{\Lambda }})$$
(A.8)
which when we go to components yields
$$\delta _{GS,V_X}=\frac{\delta _X}{8}\underset{a}{}\left(\text{Re}\lambda F^aF^a+\text{Im}\lambda F^a\stackrel{~}{F}^a\right)+\mathrm{}$$
(A.9)
The anomaly is cancelled if we choose
$$\delta _X=\frac{1}{2\pi ^2}\text{tr}T^aT^aQ_X^{}.$$
(A.10)
When combined with other terms in the lagrangian, the component form of (A.7) gives
$$D_X=\underset{i}{}q_X^i\widehat{K}_i\varphi ^i+\frac{\delta _X}{2}\mathrm{}\underset{i}{}q_X^i\widehat{K}_i\varphi ^i+\xi .$$
(A.11)
From this, we see that the FI term $`\xi `$ is given by
$$\xi =(2\mathrm{})\frac{\delta _X}{4}=\frac{2\mathrm{}}{8\pi ^2}\text{tr}T^aT^aQ_X^{}.$$
(A.12)
With the $`U(1)_X`$ generator chosen such that $`k_{Q_X^{}}=1`$, equation (A.5) gives
$$\xi =\frac{2\mathrm{}}{192\pi ^2}\text{tr}Q_X^{},$$
(A.13)
which may be recognized as the form typically quoted in the literature once it is realized that if we neglect nonperturbative corrections to the Kähler potential of the dilaton $`\mathrm{}`$, the universal coupling constant at the string scale is given by $`g^2=2\mathrm{}`$. In the FIQS normalization, the FI term is given by
$$\xi =\frac{2\mathrm{}}{192\pi ^2}\frac{1}{6\sqrt{6}}\text{tr}Q_X$$
(A.14)
which gives a significantly smaller number than if we had not accounted for the unconventional normalization of the $`U(1)_X`$ charge. In the FIQS model $`\text{tr}Q_X=1296`$, yielding
$$\xi 2\mathrm{}\times 4.7\times 10^25\times 10^2,$$
(A.15)
where we have used $`2\mathrm{}g^2`$, and $`0.5<g^2<1`$. The scale of $`U(1)_X`$ breaking is given by $`\mathrm{\Lambda }_X\sqrt{\xi }0.22m_P5\times 10^{17}`$ GeV $`\mathrm{\Lambda }_s`$, the string scale.
One must also take proper account of charge normalization for the SM hypercharge, as was pointed out in ref. . For example, in the FIQS model $`k_Y=11/3`$. Then the charge generator which will have the unified coupling at the string scale is $`Y^{}=\sqrt{3/11}Y`$. This is to be compared with the $`G_{GUT}SU(5)`$ relative factor of $`\sqrt{3/5}`$. Thus, the boundary value of the properly normalized hypercharge coupling $`g^{}`$ at the electroweak scale in the FIQS model is related to the one usually used in GUT-inspired renormalization group evolution of the couplings in the MSSM by
$$g^{}(\text{FIQS})|_{M_Z}=\sqrt{11/5}g^{}(\text{MSSM})|_{M_Z}.$$
(A.16)
This clearly does violence to unification of the couplings. In short, it is necessary to include hypercharge normalization among the criteria to be checked when searching for viable string-derived models.
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# Baryon Transport in Dual Models and the Possibility of a Backward Peak in Diffraction
## 1 Baryon transfer in particle scattering
Available experimental data
To observe the dynamics of baryon transfers over large rapidity gaps the best data involve the stopping of leading baryons to central rapidities. This requires to somehow identify the initial baryons distribution. One has to rely on the hypothesis - which presumably holds to a satisfactory degree - that the produced “sea” baryons and anti-baryons have an identical contribution and that a simple subtraction yields the distribution of the net initial baryon charge.
Unfortunately this excludes $`p\overline{p}`$ scattering, where the sea-baryon contribution cannot be determined from data. Unfortunately this precludes the use of post ISR data excepting HERA and diffractive systems with sufficiently large Pomeron-proton subenergies. The data are therefore typically quite old. Available are spectra from meson-baryon processes (compiled in ) and from ISR (compiled in ). At ISR a suitable combination of proton-proton and proton-antiproton processes can be used to separate forward and backward component of the incoming charge with a reasonable factorization assumption. An old plot of the leading baryon charge shown below is found to be consistent with a slope (in y) of $`\alpha _{Transfer}\alpha _{Pomeron}=1`$ with large error. That the central data points are somewhat on the high side can be taken as a hint of an eventual turnover to a flatter value.
To reduce the error one can integrate over the transverse momenta and to minimize the systematical error one considers the ratio:
$$A_{ISR}=\frac{\rho _{initialbaryoncharge}(y)}{\rho _{seabaryoncharge}(y)}=0.39\pm 0.05,\mathrm{\hspace{0.25em}0.33}\pm 0.05,\mathrm{\hspace{0.25em}0.23}\pm 0.05$$
given in a graph of for $`y=0.4`$, $`y=0`$ and $`y=+0.4`$. The central derivative of $`A_{ISR}`$ obtains no contribution from the (symmetric) sea-baryon distribution. Assuming the usual exponential distribution $`A\mathrm{exp}[(\alpha _{Transfer}\alpha _{Pomeron})y]`$ the quantity
$$\frac{d/dyA_{ISR}}{A_{ISR}}|_{y=0}=\alpha _{Transfer}\alpha _{Pomeron}=0.49_{0.37}^{+0.42}$$
(1)
just yields the slope. While not yet in contradiction with the initial intercept the preferred value is now considerably less.
New preliminary data come from the H1 experiment at HERA . They observed the initial baryons asymmetry at laboratory rapidities
$$A_{H1}=\frac{\rho _{initialbaryoncharge}(y)}{\rho _{seabaryoncharge}(y)}=0.08\pm 0.01\pm 0.025$$
(2)
A simple extrapolation of the ISR value accounting for the larger rapidity range spanned would have "predicted":
$$A_{H1}=0.061_{0.046}^{+0.243}$$
(3)
Hence the H1 values lie within the expected range. The trajectories required by the HERA ratio compared with its ISR value is now (see also )
$$\alpha _{Transfer}\alpha _{Pomeron}=0.4\pm 0.2.$$
If the preliminary data are finalized, this confirms the flattening of the trajectory.
### Dual Topological picture
The slowing down of the baryons is determined by the intercepts of the relevant exchanges. The basic philosophy of the Dual Topological model in the classification of such exchanges involves “materializing” or “suppressed” strings. “Materializing” means that the initial color fields are neutralized by a chain of hadronizing $`qq`$ pairs, “suppressed” means hadron-less neutralization by an exchange of a single quark. It is well-known for the Pomeron and the Reggeon exchanges where the cut Pomeron contains two materializing strings while the cut Reggeon contains one suppressed and one materializing string. As a general rule contributions with various suppressed strings have to be considered as independent and additive. For each suppressed string an extra factor $`(\sqrt{1/M_{string}})`$ appears and restricts the suppressed contribution to low energies or narrow rapidity ranges.
For a nuclear exchange one starts with a completely suppressed exchange, i. e. with the square of the quasi-elastic nucleon exchange amplitude
$$\alpha _{junction}^{III}1=2(<\alpha _{Nucleon}>1)=2$$
(4)
known from elastic backward scattering. Each of the three exchanged valence quarks can now be replaced by a “materializing” string. Corresponding to three, two, one or zero strings there are four contributions with trajectories spaced by one half. At considered energies the leading two of these “baryonium” trajectories with two and three hadronizing strings
$$\alpha _{junction}^01=0.5,\alpha _{junction}^I1=1.0$$
(5)
will be relevant. They could be responsible for the initially steep ($`1.0`$) and then possible flattening ($`0.5`$) slope observed in the data discussed above.
However the value of the final trajectory is rather uncertain. Values of $`\alpha _{junction}^I1=0.8\mathrm{}\mathrm{\hspace{0.25em}0}`$ were proposed in the literature . The correspondence to the Odderon discussed below will give support to a flat value.
### Implementation in Dual Parton model based Monte Carlo codes
At present energies nonleading Regge exchanges will have no strong effect on cross sections. To understand the final state structure local baryonium exchanges within a global Pomeron exchange have to be considered. A factorization among strings allows to ignore the quark string which is common to both trajectories and the inclusion of such exchanges is therefore straightforward:
The transition of the remaining diquark string (baryonium remnant) into an antiquark string (Pomeron remnant) can be implemented in a usual fragmentation scheme by a suitable choice of the splitting function involving quark and diquark transitions. It was implemented in most string models e.g. and it is part of the JETSET program (as diquarks or as pop-corn mechanism ). Without relying on string factorization leading- and sea-baryon exchanges are also implemented in HIJING/$`B\overline{B}`$ .
## 2 Baryon enhancement <br>in dense heavy ion scattering
### Concepts for slowing-down initial baryons
There are a number of conventional mechanisms enhancing baryon transfer and central baryon production in multiple scattering processes for string models. They are helpful in some regions but not enough to explain the large stopping in heavy ion scattering .
To understand the data it seems necessary to include interplay of string if they get sufficiently dense in transverse space. It was proposed that there are new special strings . In contrast, we shall maintain here the general factorization hypothesis between initial scattering in the quark phase and the final hadronization within standard strings.
Consider an incoming baryon. The usual Pomeron exchange in the Dual Parton model leaves a quark and a diquark for the string ends. Diquarks are no special entities and multiple scattering processes have no reason not to split them in a conventional two Pomeron interaction . It is natural to expect that diquark break-ups considerably slow down the baryons evolving. The probability for such an essentially unabsorbed process is :
$$[breakup]/[nobreakup][cutPomeronnumber]1$$
(6)
As required by the experimentally observed slowdown this is a drastic effect for heavy ion scattering while for hadron-hadron scattering multiple scattering are sufficiently rare to preserve the known hadron-hadron phenomenology. How such processes are affected was considered numerically in and no manifestly disturbing effects were found.
We emphasize that the behavior of the baryon quantum number slowed down by such a break-up is not trivial. In topological models the baryon contains Y-shaped color electric fluxes. Two Pomerons intercepting two different branches will leave two “free” valence quarks and a valence quark connected with the vortex line with the velocity of initial baryon which will subsequently form the end of the strings. The energy distribution of quarks with vortex lines (or of the fully separated vortex lines) in the structure function is a priory not known.
### Special baryon transfers in the Topological model
For a more detailed description of the slowing down we turn to the Dual Topological model introduced above. A discussion of baryon transfers in such a framework was recently given by Kharzeev . We will here emphasize topological aspects.
In topological models a Pomeron exchange corresponds to a cylinder connecting the two scattering hadrons. If one considers an arbitrary plane intersecting this exchange the intersection of the cylinder is topologically a circle. More specifically amplitudes with clockwise respectively anticlockwise orientation have to be added or subtracted depending on the charge parity. The cylinders or the circles therefore come with two orientations. This distinction is usually not very important as it is always topologically possible to attach hadrons in a matching way; except for $`C`$-parity conservation no special restrictions result.
Pomerons have a transverse extent and if they get close in transverse space they should interact. Hadronic interaction is sufficiently strong to be largely determined by geometry. It is therefore reasonable to expect that the coupling does not strongly depend on the orientation as long as there is no mechanism of suppression.
The two distinct configurations lead to different interactions. Two Pomerons with the same orientation can if they touch (starting locally at one point in the exchange-channel time) shorten their circumference and form a single circle:
This then corresponds to the usual triple Pomeron coupling experimentally well known from diffractive processes.
For two Pomerons with opposite orientation the situation is more complicated. Like for soap bubbles the two surfaces which get in contact can merge and form a single membrane. The joining inverts the orientation of the membrane. On the intersecting plane one now obtains – instead of the single circle – three lines originating in a vortex point and ending in an anti-vortex point as shown below:
Lacking a topological name for the object the term membraned cylinder will be used in the following.
How do this membraned cylinder contribute to particle production? Similar to the triple Pomeron case there are three different ways to cut through a membraned cylinder:
The cut numbered 1 which also intersects the membrane has vortex lines on both sides. They present a topological description of the baryon transfers considered above. By symmetry they contribute with a positive sign. Cuts which intersect only two sheets numbered 2 contribute to the two string contribution. Their sign is unknown. As they contain a closed internal fermion (vortex line) loop we here assume a negative sign.
### The identification with the Odderon
Even though QCD cannot presently be used to calculate soft processes the typical absence of abrupt changes in experimental distributions indicates that there is no discontinuous transition between soft and hard reactions both formulated on a partonic level in the frame work of the topological model. This provides the hope that hard processes can be used as a guide and that soft processes can be parametrized as an extrapolation of calculable hard processes.
The topological considerations are based on the $`1/N_C`$ \- expansion. Not to loose some of the younger physicist, this approximations selects contributions according to the magnitude of their color factors. To leading order $`1/N_C`$ gluons can be represented by pairs of color lines and the color factors just represent the number of coloring choices. For an amplitude of a given structure with a given number of couplings the leading contribution can be drawn without crossing. An example of a leading and a nonleading contribution is shown below:
The leading term (denoted “a”) contains a new line on the top whose color can be freely chosen. Some amplitudes require special contributions; in this situation the leading terms can be drawn without crossing on topological structures which are more complicated than the simple plane considered above. An example is the cylinder assumed to be responsible for the Pomeron contribution.
The known example of the soft hard correspondence is the connection between soft and hard Pomerons. To identify the hard partner of the soft Pomeron we first observe that the simplest representation of a Pomeron in PQCD involves the exchange of two gluons which can form color singlets with the required positive charge parity. Following this concept it can be shown that a generalization of such an exchange gives the dominant contribution at very high energies in a well defined approximation. It is called “hard” or BFKL Pomeron and involves a ladder of two exchanged Reggeized gluons linked by a number of gluons. In the topological expansion the leading structure of a BFKL Pomeron corresponds to a cylinder with the two basic gluons exchanged on opposite sides parallel to the axis:
Their matching inner color lines can be linked in front of the cylinder without color line crossing. Analogously their matching outer lines can be connected on the back of the cylinder.
Going back to the soft regime the basic assumption in topological models is that the $`1/N`$-expansion stays valid and that the soft Pomeron therefore maintains its cylindrical structure needed for the two string phenomenology of hadronic final states. If cut, soft and hard Pomerons therefore lead to similar two string final states. As difference it remains that the trajectory of the observed soft Pomeron is just shifted downward roughly by a third of a unit from hard Pomeron calculated in leading logarithmic approximation.
Can one find a similar connection for the membraned cylinder? The simplest representation spanning such a topological structure involves three gluons, one on each sheet exchanged parallel to the axis. Any gluon linking these exchanges has then to pass through a vortex line in which the three sheets join. In the $`1/N`$ expansion extended to baryons this means that the color lines have to cross passing this line. The basic structure of the membraned cylinder exchange is therefore the following:
Looking from the other side a color singlet of three gluons can have the quantum numbers of a Pomeron or an Odderon . There is a simple topological property of the Odderon. A single uncrossed gluon connection (of type “a”) would project the color structure of the pair to that of a single gluon and the exchange would have to correspond to a Pomeron-like contribution. The Odderon will therefore have to involve crossed links. Hence it has exactly the topology of the membraned cylinder.
To visualize the baryonic color structure of the Odderon with its crossed exchanges one can replace the exchanged gluons by quark antiquark pairs without changing color lines. The so modified membraned cylinder just represents an exchanged baryon antibaryon pair.
In the same QCD approximation as the “hard” Pomeron the properties of a “hard” or BKP Odderon were calculated and the predicted intercept is $`0.96`$ . Again a mismatch between this hard leading logarithmic Odderon intercept and the experimentally observed soft value (preliminarily $`0.6\pm 0.2`$) by about a third of a unit can be expected the Odderon trajectory.
## 3 Experimental consequences <br>of membraned cylinder exchanges.
Odderon in fits to total hadronic cross sections
The assumed negative contribution from the asymmetric two string cut makes the understanding of total cross sections difficult on a quantitative level. It is possible that the membraned cylinder exchange has a small or almost vanishing imaginary part. In this way there are no (presumably anyhow not serious) constraints from total cross section fits. The cancellation allows a small or vanishing Odderon to contain clearly measurable individual components. In this way data on baryon exchange can be used to determine the
### Odderon in heavy ion scattering <br>
In heavy ion scattering where the Pomerons are dense in transverse space they can join and form a Pomeron or a membraned cylinder. The individual strings are no longer independent but the general picture of particle production in separate universal strings survives. The probability of an interaction of strings and of membraned cylinder exchanges is growing proportional to the density:
$$\frac{[numberofmembranedcylinders]}{[Pomeronnumber]}\frac{[Pomeronnumber][Pomeronradius]^2}{[nucleusradius]^2}$$
(7)
The transition from a Pomeron pair to the centrally cut membraned-cylinder involves baryon antibaryon pair production. Between a proton and a Pomeron the cut membraned-cylinder is a very efficient mechanism of baryon stopping. Both effects correspond to experimental observations. As the trajectory is not well determined it is hard to obtain really reliable quantitative statements which can be tested convincingly with results in heavy ion scattering.
### The backward peak in diffraction and possibly in electro production
There is however a very specific qualitative prediction which can be tested. Consider a diffractive system whose mass exceeds ISR energies. Usually the diffractively produced particles will originate in two strings of a cut Pomeron and the baryon charge will stay on the side of the initial proton. As usual there might be some migration to the center with a slope in rapidity eventually corresponding to the difference of the Odderon and the Pomeron trajectory. Topologically it involves a horizontal cut through the following structure:
The high Odderon trajectory argued for above requires a clear suppression from the coupling constants to stay consistent with data. No such suppression is expected at a two Pomeron vertex. In consequence at a certain distance it should be more favorable for the membraned cylinder to span the total diffractive region and to utilize the more favorable coupling to the two Pomerons. In this way the initial baryon will sometimes end up exactly at the end of the string. It should be visible if one plots the rapidity distribution in relation to the inner end of the diffractive region, i.e. as function of
$$y_{\{Pomeron\}}=y_{\{CMS\}}\mathrm{ln}\frac{m\sqrt{s}}{M(diffr.)}$$
To illustrate the expected small backward peak we show the result of a calculation with the PHOJET Monte Carlo code of the incoming proton spectrum for diffractive events with a mass of $`300`$ GeV for $`pp`$-scattering of $`1.8`$ TeV with standard parameters below. To select diffractive events a lower cutoff of $`x_F=0.95`$ was used.
PHOJET contains diquark exchanges and yields reasonable baryon spectra in the forward region. To obtain the postulated backward peak we just mixed in a suitable sample of inverted events (with disabled diquark exchanges).
## 4 Conclusion
Our aim with this talk and with the more detailed paper is to encourage measurement of the initial baryon distribution in high mass diffractive systems. Similar measurements of a backward peak in electro production might also be possible at HERA.
The prediction is important as it has manifest consequences for heavy ion processes, where it offers a strong mechanism for central baryon production and for the transport of initial baryons to the central and opposite region. It might also clarify the role of the Odderon.
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# THE FUNDAMENTAL PLANE OF FIELD EARLY-TYPE GALAXIES AT INTERMEDIATE REDSHIFT
## 1. Introduction
Understanding the formation and evolution of field early-type galaxies is a cornerstone for the entire picture of galaxy formation. In fact, one of the predictions of hierarchical clustering models is that early-type galaxies and their stellar populations form much later in the field than in the core of rich clusters. This is due to the fact that for a random Gaussian initial density field the collapse of density peaks occurs earlier in the proximity of the large scale overdensities that are bound to be the location of present day clusters (Kauffmann, 1996). A number of studies have shown that the stellar populations of early-type galaxies in the core of rich clusters are homogeneously old (Bower, Lucey & Ellis, 1992; Ellis et al. 1997; Stanford et al. 1998; van Dokkum et al. 1998; Brown et al. 2000). The few results collected so far in the field environment indicate that the star formation history is more complex than what is suggested by the canonical scenario of passive evolution of an old stellar population (see e. g. Schade et al. 1999; Treu & Stiavelli, 1999, and references therein). We report here on a project that we have just completed, aimed at gathering new information on the star formation history of field early-type galaxies by studying the evolution of the Fundamental Plane with redshift as a diagnostic of stellar populations.
The Fundamental Plane (hereafter FP; Djorgovski & Davis 1987; Dressler et al. 1987) of early-type galaxies is defined as
$$\mathrm{log}R_\text{e}=\alpha \mathrm{log}\sigma +\beta \text{SB}_\text{e}+\gamma ,$$
(1)
where $`R_\text{e}`$ is the effective radius in kpc, $`\sigma `$ is the central velocity dispersion in km $`\text{s}^1`$, SB$`_\text{e}`$ is the average surface brightness within the effective radius in $`\text{mag arcsec}^2`$. In the following we refer, when needed, to $`H_0`$=h<sub>50</sub> 50 km $`\text{s}^1`$$`\text{Mpc}^1`$ and to a $`\mathrm{\Lambda }`$ cosmology ($`\mathrm{\Omega }`$=0.3; $`\mathrm{\Omega }_\mathrm{\Lambda }`$=0.7).
A simple physical interpretation of the FP (Faber et al. 1987) can be given by defining an effective mass,
$$M\frac{\sigma ^2R_\text{e}}{G},$$
(2)
suggested by the Virial Theorem, by defining the luminosity in the usual way
$$2.5\mathrm{log}L\text{SB}_\text{e}5\mathrm{log}R_\text{e}2.5\mathrm{log}2\pi ,$$
(3)
and by considering an effective mass-luminosity relation of the form
$$LM^\eta .$$
(4)
These assumptions lead directly to the FP relation given above, provided
$$\alpha 10\beta +2=0,$$
(5)
(??), with $`\eta =0.2\alpha /\beta `$. In this framework, variations of the slopes $`\alpha `$, $`\beta `$ and the intercept $`\gamma `$ (typical values at $`z0`$ are 1.25, 0.32, -8.895, in Johnson B band, so that $`\eta 0.8`$; see Bender et al. 1998) as a function of redshift are easily interpreted as general trends in the luminosity evolution of the galactic stellar population. In particular, if we assume fixed slopes and define for an individual galaxy (labeled by the superscript $`i`$)
$$\gamma ^i\mathrm{log}R_\text{e}^i\alpha \mathrm{log}\sigma ^i\beta \text{SB}_\text{e}^i,$$
(6)
the offset with respect to the prediction of the FP ($`\mathrm{\Delta }\gamma ^i\gamma ^i\gamma `$) is related to the offset of the $`M/L`$ by
$$\mathrm{\Delta }\mathrm{log}\left(\frac{M}{L}\right)^i=\frac{\mathrm{\Delta }\gamma ^i}{2.5\beta }.$$
(7)
Similarly, the scatter (rms) of the FP can be related to the scatter in the $`M/L`$ at a given $`M`$:
$$rms\left(\mathrm{log}\frac{M}{L}\right)=\frac{rms\left(\gamma \right)}{2.5\beta }.$$
(8)
Therefore the very existence of the FP constrains the similarity of the stellar populations of early-type galaxies, and it is interesting to explore how far in the past this similarity extends.
Recently, using the Hubble Space Telescope (HST) for photometry and ground-based spectroscopy from large telescopes, the FP parameters of cluster E/S0s have been measured out to $`z0.8`$ (??; ??, hereafter DF96; ??, hereafter K97; Bender et al. 1998; van Dokkum et al. 1998; Kelson et al. 2000). The data show that the FP is well-defined up to look-back times of $`7.5`$ Gyrs, indicating that the origin of the relation is hidden at still higher redshift. The evolution of the intercept is consistent with what is predicted by the passive evolution of an old stellar population in a $`\mathrm{\Lambda }`$ cosmology ($`z_f2`$; van Dokkum et al. 1998).
In a pilot study of a sample of six early-type galaxies at intermediate redshift (Treu et al. 1999; hereafter T99) we found that the FP is also well defined in the field environment out to $`z0.3`$. We have now collected a larger sample of data. In this paper we give a preliminary report of this study (the complete data-set will be presented in Treu et al. 2000a, the complete analysis in Treu et al. 2000b). In Section 2 we summarize the observational results and describe the location and scatter of the intermediate redshift field FP. In Section 3 we briefly discuss the results in terms of passive evolution of the stellar populations.
## 2. The field FP at intermediate redshift
A sample of early-type galaxies was selected from the Wide Field and Planetary Camera 2 Medium Deep Survey (Griffiths et al. 1994) archival images on the basis of morphology, color, and magnitude (see T00a for a discussion of the sample selection and data analysis). Follow-up spectroscopy was obtained at the ESO-3.6m telescope, yielding accurate velocity dispersions for 22 galaxies in the sample ($`z0.20.4`$).
In Figure 1, panels (a), (b), and (c), we plot the location in the FP space of the intermediate redshift data points binned in redshift. The data are in rest frame V band. The solid line represents the FP of the Coma Cluster (Lucey et al. 1991). In panel (d) we show the evolution of the intercept with redshift, with respect to the Coma relation. Qualitatively, two main facts are evident:
1. At any given redshift between 0 and $`0.4`$ the FP of field early type galaxies is well defined. The scatter is small. No trend of increasing scatter with redshift is noticeable within the accuracy allowed by the small number statistics available (see Figure 2).
2. The intermediate redshift data points at given SB$`_\text{e}`$ and $`\sigma `$ have larger effective radii than the value predicted by the local relation (solid line), and hence are more luminous. The average offset increases with redshift as shown in panel (d).
We can quantify these statements by describing the evolution of the intercept with a simple linear relation $`\mathrm{\Delta }\gamma =\tau z`$ (Figure 2). A least $`\chi ^2`$ fit yields $`\tau =0.54\pm 0.02`$. Taking into account selection effects (see T00b) the 68 % confidence interval is $`0.44<\tau <0.56`$. This description allows us to estimate the scatter of the FP in any redshift bin. In fact, there are three contributions to the scatter: the measurement scatter, the intrinsic scatter, and the evolutionary scatter. The latter is simply the combined effect of the thickness of the redshift bins and the evolution of the intercept. We can correct the measured scatter for this projection effect by “evolving” linearly the galaxies to the average redshift of the bin. The corrected scatter obtained in this way is remarkably small and constant (0.08-0.09 in $`\gamma `$, including measurement scatter $``$0.05-0.06). Given the limited number of data points available per redshift bin the error on the scatter is quite large ($`30`$% for a Gaussian distribution, neglecting systematics); still a substantial increase in the scatter can be ruled out. Larger samples are needed in order to measure the scatter as a function of redshift accurately.
## 3. Evolution of the stellar populations
In this section we compare the observed results on the FP at intermediate redshift with the prediction of stellar population evolution models (Bruzual & Charlot, 1993; GISSEL96 version, hereafter BC96). To this aim we assume that the only evolving factor is the stellar population, that the slopes of the FP do not change with redshift, and that the stellar mass M is proportional to $`M`$. Under these assumptions
$$\mathrm{\Delta }log\frac{M_{}}{L}=\frac{\mathrm{\Delta }\gamma }{2.5\beta },$$
(9)
where $`\mathrm{\Delta }`$ indicates the difference of the quantity with respect to the value found in Coma. In Figure 3 the evolution of the FP is compared with the prediction of single burst stellar population models formed at different formation redshifts ($`z_f`$).
A quantitative comparison is made in Figure 4 where the probability density of $`z_f`$ given the observations is shown. The probability density is obtained with a Bayesan-Montecarlo method described in T00b and takes into account the selection effects. The probability peaks at $`z_f1`$ extending from $`z_f=0.8`$ to $`z_f=1.6`$, corresponding to present ages of 7-10 Gyrs.
The single burst stellar population model is very useful as a benchmark to quantify and compare different results. However, it is clearly just a simplified picture of the star formation history of early-type galaxies. A small amount ($`<`$ 10 %) of the total mass of moderately young stars (one Gyr or so) can alter in a significant way the integrated colors and the $`M_{}/L`$ of an old stellar population. In this scenario, after a few Gyrs, the integrated colors and the $`M_{}/L`$ become totally indistinguishable from the ones of an old single burst stellar population. A wealth of observational evidence suggests that minor episodes of star formation can occur from intermediate redshifts to the present (e. g. Schade et al. 1999; Bernardi et al. 1998), and it is worth addressing this issue for our sample of field early-type galaxies. With an analysis similar to that for the single component case we obtained (T00b) the probability density for the redshift of formation of two stellar components: an older one formed at $`z_{f1}`$ and a younger one, 10 times smaller in mass, formed at $`z_{f2}`$. The contour levels of the probability density corresponding to 68% and 95% probability are shown in Figure 5. Not surprisingly, it is sufficient to have a small mass of stars formed at $`z<0.60.8`$ to make it possible for the rest of the stellar mass to be formed at the very beginning of the Universe ($`z_{f1}`$3-4, i. e. 12-13 Gyrs ago). This scenario is also consistent with the small scatter observed for the FP at low and intermediate redshift, in the sense that the scatter in the FP due to the composite ages of stellar populations is always smaller than the observed one.
## Acknowledgments
This research has been partially funded by the Ministero dell’Università e della Ricerca Scientifica e Tecnologica, by the Space Telescope Science Institute (STScI) Director Discretionary Fund grants 82216 and 82228, by the Agenzia Spaziale Italiana. It is based on observations collected at the European Southern Observatory, La Silla, Chile (Proposals 62.O-0592, 63.O-0468, 64.O-0281) and with the NASA/ESA HST, obtained at the STScI, which is operated by AURA, under NASA contract NAS5-26555.
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# The Relationship between Extremum Statistics and Universal Fluctuations
## Abstract
The normalized probability density function (PDF) of global measures of a large class of highly correlated systems has previously been demonstrated to fall on a single non Gaussian “universal” curve. We derive the functional form of the “global” PDF in terms of the “source” PDF of the individual events in the system. A single parameter distinguishes the global PDF and is related to the exponent of the source PDF. When normalized, the global PDF is shown to be insensitive to this parameter and importantly we obtain the previously demonstrated “universality” from an uncorrelated Gaussian source PDF. The second and third moments of the global PDF are more sensitive, providing a powerful tool to probe the degree of complexity of physical systems.
preprint: HEP/123-qed
The study of systems exhibiting non Gaussian statistics is of considerable current interest. These statistics are observed to arise in finite sized many body systems exhibiting correlation over a broad range of scales. The apparent ubiquitous nature of this behavior has led to interest in self organized criticality as a paradigm; other highly correlated systems include fluid turbulence. Two recent results have highlighted the connection between extremum statistics and highly correlated systems. The probability density function (PDF) of fluctuations in power needed to drive an enclosed rotating turbulent fluid at constant angular frequency has been measured over 2 decades in Reynolds number. Intriguingly, when the PDF $`P(E)`$ of these series of experiments were normalized to the first two moments they were found to fall on a single non Gaussian “universal” curve . This same universal curve was later identified in a study of the two dimensional X-Y model, a numerical model for magnetization near the critical point . To obtain the universal curve, the PDF of a global measure, namely the magnetization summed over the entire system, is again normalized to the first two moments. It was suggested that these two disparate systems share the same statistics as they are both critical. The functional form of the “universal” curve was found for the X-Y model and was shown to be of the form
$`P(E)=K(e^{ye^y})^a\text{with}y=b(Es)`$ (1)
with $`a=\pi /2`$ and $`K,b,s`$ obtained by normalizing the curve to the first two moments. Crucially, it was then demonstrated that this curve was also in reasonable agreement with appropriately chosen normalized global measures for a range of numerical models of highly correlated systems. It was suggested that this behavior is related to the extremum statistics that arises from a process that is highly correlated.
In this Letter we give a comprehensive analysis of extremum statistics in the context of finite sized systems. Our aim is to determine the relationship between the underlying “source” PDF of a given process and the PDF of some global measure. Given that events occur over a range of sizes, and that each event represents some quantity, magnetization, or energy dissipation say, we obtain a relationship between the “source” PDF of the event size, and the PDF of a global measure, the total magnetization, or energy dissipation over the system. We find, as suggested in , that the global PDF, when normalized to the first and second moments is essentially of the form of equation (1). Crucially however we find that the “universal” curve for the global PDF, that is, equation (1) with $`a=\pi /2`$ is not uniquely a property of a source PDF of a correlated process. Instead, in a finite sized system, distributions of this form with $`a`$ in the range $`[1,2]`$ arise from uncorrelated samples from a source PDF ranging from exponential through Gaussian to power law, the value of $`a`$ being determined by the source PDF. When normalized to the first and second moments these curves are only distinguishable asymptotically. Hence in reality the “universal” curve describes, to within typical experimental or numerical statistical uncertainties, distribution (1) with $`a`$ in the range $`[1,2]`$.
In many physical situations it is relatively straightforward to measure the PDF of some global quantity such as power dissipation in the driven turbulent fluid. In order to understand the underlying process we require details of the distribution of the source PDF. In particular, if this process is highly correlated, the source PDF of individual events is anticipated to be power law and we wish to i) distinguish this unambiguously from an uncorrelated Gaussian process and ii) measure the exponent. A direct measurement of the source PDF requires the challenging measurement of event sizes over many decades, but if we can relate the power law exponent to the form of the global PDF there is the possibility to remote sense this exponent. Normalizing the global PDF to the first and second moments is an insensitive method to find $`a`$; we show that for finite sized systems the higher order moments provide a more feasible method.
The first step is to obtain the PDF of some global quantity from that of the source PDF that describes individual events. Consider a finite sized system of dimension $`D`$ which at any instant in time has patches of activity on various length scales up to the system size $`L_B`$. The patches are drawn from the (time independent) source probability $`N(L)`$ of a patch of length $`L`$. These patches can represent sites involved in an avalanche in a sandpile, vortices in a turbulent fluid, ignited trees in a forest fire, or sites with nonzero magnetization in the X- Y model. Associated with the active sites is some quantity of interest, $`Q`$ say, for example energy or magnetization, which we take to be given by $`Q=L^D`$. There will be some maximum $`Q_B`$ corresponding to the (extremely rare) configuration with the highest possible value of $`Q`$, that is, highest energy or magnetization, that can be realized by the system. The total value of $`Q`$ over the system at any instant arises from the distribution of the patches at that instant $`N_j(L)`$;
$`\overline{Q}_j={\displaystyle _0^{Q_B}}QN_j(Q)𝑑Q={\displaystyle _0^{L_B}}L^DN_j(L)𝑑L`$ (2)
where $`N_j`$ is the distribution of an (unknown) realizable ensemble of patches (continuous limit $`N(L)`$) that fits within the finite sized system, and the integral is over the system. Since $`N`$ is normalized, $`N(L)dL=N(Q)dQ`$. We now wish to evaluate the PDF of the $`\overline{Q}_j`$. This arises from the many ensembles of the system, for the $`j^{th}`$ ensemble the total value of $`Q`$ can alternatively be written as a sum over the $`M_j`$ (unknown) individual patches $`\{L_i\}_j`$, $`1iM_j`$. If $`N_j(L)`$ is monotonically decreasing (from maximum $`N_{0j}`$ to zero) we can generate each of the $`\{L_i\}_j`$ by choosing $`M_j`$ random numbers $`N_i`$ in the range $`[0,N_{0j}]`$, with uniform probability distribution $`P(N_i)`$. If we then insist that $`P(N_i)P(L_i)`$, for each realization the random $`N_i`$ will each lie in one of the $`M_j`$ uniform intervals $`\delta N_i`$, giving $`L_i`$ patches which lie in corresponding (nonuniform) intervals $`\delta L_i`$ obtainable in principle by inverting $`N_j(L_i)`$. We can then write the sum of the patches in the $`j^{th}`$ ensemble:
$$\overline{Q}_j=\underset{i=1}{\overset{M_j}{}}\{L_i^D\}_j=\underset{i=1}{\overset{M_j}{}}P(L_i)\delta L_iL_i^D$$
(3)
If the gradient of $`N(L)`$ is near monotonic, $`\delta N_i/\delta L_i(dN/dL)`$ so that
$$\overline{Q}_j\underset{i=1}{\overset{M_j}{}}\frac{P(N_i)\delta N_iL_i^D}{dN/dL}_{N_{0j}}^1\frac{L^DdN}{dN/dL}$$
(4)
For a source PDF $`N(L)`$ that is exponential, Gaussian or inverse power law for large L $`dN/dL<<L^D`$ for small $`N`$, that is, large $`L`$ (large $`Q`$). Hence the dominant contribution to $`\overline{Q}_j`$ is that of the largest patch of activity. Thus the statistics of the PDF of $`\overline{Q}`$, $`P(\overline{Q})`$ will be extremum statistics, $`P(\overline{Q})=P_m(Q)`$, the normalized PDF of the maximum drawn from the ensembles. Given that the maximum for the $`j^{th}`$ ensemble is given by $`Q_j^{}=max\{Q_1,..Q_{M_j}\}`$, where $`Q_{M_j}Q_B`$, that is, $`M_j`$ finite, the PDF for $`Q^{}`$ is given by
$`P_m(Q^{})=MN(Q^{})(1N_>(Q^{}))^{M1}`$ (5)
where $`M`$ is the average of $`M_j`$ over the ensembles and
$`N_>(Q^{})={\displaystyle _Q^{}^{Q_B}}N(Q)𝑑Q{\displaystyle _Q^{}^{\mathrm{}}}N(Q)𝑑Q`$ (6)
We now obtain $`P_m`$ for large finite $`M,Q`$. For a general PDF $`N(Q)`$, $`(1N_>)^M=\mathrm{exp}(Mg(Q^{}))`$ where
$$g(Q^{})=\mathrm{ln}(1N_>(Q^{}))N_>+\frac{N_>^2}{2}$$
(7)
We now choose a characteristic value of $`Q^{}`$, namely $`\stackrel{~}{Q}^{}`$, such that for any of the $`j`$ ensembles
$$q=Mg(\stackrel{~}{Q}^{})=MN_>(\stackrel{~}{Q}^{})+M\frac{N_>^2(\stackrel{~}{Q}^{})}{2}+\mathrm{}$$
(8)
Using this definition and the form for $`g(Q^{})`$ (7) we obtain $`g^{}(\stackrel{~}{Q}^{})=N(\stackrel{~}{Q}^{})`$ to lowest order in an expansion in $`q/M`$.
We now consider specific source PDF $`N(Q)`$. If $`N(Q)`$ falls off sufficiently fast in $`Q`$, i.e. is Gaussian or exponential we can consider lowest order only giving $`g(Q^{})N_>`$ and $`q=MN_>(\stackrel{~}{Q}^{})`$. After some algebra, expanding in $`Q^{}`$ near $`\stackrel{~}{Q}^{}`$ gives
$$P(\overline{Q})=P_m(Q)P_m(Q^{})(e^{ue^u})^a$$
(9)
with
$`a={\displaystyle \frac{N^{}(\stackrel{~}{Q}^{})N_>(\stackrel{~}{Q}^{})}{N^2(\stackrel{~}{Q}^{})}}`$ (10)
$`u=\mathrm{ln}(MN_>(\stackrel{~}{Q^{}}))+{\displaystyle \frac{N(\stackrel{~}{Q}^{})}{N_>(\stackrel{~}{Q}^{})}}\mathrm{\Delta }Q^{}`$ (11)
where $`\mathrm{\Delta }Q=Q\stackrel{~}{Q^{}}`$. For $`N(Q)`$ exponential (11) gives $`a1`$ (see ). For $`N(Q)`$ Gaussian we cannot obtain $`a`$ exactly but as we shall see it is instructive to make an estimate. Given $`N(Q)=N_0\mathrm{exp}(\lambda Q^2)`$ in the above we obtain $`P_m=\overline{P_m}\mathrm{exp}(R(u))`$ with
$`R={\displaystyle \frac{\mathrm{ln}^2(q)}{4\lambda \stackrel{~}{Q}^2}}+\overline{u}\left(1+{\displaystyle \frac{2\mathrm{ln}(q)}{4\lambda \stackrel{~}{Q}^2}}\right){\displaystyle \frac{\overline{u}^2}{4\lambda \stackrel{~}{Q}^2}}e^{\overline{u}}`$ (12)
where we have used $`u=2\lambda \stackrel{~}{Q}^{}\mathrm{\Delta }Q^{}`$ and $`\overline{u}=u+\mathrm{ln}(q)`$. To lowest order in $`\mathrm{\Delta }Q^{}/\stackrel{~}{Q}^{}`$ (i.e. $`\stackrel{~}{Q}^{}\mathrm{}`$) we have PDF (1) with $`a=1`$, but to next order, that is, neglecting the term in $`\overline{u}^2`$ only in (12) we have this PDF with
$`a\left(1+{\displaystyle \frac{2\mathrm{ln}(q)}{4\lambda \stackrel{~}{Q}^2}}\right)1`$ (13)
Power law source PDF $`N(Q)`$ fall off sufficiently slowly with $`Q`$ that we need to go to next order in $`\mathrm{\Delta }Q^{}/\stackrel{~}{Q^{}}`$. If we consider normalizable source PDF
$$N(Q)=\frac{N_0}{(1+Q^2)^k}$$
(14)
then for large $`Q`$ the above method yields that $`P(\overline{Q})`$ is given by the form (9) but with
$$u=\mathrm{ln}(a)\mathrm{ln}(q)(2k1)\frac{\mathrm{\Delta }Q^{}}{\stackrel{~}{Q}^{}}(1\frac{\mathrm{\Delta }Q^{}}{2\stackrel{~}{Q}^{}})$$
(15)
and $`a=2k/2k1`$. To lowest order, neglecting the $`(\mathrm{\Delta }Q^{}/\stackrel{~}{Q}^{})^2`$ term (15) reduces to (11). Hence a power law source PDF has maximal statistics $`P_m(Q)`$ which, when evaluated to next order, have distribution (1) with a correction that is non negligible at the asymptotes, consistent with the well known result due to Frechet ().
The above results should be contrasted with that of Fischer and Tippett . Central to and later derivations is that a single ensemble of $`NM`$ patches has the same statistics as the $`N`$ ensembles (of $`M`$ patches), of which it is comprised. The fixed point of this expression for arbitrarily large $`N`$ and $`M`$ is $`a=1`$ for the exponential and Gaussian PDF, and the Frechet result for power law PDF. Here, we consider a finite sized system so that although the number of realizable ensembles of the system can be taken arbitrarily large, the number of patches $`M`$ per ensemble is always large but finite. Importantly, the rate of convergence with $`M`$ depends on the PDF $`N(L)`$. For an exponential or power law PDF we are able to resum the above expansion exactly to obtain $`a`$; and convergence will then just depend on terms $`O(1/M)`$ and above. This procedure is not possible for $`N(Q)`$ Gaussian, instead we consider the characteristic $`Q^{}`$, that is $`\stackrel{~}{Q}^{}`$ which for $`M`$ arbitrarily large should be large also. Rearranging (8) to lowest order for $`N(Q)=N_0\mathrm{exp}(\lambda Q^2)`$ yields $`\sqrt{\lambda }\stackrel{~}{Q}^{}\sqrt{ln(M)}`$ implying significantly slower convergence.
We now have the intriguing result that for a wide range of source PDF the PDF of a global measure $`P(\overline{Q})`$ is essentially a family of curves that are approximately Gumbel in form and are asymmetric with a handedness that just depends on the sign of $`Q`$; we have assumed $`Q`$ positive whereas one could choose $`Q`$ negative (with $`L`$ positive) in which case $`N(Q)N(Q)`$. The single parameter $`a`$ that distinguishes the global PDF then just depends on the source PDF of the individual events. For $`N(Q)`$ exponential we recover the well known result $`a=1`$. For a power law source PDF $`a`$ is determined by $`k`$ as above. For a Gaussian source PDF $`a1`$.
To compare these curves we normalize $`P(\overline{Q})P_m(Q^{})`$. For Gaussian and exponential source PDF we have
$`\overline{P}(y)=K(e^{ue^u})^a\text{with}u=b(ys)`$ (16)
This has moments
$$M_n=_{\mathrm{}}^{\mathrm{}}y^n\overline{P}(y)𝑑y$$
(17)
which converge for all $`n`$; we insist that $`M_0=1`$, $`M_1=0`$ and $`M_2=1`$. The necessary integrals can be expressed in terms of derivatives of the Gamma function $`\mathrm{\Gamma }(a)`$ and we obtain after some algebra
$`b^2=\mathrm{\Psi }^{}(a),K={\displaystyle \frac{b}{\mathrm{\Gamma }(a)}}e^{a\mathrm{ln}(a)},s={\displaystyle \frac{1}{b}}(\mathrm{\Psi }(a)\mathrm{ln}(a))`$ (18)
where $`\mathrm{\Psi }(a)`$ and its derivative w.r.t. $`a`$, $`\mathrm{\Psi }^{}(a)`$ have their usual meaning. The ambiguity in the sign of $`b`$ (and hence $`s`$) corresponds to the two solutions for $`P(\overline{Q})`$ for positive and negative $`Q`$.
For power law source PDF (14) we use the Frechet distribution () which we first write as (16) with
$$u=\alpha +\beta \mathrm{ln}(1+\frac{y}{G})$$
(19)
which reduces to the form of (16) for $`\mathrm{\Delta }Q^{}/\stackrel{~}{Q}^{}1`$. From (14,15) we identify $`\beta =(2k1)`$. Again we insist that $`M_0=1`$, $`M_1=0`$ and $`M_2=1`$ and obtain
$`\alpha =\beta \mathrm{ln}\left({\displaystyle \frac{a^{\frac{1}{\beta }}}{\mathrm{\Gamma }(1+1/\beta )}}\right)`$ (20)
$`K=\pm \beta a^a\left[\mathrm{\Gamma }(1+2/\beta )\mathrm{\Gamma }^2(1+1/\beta )\right]^{\frac{1}{2}}`$ (21)
$`G={\displaystyle \frac{\mathrm{\Gamma }(1+\frac{1}{\beta })}{\left[\mathrm{\Gamma }(1+\frac{2}{\beta })\mathrm{\Gamma }^2(1+\frac{1}{\beta })\right]^{\frac{1}{2}}}}`$ (22)
For $`\beta \mathrm{}`$ with $`\beta /G`$ finite these equations reduce to (18) with $`b=\beta /G`$.
We can now plot the “universal” curves, that is, normalized to the first two moments. Experimental measurements of a global PDF $`P(E)`$ normalized to $`M_0`$ would be plotted $`M_2P`$ versus $`(EM_1)/M_2`$. For the Frechet it is straightforward to show that the moments of order $`n`$ (17) exist for $`2k>n+1`$ and therefore these curves exist for power law of index $`\mathrm{}>2k>3`$ i.e. $`1<a<3/2`$. This is significant since processes exhibiting long range correlations typically have $`k`$ lower than this . Inset in Figure 1 we plot the normalized Frechet PDF for $`k=3,5,100`$ and the PDF (1) with $`a=1`$. In the limit $`k\mathrm{}`$, $`a1`$ and the normalized Frechet PDF tends to the $`a=1`$ limit of (1), hence for $`k=100`$ these are indistinguishable and differences between the PDF appear on such a plot around the mean for $`k<3`$ approximately. In the main plot we show normalized distributions of the form (1) for $`a=1,\pi /2`$ and $`2`$. It is immediately apparent that the curves are difficult to distinguish for several decades in $`\overline{P}(y)`$ and either numerical or real experiments would require good statistics over a dynamic range of about 4 decades which is not readily achievable.
Since the second moment $`M_2`$ does not exist for $`k3/2`$ we cannot consider curves of $`a3/2`$ generated by power law source PDF; however such values (in particular $`a=\pi /2`$) were identified for the “universal” curves in turbulence experiments and a variety of models of correlated systems . We now demonstrate that these are straightforward to produce. On Figure 1 we have over plotted (\*) the global PDF generated by a source PDF that is uncorrelated Gaussian, calculated numerically. We randomly select $`M`$ uncorrelated variables $`Q_j,j=1,M`$ and to specify the handedness of the extremum distribution, the $`Q_j`$ are defined negative and $`N(Q)`$ is normally distributed. This would physically correspond to a system where the global quantity $`\overline{Q}`$ is negative, i.e. power consumption in a turbulent fluid, as opposed to power generation. To construct the global PDF we generate $`T`$ ensembles, that is select $`T`$ samples of the largest negative number $`Q_i^{}=min\{Q_1..Q_M\},i=1,T`$. For the data shown in the figure $`M=10^5`$ and $`T=10^6`$; this gives $`\sqrt{\lambda }\stackrel{~}{Q}^{}\sqrt{ln(M)}3`$ so that for the Gaussian we are far from the $`a=1`$ limit .
The numerically calculated PDF lies close to $`a=\pi /2`$. Such a value of $`a`$ on these “universal” curves is therefore not strong evidence of a correlated process as suggested by . Generally, plotting data in this way is an insensitive method for determining $`a`$ and thus distinguishing the statistics of the underlying physical process.
The question of interest is whether we can determine the form of the source PDF from the global PDF from data with a reasonable dynamic range. We consider two possibilities here. First, a uniformly sampled process will have the most statistically significant values on the universal curve near the peak. For both the PDF the peak is at $`u=0`$ and is at $`\overline{P}(u=0)=Ke^a`$ with $`K`$ given by (18) and (21) respectively. The latter applies to $`k>3/2`$; for smaller $`k>1`$ we may use $`M_0=1,M_1=0`$ plus a condition on $`\overline{P}(u=0)`$ to obtain $`a`$. A more sensitive indicator may be the third moment of $`\overline{P}`$ which after some algebra can be written as
$`M_3={\displaystyle \frac{\mathrm{\Psi }^{\prime \prime }(a)}{(\mathrm{\Psi }^{}(a))^{\frac{3}{2}}}}`$ (23)
for a Gaussian or exponential source PDF i.e. with (16) and
$`M_3={\displaystyle \frac{\left[\mathrm{\Gamma }(1+\frac{3}{\beta })3\mathrm{\Gamma }(1+\frac{2}{\beta })\mathrm{\Gamma }(1+\frac{1}{\beta })+2\mathrm{\Gamma }^3(1+\frac{1}{\beta })\right]}{\left[\mathrm{\Gamma }(1+\frac{2}{\beta })\mathrm{\Gamma }^2(1+\frac{1}{\beta })\right]^{\frac{3}{2}}}}`$ (24)
for a power law source PDF i.e. with (19); the latter converging for $`2k>4`$. Again these refer to one of the two possible solutions for $`P(\overline{Q})`$; the other solution corresponding to $`yy`$, $`M_3M_3`$. We can compare these two methods by noting that for PDF of the form (1) with $`a=1,2`$ the corresponding values of $`\overline{P}_m`$ differ by $`7.9`$% whereas $`M_3`$ differs by $`32`$%. For Frechet PDF, the variation in $`\overline{P}(u=0)`$ is most significant for smaller $`k`$, for example with $`k=3,4`$ $`\overline{P}(u=0)`$ differs by $`15\%`$ whereas $`M_3`$ differs by $`30\%`$.
In conclusion, we have shown that the statistics of fluctuations in a global measure of a finite sized system, such as total energy dissipation in a turbulent fluid, or total magnetization in a ferromagnet are generally given by extremum statistics. The PDF of the global measure is then one of a family of curves whose moments have been determined in terms of a single parameter $`a`$ which in turn quantifies the PDF of the underlying “source” process, such as the PDF of individual energy release events or patches of magnetization. When normalized to the first and second moments these curves are insensitive to $`a`$ and fall close to the single “universal” curve previously identified as a property of a large class of highly correlated systems , over the range achieved by previous real or numerical experiments. In particular, we find that the global PDF of an uncorrelated Gaussian process is ’Gumbel’ (1) distributed with $`a\pi /2`$, providing a straightforward explanation for the previously demonstrated “universality”. Finally we suggest that the peak, or the third moment of the global PDF is a more sensitive indicator of the source PDF. This is a powerful tool to probe the exponents of physical systems where the source PDF is difficult to measure but provides a signature of the degree of complexity of the system.
###### Acknowledgements.
The authors thank G. King and M. Freeman for illuminating discussions. SCC was supported by PPARC.
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# 1 Bounds on sparticle masses resulting from Eq. (.)
## Acknowledgements
This research is supported in part by PPARC.
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# The Strange Behavior of Critical Branched Polymers
## Abstract
We find that 2-dimensional (2-D) critical branched polymers with no impurities conclusively belong to the same universality class as 2-D random percolation clusters, although pure critical 3-D branched polymers do not belong to the 3-D percolation universality class. We find, moreover, that the fractal dimension $`d_\text{f}`$ of critical branched polymers is independent of the presence of a random environment in 2-D, but not in 3-D. We also report that when there are no impurities the critical branching probability in 3-D is $`b_\text{c}=3.34\times 10^4\pm 0.16\times 10^4`$.
Polymerization has been an important topic of recent interest in the physics community. Polymers are extensively studied complex systems, and can be either linear or branched . In spite of being simpler, linear polymers (LP) present physical features such as non-Markovian growth and nontrivial scaling. For several decades it has been known that linear polymers in diluted solutions can be successfully modeled using self-avoiding walks (SAW). More recently it was found that large SAW chains can be more easily constructed by kinetic growth walks (KGW) that belong to the same universality class as SAW chains . Branched polymers (BP) have a larger degree of complexity and exhibit an astonishingly rich phenomenology . In the 1990’s, the KGW approach was generalized to describe the growth of BP in disordered media. Such BP are typically grown when polymerization occurs in a solution with two types of units: (i) monomers with two chemically active molecular sites (“tips”) that lead to linear polymerization and (ii) three-tipped mononers that lead to branchings (i.e., bifurcations). Many important physical properties of the resulting BP depend on the probability $`b`$ for branching to occur and the concentration $`c`$ of impurities that can block or slow the polymerization.
Here, we first address the controversial question: to which universality class does the critical polymer ($`b=b_\text{c}`$) belong for $`c=0`$? A typical phase diagram (Fig. 1(a)) for BP consists of several distinct regions. For large $`b`$ and small $`c`$, the polymer is compact and grows indefinitely (i.e., is “infinite”) with a smooth (faceted) growth surface, while for smaller values $`1>b>b_\text{c}`$ closer to the critical line $`b_\text{c}(c)`$ the resulting polymer grows indefinitely but with a rough growth surface . For yet smaller values of $`b<b_\text{c}`$ the resulting polymer is finite, due to termination of polymerization at all active growth sites. Along the critical line $`b=b_\text{c}(c)`$, the polymer is fractal (Fig. 1(a)) and has a diverging correlation length.
The commonly held belief that BP belong to the lattice animals universality class was based on the assumption that there is no random environment ($`c=0`$) and that there are only repulsive forces between chains . It was thus a remarkable finding when Bunde et al. showed, using the branched polymer growth model (BPGM), that critical BP in random environments in 2-D belong to the same universality class as 2-D random percolation clusters. One conceivable explanation for this interesting result is that in 2-D such behavior is induced by the random impurity obstacles that, indeed, form real percolation clusters with site occupation probability $`p=c.`$ Hence, it can be argued that the polymer grows as if constrained by the random obstacles, and this effect can generate an effective attraction between chains (Fig. 1(b)). For the special case $`c=0`$, however, there are no such impurities, hence this physical argument becomes irrelevant. The study by Bunde et al. is conclusive all along the critical line $`b=b_\text{c}(c)`$ for $`c0,`$ but the important special case of pure BP ($`c=0`$) still remains unconfirmed , and both the percolation as well as the lattice animals universality classes remain plausible. An additional controversy surrounding the special case $`c=0`$ arises because in many systems the introduction of disorder is sufficient to modify the universality class (and sometimes even to destroy the ordered phase). The Harris criterion predicts that the critical BP in 2-D is a candidate for experiencing a change of universality class when going from $`c=0`$ to $`c0`$.
We investigate this problem by using simulations of the BPGM . BPGM is a generalization of KGW which is able to capture the essential dynamics of branched polymerization. We briefly describe the model here (Fig. 1(b)). The BP is generated from an initiating seed at time $`t=0`$ located at the center of a $`d`$-dimensional lattice. At time step $`t=1`$, one of the vacant nearest-neighbor sites of the seed is chosen randomly and occupied to become the next active growth tip, and the seed ceases to be an active growth tip. At each step $`t`$ of the growth process, a branching also can occur at every such growth tip, by occupation of an additional (second) site with probability $`b`$. A fraction $`c`$ of randomly chosen sites is not available for growth. Hence, at time $`t+1`$, the polymer can grow from any of the active tips added at the previous step $`t`$ to empty neighboring sites (Fig. 1(b)).
Our extensive 2-D simulations, discussed below, leave no doubt that, surprisingly, random percolation is indeed the correct universality class even for $`c=0`$. We simulate critical polymers for $`c=0`$ and $`b=b_\text{c}(0)0.056`$ and find that the fractal dimension $`d_\text{f}`$, the minimum dimension $`d_{\text{min}}`$ (Fig. 2(a)) and the chemical dimension $`d_{\mathrm{}}`$ (Fig. 2(b)) are in strong agreement with the known values for a a 2-D critical percolation cluster . These dimensions are defined by the relations $`Mr^{d_\text{f}}`$, $`M\mathrm{}^d_{\mathrm{}}`$, and $`\mathrm{}r^{d_{\text{min}}}`$, where $`r`$ is radius and $`\mathrm{}`$ is the chemical distance ($`\mathrm{}=t`$ in BPGM). Clearly $`d_\text{f}=d_{\mathrm{}}d_{\text{min}}1.89`$. We conclude that the critical 2-D random percolation cluster and the critical polymer generated by BPGM unmistakably belong to the same universality class. Therefore the fractal dimension of the critical BP in 2-D is independent of the presence of a random environment—in apparent violation of the Harris criterion.
We further investigate this “paradoxical” result by performing similar 3-D simulations of pure ($`c=0`$) critical BP, because the Harris criterion favors a little more the irrelevance of the quenched disorder in 3-D. Until now, such 3-D simulations have not been possible because the precise value of $`b_\text{c}(0)`$ in 3-D was unknown.
So we now determine the value of $`b_\text{c}(0)`$ in 3-D for $`c=0`$. The fundamental difficulty of finding a precise value for $`b_\text{c}(0)`$ in 3-D is that the value is extremely small, $`b_\text{c}(0)<10^3`$ (i.e., on average more than $`10^3`$ monomer units between successive branchings), making it computationally expensive to pinpoint it using traditional numerical approaches that involve tuning $`b`$. We overcome this difficulty by mapping BPGM to the computationally less intensive fixed number of tips model (FNTM) that has been shown to generate phase diagrams identical to BPGM. FNTM differs from BPGM as follows: rather than fixing the branching probability $`b`$, instead FNTM dynamically attempts to fix the number $`N`$ of active polymerization growth tips. Hence, the branching process does not occur with an a priori fixed value $`b`$, but rather branchings occur whenever one or more of the $`N`$ existing active tips “die,” either because of impurities or due to other steric hindrance effects. We can be certain that the value of $`b_\text{c}`$ found using FNTM is identical to the one found from BPGM simulations because FNTM generates polymers in which the number of active growth tips neither vanishes nor explodes exponentially, i.e., corresponding to the critical line at $`b=b_\text{c}`$ in BPGM simulations. FNTM greatly reduces the computational burden of generating critical polymers because it automatically generates critical polymers with $`b=b_\text{c}(c)`$ and, moreover, it has only one free parameter, namely the impurity concentration $`c`$. (The arbitrary value of $`N>1`$ chosen is irrelevant in the large mass limit, as seen in Fig. 3(a)) The effective critical branching probability $`b_{\text{eff}}=B/M`$ is found by dividing the total number $`B`$ of branchings of active tips by the total number $`M`$ of occupied polymer sites. Most importantly, FNTM goes spontaneously to the critical line without the need of parameter tuning (of $`b`$), as with BP growth models that exhibit self-organized criticality (SOC) .
We simulate FNTM for $`c=0`$ in 3-D and $`N=5,10,20,30,50,80`$ and $`100`$ tips as a function of polymer mass $`M`$. In the limit $`M\mathrm{}`$, we find that the values of $`b_{\text{eff}}=B/M`$ converge towards an identical point independently of $`N`$ (Fig. 3(a)). The uncertainty for a typical point in Fig. 3(a) is $`\mathrm{\Delta }b_{\text{eff}}10^5`$. This remarkably small error bar is made possible only because FNTM dynamically seeks the fixed point attractor near $`bb_\text{c}`$, such that the effective critical branching probability $`b_{\text{eff}}=B/M`$ fluctuates around $`b_\text{c}`$. By extrapolating $`b_{\text{eff}}`$ for $`1/M0`$ we find that $`b_\text{c}(0)=3.34\times 10^4\pm 0.16\times 10^4`$.
Since it is not altogether inconceivable that the value of $`b_\text{c}(0)`$ obtained from FNTM will not carry over to BPGM, we check the above conclusions by simulating BPGM with $`2.9\times 10^4<b<4.0\times 10^4`$. We find (Fig. 3(b)) that the number of active tips either grows or decays much faster than a power law except in the approximate range $`3.3\times 10^4<b_\text{c}(0)<3.4\times 10^4`$. This value obtained for $`b_\text{c}(0)`$ is not inconsistent with the value $`b_\text{c}=3.34\times 10^4\pm 0.16\times 10^{04}`$ found using FNTM.
Using this result, we are able to perform 3-D BPGM simulations of pure critical BP. We find that $`d_\text{f}d_{\text{min}}2`$ and $`d_{\mathrm{}}1`$ in 3-D for $`b=b_\text{c}(0)`$ and $`c=0`$ (Figs. 2(c), (d)), definitely ruling out the 3-D random percolation universality class. Although this is the expected field theoretical result (see ref. ) for 3-D polymers, the Harris criterion indicates a relevant disorder and a change in the value of $`d_f`$ (because $`d\nu =1.5<2`$ for $`d=3`$ BP). This prediction is confirmed by the finding of Bunde et al. , that $`d_\text{f}2.53`$ for $`b_\text{c}=1`$, consistent with 3-D random percolation.
Somewhat surprising is our confirmation that in 2-D the fractal dimensions of the critical BP are identical for the pure case $`c=0`$ and the maximally disordered case $`b=1`$, $`c1p_\text{c}`$ (see Fig. 1(a)), where $`p_\text{c}0.59`$ is the critical site occupation probability in random percolation for a square lattice. For $`b=1`$ the critical polymer grows in a vacant space that is itself almost a critical percolation cluster (Fig. 1(a)). Since there is no random environment for $`c=0`$, a larger fractal dimension had been conceivable in principle. Our findings indicate, strangely, that the fractal dimension (but not the lacunarity) of the critical polymer in 2-D is independent of the growth environment. Why this is so is indeed an interesting question that merits further investigation, and seems to be related to a self-organizing process . Criticality occurs only when there is a delicate balance between the rates of tip deaths and branchings—a dynamic not unlike the one used in the self-organized generation of infinite 2-D percolation clusters .
We now comment briefly on our finding that, strangely, pure critical BP fall into the percolation universality class in 2-D but not in 3-D. The deeper reason behind this mysterious behavior is possibly related to the well known fact that in 2-D, encounters between different branches of a BP can result either in “trapping” of a polymer chain or else in “scattering” (bending) of the chains. In 3-D a completely different effect becomes important in the large-scale growth dynamics: interpenetration and “entanglement” of BP chains —not by independent linear chains, but by linear parts of a BP that act independently at criticality, because the chemical distance between successive branchings is much larger than the persistence length. In 2-D, there is no entanglement (for topological reasons), but in 3-D such effects can counteract the excluded volume repulsion, so that the mean square displacement of the tips grows linearly with mass (as with polymer melts, see Ch. 2 and 8 of ref. ), leading to power law scaling equivalent to ideal Gaussian chains with $`d_\text{f}=2`$. Effectively, there is a larger repulsion between chains in 2-D, thereby leading to superdiffusive motion of the tips (with $`d_\text{f}<2)`$. Note that trapping in 3-D is extremely rare, hence the very low value $`b_\text{c}(0)`$.
Finally, one hint for the cause of the observed change in the universality class (when we vary $`b`$ from $`b=b_\text{c}(0)`$ to $`b=1`$) in 3-D becomes evident by comparing the two situations with zero and maximum disorder: for $`b_\text{c}(0)=3.3\times 10^4`$ the “freedom” of motion of the tips is limited only by the sparse steric hindrance due to the polymer chains, while for $`b=1`$ the behavior is dictated by the intersticial percolation cluster available for BP growth. The arguments presented above raise new questions, such as whether other critical exponents vary along the critical line $`b=b_\text{c}(c)`$.
We thank U. L. Fulco, H. J. Hilhorst, M. L. Lyra, S. Roux, and L. R. da Silva for discussions; CNPq, PRONEX, the Projeto Nordeste de Pesquisa & FINEP for support.
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# Coulomb Drag in Coherent Mesoscopic Systems
## Abstract
We present a theory for Coulomb drag between two mesoscopic systems. Our formalism expresses the drag in terms of scattering matrices and wave functions, and its range of validity covers both ballistic and disordered systems. The consequences can be worked out either by analytic means, such as the random matrix theory, or by numerical simulations. We show that Coulomb drag is sensitive to localized states, which usual transport measurements do not probe. For chaotic 2D–systems we find a vanishing average drag, with a nonzero variance. Disordered 1D–wires show a finite drag, with a large variance, giving rise to a possible sign change of the induced current.
Moving charges in a conductor exert a Coulomb force on the charge-carriers in a nearby conductor, thus inducing a drag-current (see Fig. 1). This happens whenever the distance between the two conductors is of the same order as the average distance between charge carriers. In recent years Coulomb drag in two-dimensional systems has been studied extensively and has provided valuable information about the interactions between adjacent extended electron gases.
Coulomb drag of mesoscopic structures has been addressed in the case of 1D–systems both within the Boltzmann equation approach and for Luttinger liquids with strong interwire interactions .
The study of fluctuations in the mesoscopic regime was recently initiated by Narozhny and Aleiner , and it was established that fluctuations will dominate at temperatures smaller than the Thouless energy. This was predicted to be the case even for large extended samples, such as those used in the 2D experiments . While Ref. concentrated on structures larger than the phase-breaking length, $`\mathrm{}_\varphi `$, here we study Coulomb drag of mesoscopic samples smaller than $`\mathrm{}_\varphi `$. Experimentally there is so far only little work on drag in structures with $`L<\mathrm{}_\varphi `$ . We believe this would be an extremely promising new direction for the study of mesoscopic transport properties, since it gives an opportunity to directly study interaction and correlation effects in mesoscopic structures. Especially disordered mesoscopic systems are known to exhibit interesting and unusual physics and the same can be expected for disordered Coulomb drag systems – perhaps even more so because Coulomb drag in addition to the dependence on the transmission properties also has a strong dependence on the nature of the wave function inside the mesoscopic region. We note that Coulomb coupling also has interest in other contexts, such as capacitive coupling of a mesoscopic conductor to the environment, charge pumping in quantum dots, or spin-polarized transport .
Before presenting the technical details we state our main results. We develop a formalism for studying drag in mesoscopic systems, and apply it to a number of special cases. In the case of 1D–wires, studied numerically, we find that even a small amount of disorder induces fluctuations, such that the drag can exceed the ballistic limit, be strongly suppressed, or even change sign. The sign change is a general feature of mesoscopic drag, which we also demonstrate for chaotic systems. Here arguments based on random matrix theory show that the drag is zero on average, while the fluctuations are finite. The zero average drag can thus be taken as a test of the degree of ergodicity of the system under investigation. Furthermore, we address the importance of localized states in the sample . While localized states do not usually effect the ordinary transport properties, they turn out to be important for the transconductance. The reason is that the electron-electron interaction allows for transitions in and out of the localized states, which become visible at temperatures smaller than the level spacing, giving rise to peaks in the transconductance when they cross the Fermi level. We also find a temperature dependence which is very different than the $`T^2`$–dependence found for extended states.
General formulation – Using linear response theory similar to Refs. we find the Coulomb drag to second order in the interaction between mesoscopic subsystems, $`U_{12}`$, taking the isolated systems to be otherwise non-interacting. The general formula for the dc transconductance in the case of two mesoscopic conductors, as illustrated in Fig. 1, is given by
$`G_{21}`$ $`=`$ $`{\displaystyle \frac{e^2}{h}}{\displaystyle d𝐫_1d𝐫_2d𝐫_1^{}d𝐫_2^{}U_{12}(𝐫_1,𝐫_2)U_{12}(𝐫_1^{},𝐫_2^{})}`$ (2)
$`\times \mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega {\displaystyle \frac{\mathrm{\Delta }_1(\omega ,𝐫_1,𝐫_1^{})\mathrm{\Delta }_2(\omega ,𝐫_2,𝐫_2^{})}{2kT\mathrm{sinh}^2(\mathrm{}\omega /2kT)}},`$
where $`\mathrm{\Delta }`$ is the three point correlation function $`\widehat{I}\widehat{\rho }\widehat{\rho }`$, as explained in Ref. . Eq. (1) generalizes the results of Ref. to systems with broken translation-invariance. For the case of mesoscopic conductors it becomes
$`\mathrm{\Delta }_i(\omega ,𝐫,𝐫^{})=2i\pi ^2\mathrm{}{\displaystyle \underset{\beta }{}}\theta _\beta ^i(𝐫,𝐫^{},\epsilon _\beta \mathrm{}\omega )`$ (3)
$`\times [n_F(\epsilon _\beta \mathrm{}\omega )n_F(\epsilon _\beta )]+(𝐫𝐫^{};\omega \omega ).`$ (4)
Here
$$\theta _\beta ^i(𝐫,𝐫^{},\epsilon )=\underset{\alpha \gamma }{}I_{\alpha \gamma }^i\rho _{\alpha \beta }^i(𝐫)\rho _{\beta \gamma }^i(𝐫^{})\delta (\xi _\alpha )\delta (\xi _\gamma ),$$
(5)
where $`\xi _\alpha =\epsilon _\alpha \epsilon `$ and $`i`$ labels the subsystem. The matrix elements are given by $`I_{\alpha \gamma }^i=\alpha |\widehat{I}^i|\gamma `$ and $`\rho _{\alpha \beta }^i(𝐫)=\alpha |𝐫𝐫|\beta `$, where $`|\alpha >`$’s are the eigenstates of the uncoupled subsystem with energies $`\epsilon _\alpha `$. Using scattering states as the basis we get $`I_{\alpha \beta }^i=\frac{\mathrm{}}{2m}\delta _{\epsilon _\alpha ,\epsilon _\beta }j_{\alpha \beta }`$, where the matrix $`j`$ can be expressed in terms of the $`2N\times 2N`$ scattering matrix $`S`$ as $`j=\left(\tau ^3S^{}\tau ^3S\right)`$. Here, $`\tau _{nn^{}}^3=\pm \delta _{nn^{}}`$ with plus for $`n`$ belonging to right moving scattering states and minus for the left moving states.
Some general features immediately follow from Eq. (2). The usual cancellation of velocity and density of states, which is central in the derivation of the Landauer–Büttiker formula, occurs only for $`I_{\alpha \gamma }^i`$, whereas for $`\rho _{\alpha \beta }^i`$ this is not the case. Consequently, in contrast to individual subsystem conductances $`G_{ii}`$, $`G_{21}`$ peaks at the onset of new modes in either of the subsystems. Secondly, we notice that the sum over $`|\beta `$ mixes both propagating and evanescent modes. This means that apart from the transmission properties also localized states are probed by measuring drag conductance. Finally, we notice that the outcome of Eq. (2) can have any sign, which is directly related to lack of translation-invariance.
The low temperature limit also follows readily from Eq. (2). The factor $`\mathrm{sinh}^2`$ cuts off the frequency integration and we can expand the $`\mathrm{\Delta }`$’s to lowest order in $`\omega `$. This gives $`\mathrm{\Delta }\omega `$ with the sum over states restricted to those at the Fermi level ($`\xi _\beta ^F=\epsilon _\beta \epsilon _F`$):
$$\mathrm{\Delta }_i(\omega ,𝐫,𝐫^{})=4\omega \pi ^2\mathrm{}^2\mathrm{I}m\underset{\beta }{}\theta _\beta (𝐫,𝐫^{},\epsilon _F)\delta (\xi _\beta ^F).$$
(6)
We immediately see that the transconductance in this limit becomes proportional to $`T^2`$, in accordance with the usual Fermi liquid result for electron-electron scattering. Note however that the low temperature expansion breaks down when the temperature becomes smaller than the level spacing of the discrete, i.e. localized states, which we discuss in detail below. At higher temperatures the $`T^2`$-behavior is replaced by a weaker temperature dependence (e.g. for a quasi 1D–system, $`G_{21}T`$ for $`kT>\mathrm{}v_F/L`$ as considered in Ref. ). Here we however concentrate on the low temperature dependence.
One-dimensional wires – Next we consider as an illustrative example two disordered 1D–wires, which we solve both numerically and analytically using perturbation theory. The one-dimensional case shows that a small amount of disorder can lead to large fluctuations for the drag response and even reverse the sign. The reason for this is that inter-wire interaction induced forward scattering gives rise to a drag response provided it is combined with disorder induced backscattering. In contrast, in the case of clean wires the backscattering is induced solely by the interwire interaction, and therefore the disordered case is larger by a factor of order $`<>U_{12}(0)/U_{12}(2k_F)`$, with $`U_{12}(q)=_0^L_0^L𝑑x_1𝑑x_2e^{iq(x_1x_2)}U_{12}(x_1,x_2)`$ being the Fourier transformed interaction and $``$ is the reflection coefficient, which is inversely proportional to the mean free path $`<>L/\mathrm{}`$. We can show this explicitly by considering the lowest order perturbation theory in disorder potential, corresponding to the diagrams shown in Fig. 2, and for long wires $`k_FL1`$ we find
$$\frac{<[\delta G_{21}(\mathrm{})]^2>^{1/2}}{G_{21}(\mathrm{})}\frac{\left[2<_1><_2>U_{12}^2(2k_F)\stackrel{~}{U}_{12}^2(0)\right]^{1/2}}{U_{12}^2(2k_F)},$$
(7)
where
$`\stackrel{~}{U}_{12}^2(0){\displaystyle _0^L}{\displaystyle _0^L}{\displaystyle _0^L}{\displaystyle _0^L}𝑑x_1𝑑x_2𝑑x_1^{}𝑑x_2^{}`$ (8)
$`\times U_{12}(x_1,x_2)U_{12}(x_1^{},x_2^{})\left(1\frac{2|x_1x_1^{}|}{L}\right)\left(1\frac{2|x_2x_2^{}|}{L}\right).`$ (9)
The denominator is the result $`G_{21}(\mathrm{})U_{12}^2(2k_F)`$ for ballistic wires. For the realistic case where $`U_{12}(2k_F)\stackrel{~}{U}_{12}(0)`$ we see that the fluctuations of the drag can exceed the average value. This is in contrast to the fluctuations of the diagonal conductance $`<[\delta G_{ii}]^2>^{1/2}`$, which are vanishing compared to the mean value $`<G_{ii}>=(2e^2/h)\left(1<_i>\right)2e^2/h`$ in the limit of weak disorder. Fig. 2 displays the prediction of Eq. (7) along with the numerical results described below and very good agreement is seen.
In order to solve the 1D model numerically, we study Eq. (4) on a lattice using the method of finite differences . The method offers a way of studying disordered systems by ensemble averaging over different disorder-configurations . In our numerical example, we use a bare long-ranged Coulomb interaction and the Anderson model with diagonal disorder . We have numerically studied the drag as a function of the mean free path $`\mathrm{}`$ and the length $`L`$, choosing the Fermi energy corresponding to a quarter-filled band, and for a separation given by $`k_Fd=1`$. We calculate both $`G_{11}`$, $`G_{22}`$, and $`G_{21}`$. Since the potentials in the two wires are uncorrelated we in general have $`G_{11}G_{22}`$, but $`<G_{11}><G_{22}>`$ and $`<(\delta G_{11})^2><(\delta G_{22})^2>`$. Our numerical results for distributions, mean values, and fluctuations for $`G_{ii}`$ are in full agreement with the results of Abrikosov . In the delocalized regime $`\mathrm{}L`$ we find as expected that disorder has almost no effect on $`G_{ii}`$ and $`<G_{ii}>2e^2/h`$ with very small fluctuations. Fig. 2 shows $`<[\delta G_{21}(\mathrm{})]^2>^{1/2}`$ normalized by the drag $`G_{21}(\mathrm{})`$ in the ballistic regime as a function of $`k_F\mathrm{}`$. The expected $`1/\mathrm{}`$-dependence is born out by the numerical calculations and we also find that the fluctuations increase with the length of the wires. The inset shows a typical histogram of the drag conductance showing that depending on the disorder configuration $`G_{21}(\mathrm{})`$ can be either higher or lower than in the ballistic regime. Furthermore, note that in agreement with the arguments given above the drag conductance shows a sign reversal for some disorder realizations.
Localized states – The low temperature expansion Eq. (6), which results in a $`T^2`$-dependence, is only valid if $`|\beta `$ belongs to a continuum of states. To investigate the effects due to localized states we split $`\mathrm{\Delta }`$ in two parts, $`\mathrm{\Delta }=\mathrm{\Delta }_d+\mathrm{\Delta }_l`$, where the first term is given by Eq. (6) while the second term is due to scattering in and out of localized states,
$`\mathrm{\Delta }_l(x,y;\omega )=2i\pi ^2\mathrm{}{\displaystyle \underset{\beta \text{localized}}{}}x|\beta \beta |y`$ (10)
$`\times [\varphi (x,y;\epsilon _\beta \mathrm{}\omega )(n_F(\epsilon _\beta \mathrm{}\omega )n_F(\epsilon _\beta ))`$ (11)
$`+\varphi (y,x;\epsilon _\beta +\mathrm{}\omega )(n_F(\epsilon _\beta +\mathrm{}\omega )n_F(\epsilon _\beta ))],`$ (12)
where the localized states have been chosen to be real functions, and
$$\varphi (x,y;\epsilon )=\underset{\alpha \gamma }{}\delta \left(\xi _\alpha \right)\delta \left(\xi _\gamma \right)\gamma |\widehat{I}|\alpha \alpha |xy|\gamma .$$
(13)
At low temperatures we can approximate $`\varphi (x,y;\epsilon \pm \mathrm{}\omega )\varphi (x,y;\epsilon _F)`$, which allows the temperature dependence to be extracted by integration over $`\omega `$ in Eq. (2). Furthermore, for temperatures less than the level spacing the response will be dominated by the coupling to the localized level lying closest to the Fermi level. There are thus three different types of contributions corresponding to the response due to localized/delocalized states in each subsystem, $`G_{21}=\frac{e^2}{h}\left(g_{dd}+g_{ld}+g_{ll}\right)`$ where $`g_{dd}T^2`$. Let us consider, say, $`g_{ld}`$ in some detail. We find
$`g_{ld}`$ $``$ $`\mathrm{}{\displaystyle d\omega \frac{\mathrm{}\omega \left[n_F(\epsilon _1+\mathrm{}\omega )n_F(\epsilon _1)\right]}{kT\mathrm{sinh}^2(\mathrm{}\omega /2kT)}}`$ (14)
$``$ $`{\displaystyle \frac{5kT}{\mathrm{cosh}[0.57(\epsilon _1\epsilon _\mathrm{F})/kT]}},`$ (15)
where $`\epsilon _1`$ is the energy of the localized level lying closest to the Fermi energy. A similar calculation gives
$$g_{ll}\frac{1}{\mathrm{cosh}[(\epsilon _1\epsilon _F)/2kT]\mathrm{cosh}[(\epsilon _2\epsilon _F)/2kT]}.$$
(16)
The relative strengths of these terms can be estimated as
$$\frac{g_{dd}}{g_{ll}}\left(\frac{kT}{\epsilon _F}k_F\sqrt{𝒜}\right)^2,\frac{g_{dd}}{g_{dl}}\left(\frac{kT}{\epsilon _F}k_F\sqrt{𝒜}\right),$$
(17)
with $`𝒜`$ being the interaction area. Thus at low temperature the contributions due to localized states will dominate. The temperature dependence is very different from the usual $`T^2`$ law, and it may even be temperature independent if both $`\epsilon _1`$ and $`\epsilon _2`$ lie on the Fermi level. By adjusting the Fermi energy or system parameters, one can use the drag response to probe the properties and statistics of localized states.
Random matrix theory – We now discuss the statistical properties of the transconductance. This is important in order to determine the size of the Coulomb drag for an ensemble of disordered mesoscopic systems, such as suggested in Fig. 1. Our starting point is the low temperature result (6) (neglecting localized states). For the calculation we need the statistical properties of the $`S`$-matrix, the eigenstates, and the eigenvalues. We assume that the region where the subsystems couple by Coulomb interactions are disordered and that they can be described by random matrix theory . This means that the eigenvalues and the wave functions are assumed to be uncorrelated and furthermore that the current matrix elements $`I_{\alpha \beta }`$ are uncorrelated with the value of wave functions. The latter follows from the fact that the current matrix elements are independent of position and may be evaluated outside the disordered region, and hence do not correlate with the wave functions inside the disordered region. With these approximations
$$\frac{\mathrm{\Delta }(\omega ,𝐫,𝐫^{})}{4\omega \pi ^2\mathrm{}^2}\mathrm{I}m\underset{\alpha \beta \gamma }{}I_{\alpha \gamma }\rho _{\alpha \beta }\rho _{\beta \gamma }^{}\delta (\xi _\alpha ^F)\delta (\xi _\beta ^F)\delta (\xi _\gamma ^F).$$
(18)
The average of the current matrix element is evaluated using standard RMT , and both with and without time reversal symmetry we find $`I_{\alpha \gamma }=(\mathrm{}/2m)(\tau ^3+S^{}\tau ^3S)_{\alpha \gamma }\tau _{\alpha \gamma }^3`$, and since the second average in $`\mathrm{\Delta }`$ is symmetric with respect to interchange of $`\alpha `$ and $`\gamma `$ we get $`\mathrm{\Delta }=0`$ and of course therefore $`G_{21}`$=0. The fluctuations are, however, nonzero and involve the average $`\mathrm{\Delta }(\omega ,𝐫,𝐫^{})\mathrm{\Delta }(\stackrel{~}{\omega },𝐬,𝐬^{})`$ and hence the combination $`(S^{}\tau ^3S)_{\alpha \beta }(S^{}\tau ^3S)_{\alpha ^{}\beta ^{}}`$, which in the limit of a large $`N`$ becomes $`(2N)^2\delta _{\alpha \beta ^{}}\delta _{\alpha ^{}\beta }`$. Interestingly, again the result is not changed by breaking of time reversal symmetry, in contrast to the UCF case, where the results with or without an applied $`B`$-field differ by a factor of 2 . The variance of the $`\mathrm{\Delta }`$ then reads
$$\frac{\mathrm{\Delta }(\omega ,𝐫,𝐫^{})\mathrm{\Delta }(\stackrel{~}{\omega },𝐬,𝐬^{})}{\pi ^2\omega \stackrel{~}{\omega }}\frac{C(𝐫,𝐫^{},𝐬^{},𝐬)C(𝐫,𝐫^{},𝐬,𝐬^{})}{(2N)^2},$$
(19)
where $`C`$ is a correlation function involving four density matrices
$`C(𝐫,𝐫^{},𝐬,𝐬^{})={\displaystyle \underset{\alpha \alpha ^{}\beta \beta ^{}}{}}\rho _{\alpha \beta }(𝐫)\rho _{\beta \alpha ^{}}(𝐫^{})\rho _{\alpha ^{}\beta ^{}}(𝐬)\rho _{\beta ^{}\beta }(𝐬^{})`$ (20)
$`\times \delta (\xi _\alpha ^F)\delta (\xi _\beta ^F)\delta (\xi _\alpha ^{}^F)\delta (\xi _\beta ^{}^F)`$ (21)
$`{\displaystyle \frac{1}{(2\pi )^4}}A(𝐫,𝐫^{})A(𝐬,𝐬^{})A(𝐫,𝐬^{})A(𝐫^{},𝐬)`$ (22)
to lowest order in $`1/k_F\mathrm{}`$. Using the average spectral function relevant to the 2D case $`A(r)(m/2\mathrm{}^2)\mathrm{exp}(r/2\mathrm{})J_0(k_Fr)`$, and assuming, in addition to $`k_F\mathrm{}1`$ also that $`\mathrm{}r_s`$, where $`r_s`$ is the screening length, we obtain the estimate
$$\delta G_{21}^2^{1/2}10^4\frac{e^2}{h}\left(\frac{kT}{\epsilon _F}\frac{U_{12}(d)}{\epsilon _F}\right)^2\frac{r_s^2k_F\sqrt{𝒜}}{\mathrm{}^2N^2}.$$
(23)
With typical numbers for GaAs 2DEG structures the fluctuations of the transresistance are of the order of $`0.1`$ Ohm, which should be measurable.
In conclusion, we have studied drag in mesoscopic systems and argue that measurement of transconductance provides an interesting new method for investigation of the electronic properties of these systems.
We thank C. W. J. Beenakker and M. Brandbyge for useful discussions.
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# Yukawa Unification and Unstable Minima of the Supersymmetric Scalar Potential
## 1 Introduction
If we have to go beyond the Standard Model (SM), for which there are ample motivations, the most popular choice seems to be supersymmetry (SUSY) . With a plethora of new degrees of freedom, it is necessary to constrain them, in addition to direct searches at the colliders, in as many ways as possible so that the parameter space for the SUSY particles may be narrowed down. One of the most useful ways to put such constraints is to consider the dangerous directions of the scalar potential where the potential may be unbounded from below (UFB) or develops a charge and/or color breaking (CCB) minima . This may happen since one now has charged and colored scalar fields in the spectrum, and the possible existence of such a direction would make the standard vacuum unstable. Different directions are chosen by giving vacuum expectation value (VEV) to one or more scalar fields, while keeping the VEVs of the other scalars to zero.
Such constraints, in fact, are very powerful. This may be realized from the fact that the allowed parameter space (APS) for SUSY models is practically unrestricted as one goes for larger and larger values of the soft SUSY breaking parameters (e.g., the universal scalar and gaugino masses, the trilinear coupling, etc.)<sup>1</sup><sup>1</sup>1 Apart from the fact that they should not be more than a few TeV if we have to have an acceptable solution to the hierarchy problem, there is no hint from the theory about their actual values. beyond the kinematic limit of the current high energy colliders. On the otherhand the UFB and CCB constraints quite often acquire greater eliminating power to rule out significant parts of such regions beyond the striking range of the current experiments. Thus, there is an intricate balance between such ‘potential constraints’ and the expanding SUSY APS. For some values of the free parameters, the UFB and CCB conditions are very sharp and disallow most of the parameter space that is otherwise allowed; for some other values, they lose their constraining power.
In a very interesting paper which revived interest in UFB and CCB constraints, Casas et al investigated such constraints on SUSY models. Though their formulae are fairly model-independent, they have carried out the numerical analysis for moderate values of $`\mathrm{tan}\beta `$ (the ratio of the vacuum expectation values (VEV) of the two Higgs fields) only, when one can ignore the effects of b and $`\tau `$ Yukawa couplings in the relevant renormalization group equations (RGE’s). Further they have used the standard minimal supergravity (MSUGRA) assumption of universal soft scalar mass $`m_0`$ and universal gaugino mass $`m_{1/2}`$ at the GUT scale $`M_G`$, referred to hereafter as the ‘conventional scenario’, to determine the sparticle spectrum. Their main result was that within the framework of MSUGRA, a certain UFB constraint known as UFB-3 with VEV given in the direction of the slepton field puts the tightest bound on the SUSY parameter space that they considered (see eq. (93) of and the discussions that follows).
The purpose of this work is to extend and complement the work of by analyzing the effectiveness of the UFB constraints for large $`\mathrm{tan}\beta `$, This we have done in two models: (i) the conventional scenario and (ii) a modified version of MSUGRA within the frame of a SO(10) GUT where the sfermion soft masses $`m_{16}`$ are universal at the GUT scale, but the Higgs soft masses $`m_{10}`$ are diferent from them (this we will call the ‘nonuniversal scenario’). In course of this work we have realized that in contrast to the low $`\mathrm{tan}\beta `$ scenario, the UFB-3 constraint with squarks (eq. 31 of ) may become stronger under certain circumstances, and over a large part of the parameter space the constraint known as UFB-1 (see eq. (8)) serves as the chief restrictor of the APS.
It is well known that there are quite a few motivations for going beyond small and intermediate values of $`\mathrm{tan}\beta `$ in the context of Grand Unified Theories (GUT). If one assumes the GUT group SO(10) breaking directly to the SM gauge group SU(3)$`\times `$SU(2)$`\times `$U(1), and a minimal Higgs field content (only one 10 containing both the light Higgs doublets required in MSSM), the top, bottom and $`\tau `$ Yukawa couplings must unify to a definite GUT scale value at the scale where SO(10) breaks . Within the framework of GUTs partial $`b\tau `$ Yukawa unification is also an attractive possibility . In an SO(10) model, even if one assumes more than one 10-plet of Higgs fields, $`\tau `$ and bottom Yukawa couplings must unify, but the top Yukawa may not unify with them at the GUT scale $`M_G`$.
It can be shown that $`\mathrm{tan}\beta `$ must lie in the range 45-52 for $`tb\tau `$ Yukawa unification (for $`m_t=175`$ GeV) and in the range 30-50 for only $`b\tau `$ unification. We do not consider the possibility $`\mathrm{tan}\beta 2`$ since such low values of $`\mathrm{tan}\beta `$ are now under pressure due to the lower bound on the lightest Higgs boson mass from LEP . This justifies the enthusiasm that has been generated regarding the phenomenology of large $`\mathrm{tan}\beta `$ scenario .
To motivate the nonuniversal scenario under consideration, let us note that from a SUGRA point of view, it is natural to choose the scale at which SUSY breaks in the vicinity of the Planck scale $`M_P2.4\times 10^{18}`$ GeV. At this scale, one may have truly universal soft masses for all scalars; however, the running of the scalar masses between $`M_P`$ and $`M_G`$ can lead to a nondegeneracy at $`M_G`$ . Within the framework of an SO(10) GUT, the first two sfermion generations will still be degenerate, as they live in the same representation of SO(10) and have negligible Yukawa couplings. The Higgs fields live in a different representation, and couple to other heavy GUT fields to generate the doublet-triplet splitting; so their masses can change significantly. The third generation sfermions may have a large Yukawa coupling and hence may be nondegenerate from the first two generation of sfermions, though this effect has not been taken into account in our discussion for simplicity. Only the Higgs mass parameter ($`m_{10}`$) at $`M_G`$ is assumed to be different from the common soft sfermion mass ($`m_{16}`$) at that scale, and both are treated as free parameters.
In addition to restricting the values of $`\mathrm{tan}\beta `$, the requirement of Yukawa unification eliminates a significant region of the otherwise large APS of MSSM quite effectively. For example, this unification occurs within a rather limited region of the $`m_{16}m_{1/2}`$ plane for certain generic choices of the common trilinear soft breaking parameter $`A_0`$. This dependence arises largely through the radiative corrections to the running bottom quark mass which in turn controls the bottom quark Yukawa coupling $`\lambda _b`$ at low energies. The UFB-1 and the UFB-3 conditions further eliminate a significant region from this already restricted APS, which is one of the main results of this paper. Throughout the paper we ignore the possibility that the nonrenormalizable effective operators may stabilise the potential .
The APS obtained by requiring Yukawa unification only is quite sensitive on the choice of $`A_0`$. For example, in the conventional scenario with $`b\tau `$ unification, the APS increases quite a bit for large negative values of $`A_0`$. It is precisely these values of $`A_0`$ which makes the potential more vulnerable to the UFB conditions and many of the additional points allowed by choosing $`A_0`$ appropriately are eliminated by the UFB conditions, as will be demonstrated in a later section. Thus, there is a nice complementary behaviour: for large negative $`A_0`$, the Yukawa unification criterion is a weak condition but UFB conditions are very strong, while for small negative values of $`A_0`$ the roles are reversed. For positive $`A_0`$, none of these criteria are sufficiently strong.
Following the same procedure, significant regions of the parameter space can be eliminated for models with $`tb\tau `$ Yukawa unification. In particular, we find that $`A_00`$ is completely ruled out.
The effectiveness of Yukawa unification as a restrictor of the APS also diminishes, as expected, as the accuracy with which we require the unification to hold good is relaxed. There are several reasons why the unification may not be exact. First, there may be threshold corrections , both at the SUSY breaking scale (due to nondegeneracy of the sparticles) and at $`M_G`$, of which no exact estimates exist. Secondly, we have used two-loop RGE’s for gauge couplings as well as Yukawa couplings and one loop RGEs for the soft breaking parameters, but higher order loop corrections may be important at a few percent level at higher energy scales. Finally the success of the unification program is also dependent on the choice of $`\alpha _s(M_Z)`$ which is not known as precisely as $`\alpha _1`$ or $`\alpha _2`$. To circumvent such drawbacks, one relaxes the Yukawa unification condition to a finite amount (5%, 10% or 20%) which should indirectly take care of these possible caveats. It is interesing to note that quite often the UFB constraints rule out subtantial parts of the extended APS.
Some of the “potential” constraints analyzed here were also discussed by Rattazzi and Sarid . However, they considered the RG improved tree-level potentials only and included the possibility of stabilizing the potential through nonrenormalizable effective operators. Moreover the potent UFB-3 constraint was not available at the time of their analysis. Finally a systematic analysis of the APS in the $`m_{1/2}`$-$`m_{16}`$ plane, which is very relevant for physics studies at the Large Hadronic Collider (LHC), was not presented.
When the SO(10) symmetry breaks down to the SM symmetry, there may be a nonzero D-term, which causes the mass splitting between sfermions in 5 and $`\overline{\mathrm{𝟏𝟎}}`$ of SO(10) . Recently, a number of authors addressed to the phenomenology of the SO(10) D-terms . In this paper, we take the D-term to be zero for simplicity; with a nonzero D-term, one gets a wider variety if constraints which will be discussed in a subsequent paper .
It is well-known that there is a basic conflict between $`bs\gamma `$ and $`tb\tau `$ Yukawa unification. The latter works best for $`\mu <0`$ and large $`\mathrm{tan}\beta `$, while at the same time this region of the parameter space tends to give unacceptable contributions to the former . However, in view of the uncertainties in the long-distance corrections and the possibility of cancellation between various diagrams, we have not included this constraint in our analysis.
The plan of the paper is as follows. In the next section, we outline the various UFB directions of the supersymmetric potential, and discuss our methodolgy. The next section deals with the results and in the last section, we summarize and conclude.
## 2 UFB Directions of the SUSY potential
In this section we briefly review the necessary formulae for the UFB directions following and . We closely follow the former reference in defining the said directions.
The scalar potential of the MSSM is a function of several scalar fields. An SU(2)$`\times `$U(1) breaking minimum of this potential must exist for preserving the phenomenological successes of the SM. Moreover, one demands this real minimum $`V_{realmin}`$ to be deeper than the unwanted UFB and CCB minima. These minima are computed by giving VEV to one or more scalar components at a time; the condition is that at no point in such dangerous directions the potential should be deeper than $`V_{realmin}`$. The resulting constraints on the field space are of much importance as they can restrict the soft SUSY breaking parameters, and hence the sparticle masses and couplings . Let us see how these dangerous field directions arise.
The tree level scalar potential in the MSSM can be written as the sum of the D-term, the F-term and the soft mass term:
$$V_0=V_F+V_D+V_{soft}$$
(1)
where
$`V_F`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}\left|{\displaystyle \frac{W}{\varphi _\alpha }}\right|^2,`$
$`V_D`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{a}{}}g_a^2\left({\displaystyle \underset{\alpha }{}}\varphi _\alpha ^{}T^a\varphi _\alpha \right)^2,`$
$`V_{soft}`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}m_{\varphi _\alpha }^2|\varphi _\alpha |^2+{\displaystyle \underset{i}{}}\{A_{u_i}\lambda _{u_i}Q_iH_2u_i+A_{d_i}\lambda _{d_i}Q_iH_1d_i+A_{e_i}\lambda _{ei}L_iH_1e_i+h.c.\}`$ (2)
$`+(B\mu H_1H_2+h.c)`$
and the superpotential $`W`$ is defined as
$$W=\underset{i}{}\{\lambda _{u_i}Q_iH_2u_i+\lambda _{d_i}Q_iH_1d_i+\lambda _{ei}L_iH_1e_i\}+\mu H_1H_2.$$
(3)
Here, $`\varphi _\alpha `$ are the generic scalar fields, $`T_a`$ and $`g_a`$ are the gauge group generators and the gauge couplings respectively, and $`\lambda `$’s are the respective Yukawa couplings. $`A_i`$, $`B`$ and $`m_i`$ are the soft SUSY breaking parameters, and $`\mu `$ is the Higgsino mass term. In eq. (3) $`Q_i`$ and $`L_i`$ stand for SU(2) doublet squark and slepton superfields while $`u,d`$ and $`e`$ are the corresponding singlet superfields. The generation index $`i`$ runs from 1 to 3.
The neutral part of the Higgs potential in the MSSM is given by
$$V_{Higgs}=m_1^2\left|H_1^2\right|+m_2^2\left|H_2^2\right|2\left|m_3^2\right|\left|H_1\right|\left|H_2\right|+\frac{1}{8}(g^2+g_2^2)(\left|H_2\right|^2\left|H_1\right|^2)^2$$
(4)
where $`m_1^2=m_{H_1}^2+\mu ^2`$, $`m_2^2=m_{H_2}^2+\mu ^2`$, and $`m_3^2=\mu B`$, with $`m_{H_1}`$ and $`m_{H_2}`$ being the mass terms of the two doublets. The standard GUT normalization is used for the gauge couplings: $`g_3=g_2=g_1=\sqrt{5/3}g^{}`$ at $`M_G`$. Minimization of this tree-level potential yields
$$V_{realmin}=\frac{\{[(m_{1}^{}{}_{}{}^{2}+m_{2}^{}{}_{}{}^{2})^24\left|m_3\right|^4]^{1/2}m_{1}^{}{}_{}{}^{2}+m_{2}^{}{}_{}{}^{2}\}^2}{2(g_{}^{}{}_{}{}^{2}+g_{2}^{}{}_{}{}^{2})}.$$
(5)
At any scale $`Q`$, there is a significant radiative correction to this potential. Including the one-loop corrections, the potential becomes
$$V=V_0+\mathrm{\Delta }V_1$$
(6)
where
$$\mathrm{\Delta }V_1=\underset{\alpha }{}\frac{n_\alpha }{64\pi }M_\alpha ^4\left[\mathrm{ln}\frac{M_\alpha ^2}{Q^2}\frac{3}{2}\right],$$
(7)
with $`n_\alpha =(1)^{2s_\alpha }(2s_\alpha +1)`$, $`s_\alpha `$ being the spin of the corresponding field. One ensures the minima of $`H_1`$ and $`H_2`$ at $`\left|H_1\right|=v_1`$ and $`\left|H_2\right|=v_2`$ with $`M_W^2=\frac{1}{2}g^2(v_1^2+v_2^2)`$.
As was pointed out, the constraints on the APS arise from directions in the field space along which the potential becomes lower than $`V_{realmin}`$ (and may become unbounded from below). However, the minimization of the full potential $`V`$ is rather cumbersome. On the otherhand, just minimizing the tree-level potential at the weak scale neglecting $`\mathrm{\Delta }V_1`$ can lead to quite erroneous conclusions about the minimum point in the field space as was shown by . As a compromise, one evaluates $`V`$ at a judiciously chosen scale $`\widehat{Q}`$ where the one-loop correction is minimum. As is evident from eq. (7), this scale should be about the typical SUSY mass scale $`M_S`$ so that the large logarithmic terms tend to vanish.
The dangerous directions are selected in such a way that the positive definite F-terms vanish and the D-terms either cancel each other or their magnitudes can be kept under control. There are several other guidelines as discussed in . Using these conditions one can get the following UFB potentials.
UFB-1: The condition
$$m_{1}^{}{}_{}{}^{2}+m_{2}^{}{}_{}{}^{2}2\left|m_3^2\right|,$$
(8)
which is known as UFB-1, must be satisfied at any scale $`\widehat{Q}>M_S`$, and particularly at the unification scale $`\widehat{Q}=M_G`$, to have a realistic minimum of the scalar potential. From eq. (8), small $`m_{H_1}^2`$ (and $`m_{H_2}^2`$) makes the UFB-1 condition severely restrictive. This may be the case for large $`\mathrm{tan}\beta `$ and large negative values of $`A_0`$. The variations of $`m_{H_1}^2`$ and $`m_{H_2}^2`$ in the conventional scenario with respect to the common trilinear coupling $`A_0`$ for $`\mathrm{tan}\beta =30`$ and 45, corresponding to $`b\tau `$ and $`tb\tau `$ Yukawa unification respectively, are illustrated in fig. 1. From the figure we find that negative values of $`A_0`$ drive $`m_{H_2}^2`$ to large negative values in both the cases. For $`m_{H_1}^2`$ the effect is prominent for large $`\mathrm{tan}\beta `$, while for $`\mathrm{tan}\beta =30`$, $`m_{H_1}^2`$ remains positive for most of the range of $`A_0`$ that we have studied. The plot of $`m_H^2`$ vs $`A_0`$ is also given for different values of $`m_{1/2}`$ and $`m_{16}`$ for $`\mathrm{tan}\beta =45`$ in figures 2 and 3 respectively. From these plots it is clear that for large $`m_{1/2}`$ and/or $`m_{16}`$, and large negative $`A_0`$, $`m_{H_1}^2`$ and $`m_{H_2}^2`$ decrease significantly, so that these values of $`m_{1/2}`$ and $`m_{16}`$ become vulnerable to UFB-1. This is the reason why in the large $`\mathrm{tan}\beta `$ case the UFB-1 constraint plays a very significant role in restricting the APS.
UFB-2: The doublet slepton (along the sneutrino direction) and both $`H_1`$ and $`H_2`$ are given nonzero VEVs. For any value of $`\left|H_2\right|<M_G`$ satisfying
$$\left|H_2\right|^2>\frac{4m_{L_i}^2}{(g_{}^{}{}_{}{}^{2}+g_{2}^{}{}_{}{}^{2})[1\frac{\left|m_3\right|^4}{\mu ^4}]},$$
(9)
and provided that
$$\left|m_3^2\right|<\mu ^2(=m_{1}^{}{}_{}{}^{2}m_{L_i}^{}{}_{}{}^{2}),$$
(10)
the UFB-2 potential is given by
$$V_{UFB2}=[m_{2}^{}{}_{}{}^{2}+m_{L_i}^2\frac{\left|m_3\right|^4}{\mu ^2}]\left|H_2\right|^2\frac{2m_{L_i}^4}{g_{}^{}{}_{}{}^{2}+g_{2}^{}{}_{}{}^{2}}$$
(11)
At any momentum scale $`\widehat{Q}`$, this should be greater than $`V_{realmin}`$ for a stable configuration:
$$V_{UFB2}(Q=\widehat{Q})>V_{realmin}(Q=M_S)$$
(12)
where $`\widehat{Q}Max(g_2\left|H_2\right|,\lambda _{top}\left|H_2\right|,M_S)`$. However, we find that UFB-2 hardly rules out any further region of the APS which passes the UFB-1 and UFB-3 constraints, so it is of limited interest to us.
UFB-3: The convention is to choose $`H_1=0`$ and to cancel the $`H_1`$ F-term (which is a combination of $`H_2`$ and $`d_{L_i},d_{R_i}`$ or $`e_{L_i},e_{R_i}`$) with suitable VEVs to $`H_2`$ and the abovementioned slepton or squark directions. However, it is economical to give VEVs to the doublet fields (along $`T_3=1/2`$ direction) rather than the singlet fields to cancel both SU(2) and U(1) D-terms at the same stroke. Suppose the sleptons are given VEV; then for any values of $`\left|H_2\right|<M_G`$ satisfying
$$\left|H_2\right|>\sqrt{\frac{\mu ^2}{4\lambda _{e_j}}+\frac{2m_{L_i}^2}{g_{}^{}{}_{}{}^{2}+g_{2}^{}{}_{}{}^{2}}}\frac{\left|\mu \right|}{2\lambda _{e_j}},$$
(13)
the UFB-3 potential is defined as
$$V_{UFB3}=[m_{2}^{}{}_{}{}^{2}\mu ^2+m_{L_i}^2]\left|H_2\right|^2+\frac{\left|\mu \right|}{\lambda _{e_j}}[m_{L_j}^2+m_{e_j}^2+m_{L_i}^2]\left|H_2\right|\frac{2m_{L_i}^{}{}_{}{}^{4}}{g_{}^{}{}_{}{}^{2}+g_{2}^{}{}_{}{}^{2}}.$$
(14)
If $`\left|H_2\right|`$ does not satisfy (13), the formula changes to
$$V_{UFB3}=[m_{2}^{}{}_{}{}^{2}\mu ^2]\left|H_2\right|^2+\frac{\left|\mu \right|}{\lambda _{e_j}}[m_{L_j}^2+m_{e_j}^2]\left|H_2\right|+\frac{1}{8}(g_{}^{}{}_{}{}^{2}+g_{2}^{}{}_{}{}^{2})\left[\left|H_2\right|^2+\frac{\left|\mu \right|}{\lambda _{e_j}}\left|H_2\right|\right]^2$$
(15)
with $`ij`$. Note that we could substitute squarks for sleptons, where $`i=j`$ is allowed. The constraints on the parameter space arise from the requirement
$$V_{UFB3}(Q=\widehat{Q})>V_{realmin}(Q=M_S)$$
(16)
where $`\widehat{Q}`$ is chosen to be $`\widehat{Q}Max(g_2\left|e\right|,g_2\left|H_2\right|,\lambda _{top}\left|H_2\right|,g_2\left|L_i\right|,M_S)`$ to minimize $`\mathrm{\Delta }V_1`$. The VEVs are not arbitrary; they satisfy
$$\left|e\right|=\sqrt{\left|H_2\right|\left|\mu \right|/\lambda _{ej}},\left|L_i^2\right|=\left(\left|H_2\right|^2+\left|e\right|^2\right)4\frac{m_{L_i}^2}{(g_{}^{}{}_{}{}^{2}+g_{2}^{}{}_{}{}^{2})}.$$
(17)
As can be seen from eq. (14), the region of the parameter space where $`m_{H_2}^2=m_{2}^{}{}_{}{}^{2}\mu ^2`$ is large and negative is very susceptible to be ruled out by the UFB-3 condition. This is because the first term of eq. (14) may become negative in this case. However, the second term in (14), which is positive definite, may become competitive in certain cases (e.g., for $`j=1`$, when the Yukawa coupling in the denominator is small), which directions one should avoid when looking for the dangerous minima.
$`V_{UFB3}`$ with sleptons was found to be the strongest among all the UFB and CCB constraints in the low $`\mathrm{tan}\beta `$ case . In order to get the optimum result one has to take the largest $`\lambda _{e_j}`$ in the second term of eq. (14), which leads to the choice $`e_j=\stackrel{~}{\tau }_R`$. Now the restriction $`ij`$ requires $`L_i=\stackrel{~}{e}_L`$ or $`\stackrel{~}{\mu }_L`$ and excludes the choice $`\stackrel{~}{\tau }_L`$. In the low $`\mathrm{tan}\beta `$ case this restriction, however, is of little consequence since all the left sleptons are degenerate to a very good approximation.
The UFB -3 constraint with squarks may also be imposed by the following replacements in (14):
$$ed,\lambda _{e_j}\lambda _{d_j},L_jQ_j$$
(18)
(see eq. (31) of ). Now $`i`$ may be equal to $`j`$ and $`\widehat{Q}Max(g_2\left|d\right|,g_2\left|H_2\right|,\lambda _{top}\left|H_2\right|,g_2\left|L_i\right|,M_S)`$.
Now the optimum choice is $`d_j=\stackrel{~}{b}_R`$. However, since the choice $`i=j`$ is permitted, $`L_i=\stackrel{~}{\tau }_L`$ is not excluded. At high negative $`A_0`$ and at high $`\mathrm{tan}\beta `$, $`m_{\stackrel{~}{\tau _L}}`$ becomes smaller than the corresponding mass parameters of the first two generations. The variation of left-handed slepton mass parameters with $`A_0`$ (for $`\mathrm{tan}\beta =45`$) is shown in fig. 4. This relatively small $`m_{\stackrel{~}{\tau _L}}`$ at high $`\mathrm{tan}\beta `$ may make the alternative choice (18) more restrictive than the UFB-3 condition with sleptons. This, in fact, has been supported by our numerical computations.
## 3 Results
We now briefly review our methodology for implementing the Yukawa unification and computing the spectrum which is based on the computer program ISASUGRA, a part of the ISAJET package. We use the ISAJET version 7.48 .
To calculate Yukawa couplings at $`\widehat{Q}=M_Z`$, we start with the pole masses $`m_b=4.9`$ GeV, $`m_\tau =1.784`$ GeV and $`m_t=175`$ GeV. At $`\widehat{Q}=M_Z`$ the SUSY loop corrections to $`m_b`$ and $`m_\tau `$ is included using the approximate formulae from ref. . For the top quark Yukawa coupling this correction is added at $`\widehat{Q}=m_t`$. Starting with the three gauge couplings and the $`t,b`$ and $`\tau `$ Yukawa couplings, we evolve them upto the energy scale $`M_G`$. Now the boundary conditions are imposed on the soft breaking parameters according to the conventional or the nonuniversal scenario, while trial values for the $`\mu `$ and $`B`$ parameters are taken. Then all parameters are evolved down to the weak scale $`M_Z`$. The parameters $`\mu `$ and $`B`$ are then tentatively fixed at $`\widehat{Q}=\sqrt{m_{\stackrel{~}{t_L}}m_{\stackrel{~}{t_R}}}`$ by the radiative SU(2)$`\times `$U(1) breaking conditions. Using the particle spectrum so obtained, we compute the radiative corrections to the SU(2)$`\times `$U(1) breaking condition, and hence obtain the corrected result for $`\mu `$ and $`B`$. The whole proccess is then repeated iteratively until a stable solution within a reasonable tolerance is achieved. While running down from $`M_G`$, the SUSY thresholds are properly taken care of. The renormalization group (RG) equations that we use are upto two-loop for both the gauge couplings and the Yukawa couplings.
The demand of the Yukawa coupling unification at $`M_G`$ puts an extra constraint on $`\mathrm{tan}\beta `$. We require unification within an accuracy of 5% for $`Y_b`$ and $`Y_\tau `$ and 10% for $`Y_t,Y_b`$ and $`Y_\tau `$. The accuracy for the latter is relaxed since there are more uncertain factors, e.g., the choice of the Higgs sector. We define three variables $`r_{b\tau },r_{tb}`$ and $`r_{t\tau }`$ where generically $`r_{xy}=Max(Y_x/Y_y,Y_y/Y_x)`$. To check whether the couplings unify, we select only those points in the parameter space where $`Max(r_{b\tau },r_{tb},r_{t\tau })<1.10`$ (for $`tb\tau `$ unification) and $`r_{b\tau }<1.05`$ (for $`b\tau `$ unification). The quark Yukawa couplings depend on $`\alpha _s(M_Z)`$ which comes out from the gauge unification conditions to be $`0.118`$. Then we impose the experimental constraints $`m_{\chi ^+}>95`$ GeV, $`m_h>85.2`$ GeV and $`m_{\tau _1}>73`$ GeV, and require the lightest neutralino to be the lightest SUSY particle (LSP). These constraints filter out the APS on which the potential minima conditions UFB-1 and UFB-3 should apply.
Using $`\mu ,B`$, the gauge and the Yukawa couplings at the GUT scale alongwith the boundary conditions there, we generate the mass spectrum at any scale $`\widehat{Q}`$ using the 26 RG equations of the MSSM. In fig. 5, we show the lightest $`\stackrel{~}{\tau }`$, $`\stackrel{~}{t}`$ and $`\stackrel{~}{b}`$ masses at the weak scale as functions of $`A_0`$ for the conventional scenario with $`b\tau `$ Yukawa unification at $`m_{16}=1`$ TeV, $`m_{1/2}=500`$ GeV (this particular point, for the range of $`A_0`$ shown, is allowed by all constraints that we have considered). Note that for $`A_0<1.8`$ TeV, the lightest stop is the next lightest SUSY particle (NLSP), and is perfectly in the accessible range of the LHC.
We demand the electroweak symmetry to be unbroken at $`M_G`$. The Higgs potential is minimized at $`\widehat{Q}=\sqrt{m_{\stackrel{~}{t_L}}m_{\stackrel{~}{t_R}}}`$. The proper scale for the UFB potential where the one-loop effects are minimized, as discussed after eq. (15), is chosen by an iterative process within 1% accuracy. Usually a few iterations are sufficient. The UFB potential is calculated for different $`|H_2|`$ values ranging from zero to $`M_G`$, using a logarithmic scale. For each value of $`\left|H_2\right|`$ we compare the UFB potential with the scalar potential of MSSM, and whenever $`V_{UFB}<V_{realmin}`$, that particular region in the parameter space is marked as disallowed.
It should be emphasized that if the model is subject to the constraint of $`b\tau `$ Yukawa unification alone, the allowed region of the $`m_{1/2}m_{16}`$ parameter space increases as $`A_0`$ becomes more negative. The additional regions of the parameter space thus opened up are, however, severely restricted by the stability conditions on the potential. As a result the region allowed by Yukawa unification in conjunction with the stability of the potential is restricted to a rather small region even for large negative values of $`A_0`$. This will be illustrated by the following numerical results.
We begin our discussion for the allowed parameter space in the $`m_{1/2}`$-$`m_{16}`$ plane for$`A_0=2m_{16}`$ (see fig. 6) in the conventional scenario. At each point $`\mu `$ and $`\mathrm{tan}\beta `$ have been fixed by the radiative electroweak breaking condition and $`b\tau `$ Yukawa unification (at an accuracy of 5%) respectively. As expected from the discussions of the last section, the UFB-1 condition severely restricts the APS for relatively large $`m_{1/2}`$ and $`m_{16}`$. For smaller values of these parameters, the UFB-3 condition takes over; it is interesting to note that for relatively small $`m_{1/2}`$ and $`m_{16}`$, relevant for SUSY searches at the LHC, this condition rules out a small but interesting region of the parameter space. As a result for each $`m_{1/2}`$ there is an upper limit on $`m_{16}`$ and vice versa. Thus for $`m_{16}`$ = 500 (700, 900) GeV we find the gluino mass $`m_{\stackrel{~}{g}}`$ to be definitely less than 749 (1189, 1820) GeV respectively. It may be recalled that in the conventional scenario there is already a lower limit of approximately 300 GeV on $`m_{\stackrel{~}{g}}`$ from the direct searches at the Tevatron . On the otherhand, for $`m_{1/2}`$ = 200 (400, 800) GeV both upper and lower bounds on $`m_{16}`$ emerge, and we get 590 GeV (1010, 1775) $`<m_{\stackrel{~}{q}}<`$ 1170 GeV (1690, 2200) where $`m_{\stackrel{~}{q}}`$ is the average squark mass. Once SUSY signals are seen at the LHC this highly predictive model can be tested.
It may be argued that the accuracy to which Yukawa unification holds is worse than 5% due to the uncertainties discussed in the introduction. Relaxing the accuracy to 10% the region of the parameter space allowed by Yukawa unification alone expands. On imposing various stability conditions, we find that the UFB-1 constraints become somewhat weaker. However, the additional points allowed, especially the ones for low $`m_{16}`$, are disallowed by the UFB-3 condition which becomes stronger in this case. As a result the upper bounds on $`m_{1/2}`$ for relatively small values of $`m_{16}`$ presented in the last paragraph remain more or less unaltered.
For smaller negative values of $`A_0=m_{16}`$ the UFB constraints become less effective as may be seen from fig. 7. However, the APS is already quite restricted due to the requirement of Yukawa unification alone (this is the complementarity that we have talked about in the introduction). Although the bulk of the restricted APS can be probed at the LHC, a significant region remains inaccessible to it.
As we keep on increasing $`A_0`$ (in an algebraic sense) the UFB conditions start losing their effectiveness. For $`A_0=0`$ none of these conditions have any further usefulness in constraining the APS; see fig. 8. However, the stranglehold of Yukawa unification on the APS suffices by itself to predict a restrictive mass spectrum. The $`m_{16}`$-$`m_{1/2}`$ plot is bounded from both below and above, and a significant part of this APS can be probed at the LHC.
For $`A_0>0`$ the UFB conditions become ineffective. Yukawa unification alone yields a loosely restricted APS but most of it lies beyond the kinematic reach of the LHC.
We next focus our attention on the nonuniversal scenario $`m_{10}m_{16}`$. For a given $`A_0`$, the parameter space allowed by $`b\tau `$ Yukawa unification expands considerably from the conventional scenario for $`m_{10}<m_{16}`$. This is illustrated for $`A_0=2m_{16}`$ in fig. 9 which should be compared with fig. 6. In this case Yukawa unification is achieved for relatively low $`\mathrm{tan}\beta `$, which in turn makes $`m_{H_1}^2`$ less negative and hence the UFB-1 constraint weaker to some extent. However, many of the new points so allowed are eaten up by the UFB-3 condition. As a result, there is an upper bound on the allowed values of $`m_{1/2}`$ for the range of $`m_{16}`$ studied by us $`(m_{16}<3`$ TeV). Moreover, the gluino is most likely to be observed at the LHC for this entire range. Also the theoretical lower bound on $`m_{16}`$ gets stronger. However, the UFB conditions become ineffective as the magnitude of $`A_0`$ is reduced keeping its sign negative. At $`A_0=0`$ hardly any point is ruled out by these UFB conditions. This trend is similar to what we obtained for the conventional scenario.
On the otherhand, for $`m_{10}>m_{16}`$ the APS due to Yukawa unification alone is reduced quite a bit. This is illustrated in fig. 10 with $`m_{10}=1.2m_{16}`$, which should be compared with figures 6 and 9. For relatively large $`m_{16}`$, UFB-1 is a strong constraint as before, while some portion in the low $`m_{16}`$ region may be ruled out by the UFB-3 condition. We see that the lower bound on $`m_{16}`$ is significantly weaker than that in the previous case and $`m_{16}`$ as low as 300 GeV is allowed. The upper bound on $`m_{1/2}`$ is also weakened considerably. Yet an observable gluino is predicted over most part of the APS.
We now consider the scenario with $`tb\tau `$ Yukawa unification (within an accuracy of 10%) in the conventional scenario. The UFB-1 condition completely rules out the APS allowed by the unification criterion alone for $`A_00`$. (UFB-2 and UFB-3 conditions do not play any major role in constraining the APS.) For $`A_0>0`$, the APS (allowed by Yukawa unification) expands gradually; though a portion of it is ruled out by the UFB-1 constraint, a significant amount still remains allowed, and a sizable fraction of it is accessible at the LHC. The UFB-1 condition gets weaker as we go to larger values of $`A_0`$. In fig. 11, we show the allowed region for $`A_0=0.3m_{16}`$ and $`m_{1/2}=m_{16}`$; in fig. 12, we introduce nonuniversality by setting $`m_{10}=1.2m_{16}`$. Note that in the latter case the APS allowed by Yukawa unification alone is somewhat smaller than that in the conventional MSUGRA scenario.
Lastly, if the accuracy of the Yukawa unification is reduced to some extent (say, to 20%) the APS allowed by the unification criterion alone is significantly enhanced. However, the UFB-1 constraint rules out a large amount of this space, and only a small portion survives for negative $`A_0`$.
To summarize, the APS for large negative $`A_0`$ is so restricted by the UFB conditions that one should be able, with a bit of luck, to test the Yukawa unification models that we have discussed at the LHC by checking the squark and gluino masses. This restriction weakens if one goes to algebraically larger values of $`A_0`$. The quantitative nature obviously depends on the model chosen.
## 4 Conclusions
We have analyzed the consequences of both $`b\tau `$ and $`tb\tau `$ Yukawa unifications in conjunction with the UFB conditions in the MSUGRA scenario. In the forrmer case, for $`A_0<0`$, these two constraints nicely complement each other in restricting the APS; when one is weak, the other is sufficiently strong (see figures 6 and 7). For $`A_00`$ the UFB constraints are rather weak. However, the requirement of Yukawa unification at an accuracy less than $`5\%`$ by itself squeezes the APS sufficiently. As a result, both $`m_{1/2}`$ and $`m_{16}`$ are bounded (see fig. 8 for details) from below as well as above. Bulk of this restricted APS is within the striking range of the LHC. For large positive values of $`A_0`$, both the UFB conditions and the Yukawa unification constraint weaken and a large region of the parameter space accessible at the LHC is permitted.
The most restrictive model that we have studied is the one with $`A_0=2m_{16}`$. Here, mainly due to the UFB-3 constraint, one obtains $`m_{\stackrel{~}{g}}<2`$ TeV for $`m_0<1`$ TeV. Such gluinos are obviously within the reach of the LHC. For $`A_0>0`$, the UFB constraints lose their effectiveness and the loosely restricted APS is rather large.
If the accuracy of $`b\tau `$ unification is relaxed, the APS tends to increase as expected. However, the upper bound as mentioned above, viz., $`m_{\stackrel{~}{g}}<2`$ TeV for $`m_0<1`$ TeV more or less holds for large negative $`A_0`$, thanks to the UFB-3 condition.
The requirement of Yukawa unification is less effective in the nonuniversal scenario with $`m_{10}<m_{16}`$. Nevertheless the model on the whole is quite restrictive due to the UFB constraints. This is illustrated in fig. 9 for $`m_{10}=0.6m_{16}`$ and $`A_0=2m_{16}`$. Here a gluino observable at the LHC is almost definitely predicted for $`m_{16}<3`$ TeV. On the otherhand, for $`m_{10}>m_{16}`$, Yukawa unification by itself strongly constrains the APS (see fig. 10) and the UFB constraints play a subdominant role. Again $`m_{\stackrel{~}{g}}`$ is predicted to be observable at the LHC over most of the APS.
The masses of the third generation of sfermions are expected to be considerably lower than that of the first two generations for large values of $`\mathrm{tan}\beta `$. In fig. 5 we display in the $`b\tau `$ unification scheme, alongwith the UFB conditions, the masses of the lighter stop ($`\stackrel{~}{t}_1`$), sbottom ($`\stackrel{~}{b}_1`$) and stau ($`\stackrel{~}{\tau }_1`$) mass eigenstates as functions of $`A_0`$. They are indeed found to be considerably lighter than the sparticles belonging to the first two generations; in fact, the lighter stop could very well be the second lightest SUSY particle. Thus in spite of the restrictions imposed by the UFB conditions and Yukawa unification, light third generation sfermions can be accomodated. In particular, the possibility that the lighter stop is the NLSP is open for large negative $`A_0`$.
The $`tb\tau `$ Yukawa unification models, with a unification accuracy of 10%, are definitely ruled out for $`A_00`$ in the conventional scenario. For positive values of $`A_0`$, the UFB-1 condition is less severe, and a portion of the parameter space remains allowed, of which a sizable fraction should be accessible at the LHC. If we relax the accuracy for unification, the APS increases, most of which could be ruled out by the UFB-1 condition.
Acknowledgements
AD thanks Prof. E. Reya for hospitalities at the University of Dortmund. He further thanks both E. Reya and M. Glück for many discussions on the dangerous directions of the scalar potential in supersymmetric theories, and A. Stephan for helps in computation. His work was supported by DST, India (Project No. SP/S2/k01/97) and BRNS, India (Project No. 37/4/97 - R & D II/474). AS acknowledges CSIR, India, for a research fellowship.
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# On |𝑉_{𝑢𝑏}| from the 𝐵̄→𝑋_𝑢ℓ𝜈̄ dilepton invariant mass spectrum^∗
## Acknowledgements
This work was supported in part by the Natural Sciences and Engineering Research Council of Canada and the Sloan Foundation. Fermilab is operated by Universities Research Association, Inc., under DOE contract DE-AC02-76CH03000.
## References
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# I Introduction
## I Introduction
Supergravity solutions representing D-branes, NS 5-branes, and M-branes have been known for a long time now . It is also known that strings can end on D-branes and branes can end on other branes . However, supergravity solutions representing such intersecting objects are, by and large, unknown.
Previous work in the literature has mostly focused on cases where one of the branes is delocalized in some direction . There are only a few cases where all the branes have been localized , but these appear to be found on a case by case basis.
In this paper, we will provide a completely general method for finding supergravity solutions representing intersecting branes. This is based on the method introduced by (See also ). Using this approach, which we describe in more detail below, we are able to write the solution in terms of the solution to a single nonlinear differential equation. Unfortunately, we have been unable to find explicit solutions of this differential equation. Nevertheless, we regard this as a solution in principle, and it is far simpler than attempting to solve the coupled Maxwell-Einstein equations.
The idea in this approach is to use the supersymmetry equations (the equations for conserved Killing spinors) as linear equations determining the field strengths in terms of the metric. This enables us to write all the fields in terms of the parameters of the metric.
It is then straightforward to substitute these expressions into the equations of motion. By the miracle of supersymmetry, we find that all the equations of motion simplify into a few simple equations. All the parameters can then be determined in terms of one function, which is in turn a solution to a nonlinear differential equation. Although we have been unable to find explicit solutions to this nonlinear equation, nevertheless, we emphasize that the solution is completely determined in principle by this differential equation.
We carry out these steps in great detail for the case of M2-branes ending on M5-branes. This method can be easily generalized to all other systems of intersecting branes in M-theory and string theory, and as a further example, we present solutions for strings ending on Dp-branes. Finally, we list several open problems.
## II M2-branes ending on M5-branes
### A Setup
The notation will be as follows. Indices with a tilde over them will denote curved (world) indices, and indices without a tilde will be tangent indices.
The M5-brane is oriented along the directions $`x_0`$ (time), $`x_1,x_2,x_3,x_4,x_5`$. The M2-brane is oriented along the directions $`x_0,x_1,x_6`$. Indices $`(2,3,4,5)`$ will be labeled by $`a,b,\mathrm{}`$, and indices $`(7,8,9,10)`$ will be labeled $`\alpha ,\beta ,\mathrm{}`$.
The M5-brane and M2-brane act as sources for $`a_{\alpha \beta \gamma }^{(3)}`$ and $`a_{016}^{(3)}`$ respectively. Correspondingly, we expect the field strengths $`f_{\alpha \beta \gamma \delta },f_{\alpha \beta \gamma a},f_{\alpha \beta \gamma 6}`$, and $`f_{01a6},f_{016\alpha }`$ to be nonzero. Furthermore, the M2 brane acts as a source for the $`U(1)`$ field on the worldvolume of the M5-brane. If the worldvolume field strength is denoted by $`h`$, then $`h_{01a}`$ and $`h_{abc}`$ are nonzero. This field acts as a source for $`a_{abc}^{(3)}`$ and $`a_{01a}^{(3)}`$. Correspondingly, we also expect $`f_{01a\alpha }`$, and $`f_{2345},f_{abc6},f_{abc\alpha }`$ to be nonzero.
Also, since $`_aX^6`$ is nonzero on the M5 brane worldvolume, we expect $`g_{a6}`$ to be nonzero.
The nonzero components of the field strengths will be denoted
$`F_{\alpha \beta \gamma \delta },F_{\alpha \beta \gamma a},F_{\alpha \beta \gamma 6};G_{01a6},G_{01a\alpha },G_{016\alpha };H_{2345},H_{abc6},H_{abc\alpha }`$ (1)
We have chosen a notation where different components of the field strength are denoted by $`F,G,`$ or $`H`$ depending on the indices.
The symmetry further dictates that $`(U=\sqrt{x_\alpha ^2})`$
$`F_{8910a}={\displaystyle \frac{x_7}{U}}F_aF_{7910a}={\displaystyle \frac{x_8}{U}}F_aF_{7810a}={\displaystyle \frac{x_9}{U}}F_aF_{789a}={\displaystyle \frac{x_{10}}{U}}F_a`$ (2)
$`F_{89106}={\displaystyle \frac{x_7}{U}}F_6F_{79106}={\displaystyle \frac{x_8}{U}}F_6F_{78106}={\displaystyle \frac{x_9}{U}}F_6F_{7896}={\displaystyle \frac{x_{10}}{U}}F_6`$ (3)
$`H_{abc\alpha }={\displaystyle \frac{x_\alpha }{U}}H_{abc}G_{016\alpha }={\displaystyle \frac{x_\alpha }{U}}G_6G_{01a\alpha }={\displaystyle \frac{x_\alpha }{U}}G_a`$ (4)
### B The Equations for Supersymmetry
We can now analyse the equation for conserved supercharges
$`\delta \psi _{\stackrel{~}{\mu }}=_{\stackrel{~}{\mu }}ϵ{\displaystyle \frac{1}{4}}\omega _{\stackrel{~}{\mu }}^{ab}\gamma _{ab}+{\displaystyle \frac{i}{288}}(\gamma _{\stackrel{~}{\mu }}^{\alpha \beta \gamma \delta }8\delta _{\stackrel{~}{\mu }}^\alpha \gamma ^{\beta \gamma \delta })f_{\alpha \beta \gamma \delta }ϵ=0`$ (5)
These equations are rather cumbersome. For example, for $`\mu =1`$, we have
$`\left({\displaystyle \frac{1}{2}}\omega _1^{1a}\gamma _a+{\displaystyle \frac{1}{2}}\omega _1^{16}\gamma _6+{\displaystyle \frac{1}{2}}\omega _1^{1\alpha }\gamma _\alpha \right)ϵ{\displaystyle \frac{i}{12}}F_{78910}\gamma ^{78910}ϵ{\displaystyle \frac{i}{12}}H_{2345}\gamma ^{2345}ϵ`$ (6)
$`{\displaystyle \frac{i}{12}}\left(F_{789a}\gamma ^{789a}+F_{7810a}\gamma ^{7810a}+F_{7910a}\gamma ^{7910a}+F_{8910a}\gamma ^{8910a}\right)ϵ`$ (7)
$`{\displaystyle \frac{i}{12}}\left(F_{7896}\gamma ^{7896}+F_{78106}\gamma ^{78106}+F_{79106}\gamma ^{79106}+F_{89106}\gamma ^{89106}\right)ϵ`$ (8)
$`{\displaystyle \frac{i}{12}}\left(H_{234\alpha }\gamma ^{234\alpha }+H_{235\alpha }\gamma ^{235\alpha }+H_{245\alpha }\gamma ^{245\alpha }+H_{345\alpha }\gamma ^{345\alpha }\right)ϵ`$ (9)
$`{\displaystyle \frac{i}{12}}\left(H_{2346}\gamma ^{2346}+H_{2356}\gamma ^{2356}+H_{2456}\gamma ^{2456}+H_{3456}\gamma ^{3456}\right)ϵ`$ (10)
$`+{\displaystyle \frac{i}{6}}\left(G_{01a\alpha }\gamma ^{01a\alpha }+G_{01a6}\gamma ^{01a6}+G_{016\alpha }\gamma ^{016\alpha }\right)ϵ=0`$ (11)
All indices above are tangent space indices; we have defined $`\omega _a^{bc}=e_a^{\stackrel{~}{\mu }}\omega _{\stackrel{~}{\mu }}^{bc}`$.
To simplify the problem, we note that since we are looking for BPS solutions, we expect to be able to superpose them. In this case, we should be able to superpose M2 branes with arbitrary coordinates in the directions $`x_a`$. We can therefore separate the SUSY equations into terms which are even in $`x_a`$, and terms which are odd in $`x_a`$. These equations should be separately satisfied. The terms involving a $`_a`$ clearly produces a term odd in $`x_a`$, while $`_6,_\alpha `$ produce terms even in $`x_a`$.
In addition, we will take the ansatz for the spinor to be
$`ϵ=(g_{11})^{\frac{1}{4}}ϵ_0`$ (12)
where $`ϵ_0`$ is a constant spinor. This ansatz is justified in .
The parts of the equations odd in $`x_4`$ are
$`{\displaystyle \frac{1}{2}}\omega _1^{14}\gamma _4ϵ+\left({\displaystyle \frac{i}{12}}H_{235\alpha }\gamma ^{235\alpha }{\displaystyle \frac{i}{12}}H_{2356}\gamma ^{2356}+{\displaystyle \frac{i}{6}}G_{014\alpha }\gamma ^{014\alpha }+{\displaystyle \frac{i}{6}}G_{0146}\gamma ^{0146}\right)ϵ`$ (13)
$`{\displaystyle \frac{i}{12}}\left(F_{7894}\gamma ^{7894}+F_{78104}\gamma ^{78104}+F_{79104}\gamma ^{79104}+F_{89104}\gamma ^{89104}\right)ϵ=0`$ (14)
$`{\displaystyle \frac{1}{2}}(\omega _5^{54}\omega _1^{14})\gamma _4ϵ+\left({\displaystyle \frac{i}{4}}H_{235\alpha }\gamma ^{235\alpha }+{\displaystyle \frac{i}{4}}H_{2356}\gamma ^{2356}{\displaystyle \frac{i}{4}}G_{014\alpha }\gamma ^{014\alpha }{\displaystyle \frac{i}{4}}G_{0146}\gamma ^{0146}\right)ϵ=0`$ (15)
$`{\displaystyle \frac{1}{2}}\left(\omega _5^{54}\gamma _4+\omega _4^{67}\gamma _{467}ϵ\right)+\left({\displaystyle \frac{i}{4}}H_{235\alpha }\gamma ^{235\alpha }+{\displaystyle \frac{i}{4}}H_{2356}\gamma ^{2356}\right)ϵ`$ (16)
$`\left({\displaystyle \frac{i}{4}}G_{014\alpha }\gamma ^{014\alpha }{\displaystyle \frac{i}{4}}G_{0146}\gamma ^{0146}{\displaystyle \frac{i}{4}}F_{89104}\gamma ^{89104}\right)ϵ=0`$ (17)
$`{\displaystyle \frac{1}{2}}(\omega _6^{64}\omega _1^{14})\gamma _4ϵ{\displaystyle \frac{1}{2}}\omega _6^{4\alpha }\gamma _{64\alpha }ϵ+\left({\displaystyle \frac{i}{4}}H_{2356}\gamma ^{2356}{\displaystyle \frac{i}{4}}G_{014\alpha }\gamma ^{014\alpha }\right)ϵ=0`$ (18)
$`{\displaystyle \frac{1}{2}}\omega _7^{64}\gamma _{764}+\left({\displaystyle \frac{i}{4}}G_{0147}\gamma ^{0147}+{\displaystyle \frac{i}{4}}H_{2357}\gamma ^{2357}{\displaystyle \frac{i}{4}}F_{89104}\gamma ^{89104}\right)ϵ=0`$ (19)
The parts of the equations even in all $`x_a`$ are
$`{\displaystyle \frac{1}{2}}\omega _1^{16}\gamma _6ϵ+{\displaystyle \frac{1}{2}}\omega _1^{1\alpha }\gamma _\alpha ϵ+\left({\displaystyle \frac{i}{6}}G_{016\alpha }\gamma ^{016\alpha }{\displaystyle \frac{i}{12}}H_{2345}\gamma ^{2345}{\displaystyle \frac{i}{12}}F_{78910}\gamma ^{78910}\right)ϵ`$ (20)
$`{\displaystyle \frac{i}{12}}\left(F_{7896}\gamma ^{7896}+F_{78106}\gamma ^{78106}+F_{79106}\gamma ^{79106}+F_{89106}\gamma ^{89106}\right)ϵ=0`$ (21)
$`{\displaystyle \frac{1}{2}}(\omega _5^{56}\omega _1^{16})\gamma _6ϵ+{\displaystyle \frac{1}{2}}(\omega _5^{5\alpha }\omega _1^{1\alpha })\gamma _\alpha ϵ+\left({\displaystyle \frac{i}{4}}H_{2345}\gamma ^{2345}{\displaystyle \frac{i}{6}}G_{016\alpha }\gamma ^{016\alpha }\right)ϵ=0`$ (22)
$`{\displaystyle \frac{1}{2}}(\omega _6^{6\alpha }\omega _1^{1\alpha })\gamma _\alpha ϵ+{\displaystyle \frac{i}{4}}\left(F_{7896}\gamma ^{7896}+F_{78106}\gamma ^{78106}+F_{79106}\gamma ^{79106}+F_{89106}\gamma ^{89106}\right)ϵ=0`$ (23)
$`{\displaystyle \frac{1}{2}}(\omega _1^{17}\omega _8^{87})\gamma _7ϵ+{\displaystyle \frac{i}{4}}G_{0167}\gamma ^{0167}ϵ{\displaystyle \frac{i}{4}}F_{89106}\gamma ^{89106}ϵ=0`$ (24)
These equations should preserve one quarter of the supersymmetries. In other words, each equation should be proportional to a linear combination of $`P_1ϵ_0`$ and $`P_2ϵ_0`$, where $`P_1ϵ_0=P_2ϵ_0=0`$ and $`P_1,P_2`$ are projection operators.
A glance at the equations (13) and (20) shows that the only possible choice of projection operators is (upto signs)
$`P_1ϵ_0(1+i\gamma ^{678910})ϵ_0=0`$ (25)
$`P_2ϵ_0(1+i\gamma ^{016})ϵ_0=0`$ (26)
With these projection operators, we can reduce the matrix equations (13) and (20) to a set of algebraic equations. We will treat these algebraic equations as equations determinining the various field strengths as functions of the spin connections (and thereby implicitly as functions of the metric.)
The equations for the field strengths are then
$`G_{0164}=4w_1^{14}+2w_5^{54}`$ (27)
$`G_{0167}=2w_1^{17}2w_4^{47}`$ (28)
$`G_{0147}=2w_4^{67}`$ (29)
$`H_{2356}=4w_5^{54}+2w_1^{14}`$ (30)
$`H_{2357}=2w_4^{67}`$ (31)
$`H_{2345}=2w_4^{46}2w_1^{16}`$ (32)
$`F_{48910}=2w_4^{67}`$ (33)
$`F_{78910}=2w_4^{46}w_1^{16}`$ (34)
$`F_{68910}=4w_4^{47}+2w_1^{17}`$ (35)
In addition, we get the constraints
$`w_1^{16}+w_4^{46}+w_7^{76}=0`$ (36)
$`w_1^{18}+w_4^{48}+w_7^{78}=0`$ (37)
$`w_1^{14}+w_5^{54}+w_7^{74}=0`$ (38)
$`w_4^{67}=w_6^{47}=w_7^{64}`$ (39)
$`w_6^{67}+2w_4^{47}=0`$ (40)
$`w_6^{64}=2w_5^{54}+2w_1^{14}`$ (41)
In all the equations above, all the indices are tangent space indices; we have defined $`w_c^{ab}=e_c^{\stackrel{~}{\mu }}w_{\stackrel{~}{\mu }}^{ab}`$.
We can solve all the constraints above by the metric ansatz
$`e_{0\stackrel{~}{0}}=e_{1\stackrel{~}{1}}=\lambda ^{\frac{1}{3}}H^{\frac{1}{6}}`$ (42)
$`e_{2\stackrel{~}{2}}=e_{3\stackrel{~}{3}}=e_{4\stackrel{~}{4}}=e_{5\stackrel{~}{5}}=\lambda ^{\frac{1}{6}}H^{\frac{1}{6}}`$ (43)
$`e_{6\stackrel{~}{6}}=\lambda ^{\frac{1}{3}}H^{\frac{1}{3}}`$ (44)
$`e_{7\stackrel{~}{7}}=e_{8\stackrel{~}{8}}=e_{9\stackrel{~}{9}}=e_{10\stackrel{~}{10}}=\lambda ^{\frac{1}{6}}H^{\frac{1}{3}}`$ (45)
$`e_{6\stackrel{~}{a}}=\varphi _ae_{6\stackrel{~}{6}}`$ (46)
with the constraint
$`_6(H\varphi _a)=_aH`$ (47)
### C The Equations of Motion
We now look at the equations of motion. Away from the M5-brane, we can look at the vacuum equations of motion, which for the gauge fields is
$`_\mu f^{\mu \nu \rho \sigma }={\displaystyle \frac{1}{2.(24)^2}}ϵ^{\nu \rho \sigma \stackrel{~}{a}\stackrel{~}{b}\stackrel{~}{c}\stackrel{~}{d}\stackrel{~}{e}\stackrel{~}{f}\stackrel{~}{g}\stackrel{~}{h}}f_{\stackrel{~}{a}\stackrel{~}{b}\stackrel{~}{c}\stackrel{~}{d}}f_{\stackrel{~}{e}\stackrel{~}{f}\stackrel{~}{g}\stackrel{~}{h}}`$ (48)
We can now substitute the field strengths found above into these equations. Remarkably, all the equations of motion collapse to the equations
$`_a\left(\varphi _a\right){\displaystyle \frac{1}{H}}_6\lambda {\displaystyle \frac{1}{2}}_6\varphi _a^2=0`$ (49)
and
$`_\alpha ^2\left(H\right)+_6^2(H\lambda )=0`$ (50)
When we include sources, we need to modify the second equation. Defining
$`_\alpha ^2\left(H\right)+_6^2(H\lambda )=_6Q`$ (51)
we have
$`_7F_{a8910}_8F_{a7910}+_9F_{a7810}_{10}F_{a789}_aF_{78910}=_aQ`$ (52)
$`_\mu G^{\mu 01a}+{\displaystyle \frac{1}{2.(24)^2}}ϵ^{\stackrel{~}{0}\stackrel{~}{1}\stackrel{~}{a}\stackrel{~}{A}\stackrel{~}{B}\stackrel{~}{C}\stackrel{~}{D}\stackrel{~}{E}\stackrel{~}{F}\stackrel{~}{G}\stackrel{~}{H}}f_{\stackrel{~}{A}\stackrel{~}{B}\stackrel{~}{C}\stackrel{~}{D}}f_{\stackrel{~}{E}\stackrel{~}{F}\stackrel{~}{G}\stackrel{~}{H}}=_aQi\varphi _6Q`$ (53)
$`_6\left(_\mu G^{\mu 016}+{\displaystyle \frac{1}{2.(24)^2}}ϵ^{\stackrel{~}{0}\stackrel{~}{1}\stackrel{~}{6}\stackrel{~}{A}\stackrel{~}{B}\stackrel{~}{C}\stackrel{~}{D}\stackrel{~}{E}\stackrel{~}{F}\stackrel{~}{G}\stackrel{~}{H}}f_{\stackrel{~}{A}\stackrel{~}{B}\stackrel{~}{C}\stackrel{~}{D}}f_{\stackrel{~}{E}\stackrel{~}{F}\stackrel{~}{G}\stackrel{~}{H}}\right)=_a^2Q`$ (54)
$`Q`$ thus parametrizes the perturbation of the original M5-brane due to the addition of the M2-brane.
### D Linearized analysis
At the linearized level, we can drop the terms involving $`\varphi `$. If we then take
$`_\mu G^{\mu \stackrel{~}{0}\stackrel{~}{1}\stackrel{~}{a}}+{\displaystyle \frac{1}{2.(24)^2}}ϵ^{\stackrel{~}{0}\stackrel{~}{1}\stackrel{~}{a}\stackrel{~}{A}\stackrel{~}{B}\stackrel{~}{C}\stackrel{~}{D}\stackrel{~}{E}\stackrel{~}{F}\stackrel{~}{G}\stackrel{~}{H}}f_{\stackrel{~}{A}\stackrel{~}{B}\stackrel{~}{C}\stackrel{~}{D}}f_{\stackrel{~}{E}\stackrel{~}{F}\stackrel{~}{G}\stackrel{~}{H}}=\delta (x_\alpha )\delta (x_6)h^{01a}`$ (55)
$`_\mu G^{\mu \stackrel{~}{0}\stackrel{~}{1}\stackrel{~}{6}}+{\displaystyle \frac{1}{2.(24)^2}}ϵ^{\stackrel{~}{0}\stackrel{~}{1}\stackrel{~}{6}\stackrel{~}{A}\stackrel{~}{B}\stackrel{~}{C}\stackrel{~}{D}\stackrel{~}{E}\stackrel{~}{F}\stackrel{~}{G}\stackrel{~}{H}}f_{\stackrel{~}{A}\stackrel{~}{B}\stackrel{~}{C}\stackrel{~}{D}}f_{\stackrel{~}{E}\stackrel{~}{F}\stackrel{~}{G}\stackrel{~}{H}}=\{\begin{array}{cc}\delta ^4(x_a)\delta ^4(x_\alpha ),\hfill & x_60\text{;}\hfill \\ 0,\hfill & \text{else.}\hfill \end{array}`$ (56)
the above equations are satisfied provided
$`_ah^{01a}=\delta (x_a)`$ (57)
In other words, we have a gauge field strength on the M5-brane corresponding to the field of an electric charge on a line at $`x_a=0`$. Also, we have a membrane (the source for $`A_{016}`$) extending along the positive $`x_6`$ axis. This is exactly the configuration of a M2-brane ending on an M5-brane.
At the linearized level, the equations (47), (49) and (51) also simplify to
$`_6(H\varphi _a)=_aH`$ (58)
$`_a(H\varphi _a)=_6\lambda `$ (59)
$`_\alpha ^2(H\varphi _a)+_a_6(H\lambda )=_aQ`$ (60)
which can be simplified to the linear equation
$`H_a^2\lambda +_6^2\lambda +_\alpha ^2\lambda ={\displaystyle \frac{1}{_6}}_a^2Q=\{\begin{array}{cc}\delta ^4(x_a)\delta ^4(x_\alpha ),\hfill & x_60\text{;}\hfill \\ 0,\hfill & \text{else.}\hfill \end{array}`$ (61)
More generally, we have multiple M2 branes ending at different points on the M5-brane. The general solution, at the linear level, is then
$`\lambda ={\displaystyle 𝑑x^{}K(x,x^{})\rho (x^{})}`$ (62)
where we have introduced the Green’s function K, satisfying
$`(H_0_a^2+_6^2+_\alpha ^2)K(x,x^{})=\delta (xx^{})`$ (63)
$`H_0`$ is the harmonic function describing the M5-branes without the M2-branes and $`\rho (x^{})`$ is the membrane source density.
The full nonlinear solution no longer has pointlike sources. The source $`Q`$ is now a nontrivial function of $`x_a`$ and $`x_6`$, and even the identification of $`Q`$ as a worldvolume field is a little problematic.
However, at long distances from the M2-branes, the linearized analysis still applies, since $`\lambda `$ is small. Also, very close to any M2-brane, where $`\lambda `$ diverges, the solution (62) should still apply, since the effects of the membrane dominate.
At the nonlinear level, we can then say that $`\lambda `$ is a solution to the equations (47), (49) and (51), subject to the condition that
$`\lambda ={\displaystyle 𝑑x^{}K(x,x^{})\rho (x^{})}`$ (64)
both when $`\lambda `$ is small or when it is very large. This is expected to completely specify the solution, although we cannot explicitly solve the equations.
### E Comments
1. One might ask why one needs to look at the equations of motions at all. After all, once we impose the constraint (47), we have a solution that is supersymmetric. It is well known that a solution preserving a supersymmetry should also satisfy the equations of motion.
The point is that imposing the supersymmetry constraints does not completely specify the sources. The M2-branes along $`x_0,x_1,x_6`$, and the M5-branes along $`x_0,x_1,x_2,x_3,x_4,x_5,`$ preserve the supersymmetries (25). But one can also add M5-branes oriented along $`x_0,x_1,x_7,x_8,x_9,x_{10}`$ without breaking any more supersymmetries. The equation (49) precisely sets the number of these M5-branes to zero, so that we have the sources we started with.
2. The second equation (51) determines the positions of the M5-brane and M2-brane sources.
One expects on physical grounds that the M5-branes should not deform in the $`x_\alpha `$ directions. In this case, one should find that $`Q=\delta ^4(x_\alpha )Q_1(x_a,x_6)`$.
It is however not clear how to show that the differential equations (47), (49), (51) are consistent with this expectation. One could, in principle, worry that when one tried to integrate in the equations from $`x_\alpha =\mathrm{}`$, one gets a singularity away from the origin. The solution would then be valid only outside this region (somewhat like the enhancon phenomenon ). It would be extremely surprising (at least to this author) if such a phenomenon should occur in this system, so we shall assume that in fact, $`Q=\delta ^4(x_\alpha )Q_1(x_a,x_6)`$.
### F Summary
The solution for M2-branes ending on M5-branes is given by the metric (42)
$`ds^2=\lambda ^{\frac{2}{3}}H^{\frac{1}{3}}(dx_0^2+dx_1^2)+\lambda ^{\frac{1}{3}}H^{\frac{1}{3}}(dx_2^2+dx_3^2+dx_4^2+dx_5^2)`$ (65)
$`+\lambda ^{\frac{1}{3}}H^{\frac{2}{3}}(dx_7^2+dx_8^2+dx_9^2+dx_{10}^2)+\lambda ^{\frac{2}{3}}H^{\frac{2}{3}}(dx_6+\varphi _adx_a)^2`$ (66)
The field strengths are given in equations (8)-(16).
We also have the equations (47), (49), and (51) which determine the parameters
$`_6(H\varphi _a)=_aH`$ (67)
$`_a\left(\varphi _a\right){\displaystyle \frac{1}{H}}_6\lambda {\displaystyle \frac{1}{2}}_6\varphi _a^2=0`$ (68)
$`_\alpha ^2\left(H\right)+_6^2(H\lambda )=\delta ^4(x_\alpha )_6Q`$ (69)
where $`Q`$ is a function of $`x_a,x_6`$.
Far away from the M2-branes, we have the boundary condition
$`HH_0\varphi _a0\lambda 1`$ (70)
The first linear perturbation in $`\lambda `$ is
$`\lambda ={\displaystyle 𝑑x^{}K(x,x^{})\rho (x^{})}+1`$ (71)
where $`\rho (x^{})`$ is the membrane source density, and $`K(x,x^{})`$ is the Green’s function satisfying
$`(H_0_a^2+_6^2+_\alpha ^2)K(x,x^{})=\delta (xx^{})`$ (72)
where $`H_0`$ is the harmonic function describing the M5-brane background without the M2-branes. Furthermore, $`\lambda `$ is expected to behave as the above equation (71) in regions where $`\lambda `$ diverges.
The system of equations (47), (49), and (51) can further be simplified to the following equations
$`H\varphi _a=_a_6\tau `$ (73)
$`H=_6^2\tau `$ (74)
$`\lambda =_a^2\tau H\varphi _a^2`$ (75)
where the function $`\tau `$ satisfies the nonlinear differential equation
$`_\alpha ^2\tau +_6^2\tau _a^2\tau (_a_6\tau )^2=\delta ^4(x_\alpha ){\displaystyle \frac{1}{_6}}Q`$ (76)
## III Strings ending on Dp-branes
We can carry out a very similar analysis for strings ending on Dp-branes. We shall take the Dp-branes oriented along the directions $`x_0..x_p`$, the string is oriented along $`x_0,x_9`$.
The Killing spinors preserved by this configuration are given by $`ϵ=(g_{00}^{1/4})ϵ_0`$, where $`ϵ_0`$ is a constant spinor satisfying
$`\gamma ^{09}ϵ_0=\eta _1ϵ_0`$ (77)
$`\gamma ^{01..p}ϵ_0=\eta _2ϵ_0`$ (78)
where $`\eta _1,\eta _2`$ are constants.
The Killing equations can be written in the generic form
$`w_\mu ^{ab}\gamma _{ab}ϵ+(a_GG_{\mu ab}\gamma ^{ab}+b_GG_{abc}\gamma _\mu ^{abc})ϵ+(a_HH_{\mu a_1..a_{p2}}\gamma ^{a_1..a_{p2}}+b_HH_{a_1..a_{p1}}\gamma _\mu ^{a_1..a_{p1}})ϵ`$ (79)
$`+(a_FF_{\mu a_1..a_p}\gamma ^{a_1..a_p}+b_FF_{a_1..a_p}\gamma _\mu ^{a_1..a_p})ϵ=0`$ (80)
The metric components are found by noting that the diagonal components are the products of the corresponding component in the metric of the string and the p-brane metric (i.e. $`g_{ii}=g_{ii}^{(string)}g_{ii}^{(pbrane)}`$.) (We will work in the string frame)
$`e_{0\stackrel{~}{0}}=\lambda ^{\frac{1}{2}}H^{\frac{1}{4}}`$ (81)
$`e_{1\stackrel{~}{1}}=e_{2\stackrel{~}{2}}=..=e_{p\stackrel{~}{p}}=H^{\frac{1}{4}}`$ (82)
$`e_{9\stackrel{~}{9}}=\lambda ^{\frac{1}{2}}H^{\frac{1}{4}}`$ (83)
$`e_{(p+1)\stackrel{~}{(p+1)}}=..=e_{8\stackrel{~}{8}}=H^{\frac{1}{4}}`$ (84)
In addition, we have an off diagonal component.
$`e_{9\stackrel{~}{a}}=\varphi _ae_{9\stackrel{~}{9}}`$ (85)
The dilaton is
$`e^\varphi =\lambda ^{1/2}H^{\frac{p3}{4}}`$ (86)
The equations for the field strengths are then (all indices are tangent space indices)
$`(a_G\eta _1)G_{019}=w_2^{21}w_0^{01}`$ (87)
$`(a_G\eta _1)G_{098}=w_0^{08}w_2^{28}`$ (88)
$`(a_G\eta _1)G_{018}=w_8^{19}`$ (89)
$`(a_H\eta _1\eta _2)e^\varphi H_{92..p}=w_2^{21}`$ (90)
$`(a_H\eta _1\eta _2)e^\varphi H_{82..p}={\displaystyle \frac{1}{2}}w_8^{19}`$ (91)
$`(a_H\eta _1\eta _2)e^\varphi H_{1..p}={\displaystyle \frac{1}{2}}(w_2^{29}w_0^{09})`$ (92)
$`(a_F\eta _2)e^\varphi F_{801..p}=w_2^{28}`$ (93)
$`(a_F\eta _2)e^\varphi F_{8902..p}=w_8^{19}`$ (94)
$`(a_F\eta _2)e^\varphi F_{901..p}={\displaystyle \frac{1}{2}}(w_0^{09}+w_2^{29})`$ (95)
The equations relating the various parameters are
$`_9(H\varphi _a)=_aH`$ (96)
$`H^1_9\lambda +{\displaystyle \frac{1}{2}}_9(\varphi _a^2)=_a\varphi _a`$ (97)
$`_9^2(H\lambda )+_\alpha ^2(H)=\delta (x_\alpha )_9Q`$ (98)
which can be simplified in terms of a single function $`\tau `$
$`H=_9^2\tau `$ (99)
$`H\varphi _a=_a_9\tau `$ (100)
$`\lambda =_a^2\tau H\varphi _a^2`$ (101)
where $`\tau `$ satisfies
$`_\alpha ^2\tau +_9^2\tau _a^2\tau (_a_6\tau )^2=\delta (x_\alpha ){\displaystyle \frac{1}{_9}}Q`$ (102)
which is the same differential equation as in the previous section.
## IV Open questions
We have reduced the construction of intersecting brane systems to solving a single nonlinear partial differential equation. Much of the physics is contained in the solution to this equation, which, unfortunately, we have been unable to find exactly. It would be useful to be able to extract some information from the equation, for instance, to know whether it admits multiple solutions in some cases.
There are several potential applications of these solutions. By the AdS/CFT correspondence , some of these solutions can be mapped to configurations in strongly coupled gauge theories. For example, the solution for $`k`$ D-strings ending on $`N`$ D3-brane is dual to a configuration of $`k`$ monopoles in $`SU(N)`$ $`𝒩=4`$ Yang-Mills theory. We can then extract quantities like the quark monopole potential etc. in this limit.
Similarly, these configurations, can be used to construct supergravity duals for any gauge theory that can be engineered by putting branes on orbifolds. These include pure $`𝒩=2`$ gauge theories and $`𝒩=1`$ gauge theories. One can also study solitons in these theories.
One very interesting open question is to find the near extremal solutions. It is, of course, not possible to use the supersymmetry equations there, but the solutions found here should be a useful starting point to get the nonextremal solutions. Such solutions would give us information about the thermodynamics of theories with less supersymmetry.
## V Acknowledgements
This research was supported in part by DOE grant DE-FG02-96ER40559.
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# Low frequency response of a collectively pinned vortex manifold
## I Introduction
In various situations the theory of weak collective pinning (for a review see Ref. ) treats a vortex medium in a type-II superconductor like an elastic manifold interacting with a random pinning potential. In particular, a vortex contribution to a low frequency response to an applied current (impedance) is ascribed to the classical (thermally activated) tunneling of large portions of the vortex manifold between different minima of the uncorrelated random potential . In the framework of this approach the pairs of states that allow for thermally activated tunneling between them are treated as current-biased two-level systems. Earlier analogous ideas were applied for the description of the dynamic response of spin glasses and of randomly pinned dislocations and interfaces .
The frequency dependence of the two-level systems contribution to specific impedance $`z_{tl}(\omega )\rho _{tl}(\omega )+i\omega l_{tl}(\omega )`$ of a type-II superconductor penetrated by external magnetic field has been found by different groups of authors to be of the form:
$`l_{tl}(\omega )`$ $``$ $`[\mathrm{ln}(1/\tau _0|\omega |)]^y;`$ ()
$`\rho _{tl}(\omega )`$ $``$ $`|\omega |[\mathrm{ln}(1/\tau _0|\omega |)]^{y1}.`$ ()
but the conclusions of these three groups on the value of the exponent $`y`$ are not compatible with each other.
Comparison shows that the discrepancy appears because the calculation of Fisher, Fisher and Huse is based on a semi-phenomenological conjecture that a contribution of an active two-level system to inductance can be estimated (from above) by replacing this two-level system with an insulating hole of the same volume. This assumption does not take into account that the motion of vortices inside of a two-level system leads to a change of a phase distribution in superconductor (and of the energy of electric current in it) also outside of the limits of this particular two-level system. Accordingly, it turns out to be in direct contradiction with the explicit expression for a two-level system contribution to impedance derived by Koshelev and Vinokur and used in their calculation of $`z_{tl}(\omega )`$, and naturally leads to a different value of $`y`$. The less pronounced descrepancy between the results of Ref. and Ref. can be explained, in particular, by the difference in assumptions on statistics of metastable states.
In the present work we revisit this problem following the more reliable (as is demonstrated below) approach of Koshelev and Vinokur . We start by deriving the value of the exponent $`y`$ for the generic case of a vortex manifold with a single elastic modulus (but with arbitrary dimensionality). Any more complex system with different (but local) elastic moduli will be characterized by the same value of $`y`$. In Ref. the value of $`y`$ in non-dispersive system has been found only for the particular case of a single vortex pinning.
We then suggest a new independent way to check the validity of the numerous assumptions involved in calculation of the two-level systems contribution to a linear dynamic response by using the same set of assumptions for calculation of another quantity (the amplitude of thermal fluctuations), which on the other hand can be calculated exactly, because in the framework of a random manifold description it has to be the same as in absence of disorder . It turns out that application of the assumptions used in the previous calculation indeed produces an answer which is in agreement with the well known result for the pure case.
The possibility of the additional check turns out to be very useful when we address the case of a thin superconducting film in which the compressibility modulus $`c_{11}`$ of vortex manifold is always non-local. It allows to draw some conclusions about size and shape distribution of two-level systems which otherwise would be unavailable. Inclusion of these conclusions into calculation leads \[for $`c_{11}(q)1/q^2`$\] to very weak (double logarithmic) frequency dependence of $`l_{tl}(\omega )`$ corresponding to $`y=0`$.
The case of a thin superconducting film is also of a special interest in relation with recent experimental investigation of $`z(\omega )`$ in ultrathin YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> films . Note that in a two-dimensional (2D) geometry with complete penetration of the magnetic field the frequency-dependent specific (sheet) impedance directly determines the response of a film to external field, whereas in the case of a bulk superconductor it has to be extracted from the surface impedance.
In all the cases we have considered $`l_{tl}(\omega )`$ diverges for $`\omega 0`$. Thus the important consequence of our results is that the random manifold approximation can not be sufficient for the description of a truely superconducting vortex glass phase with finite superfluid density.
The outline of the article can be summarized as follows. In Sec. II the random manifold problem is briefly introduced and some situations when it is appropriate for the description of a vortex manifold in a superconductor are specified. In Sec. III the statistical properties of the metastable states which have to be taken into account in the framework of the two-level system approach are discussed. In Sec. IV the application of this approach to calculation of the vortex medium contribution to impedance is presented in the systematic form for the non-dispersive case.
In Sec. V the same approach is used for the analysis of thermal fluctuations. We show that (with the same set of assumptions as used earlier in calculation of a linear dynamic response) it indeed produces an answer which is consistent with expectations based on consideration of pure system.
Sec. VI is devoted to the discussion of a thin superconducting film in which the non-locality of the compressibility modulus is always important. We show that if one assumes that the dominant two-level systems in this case are strongly anisotropic (as suggested by the energy balance estimates used in the analysis of the non-linear creep ), the expression for thermal fluctuations amplitude turns out to be convergent in contrast to its logarithmic divergence in the pure system. The only way to resolve this contradiction consists in assuming that the statistics of metastable states is dominated by the presence of hierarchical sequence of quasi-isotropic two-level systems. The same distribution is then used for the calculation of the components of $`z_{tl}(\omega )`$. In Sec. VII the results are summarized and discussed.
## II Random manifold problem
An elastic manifold (with internal dimension $`D`$) interacting with a random pinning potential can be described by the Hamiltonian:
$`H`$ $`=`$ $`H_{\mathrm{el}}+H_\mathrm{d}`$ (2)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^D𝐱_1d^D𝐱_2[G_0^1(𝐱_1𝐱_2)]^{ab}u^a(𝐱_1)u^b(𝐱_2)}`$ (4)
$`+{\displaystyle d^D𝐱v[𝐱,𝐮(𝐱)]}`$
where the $`N`$-dimensional vector $`𝐮(𝐱)u^a(𝐱)`$ is the displacement of the manifold. The first term in Eq. (4) describes (in a most general form) the elastic energy of the manifold and the second the energy of its interaction with a random pinning potential $`v(𝐱,𝐮)`$.
The simplest assumption would consist in assuming that the random potential $`v(𝐱,𝐮)`$ has a Gaussian distribution with
$`v(𝐱,𝐮)_\mathrm{d}`$ $`=`$ $`0`$ (5)
$`v(𝐱_1,𝐮_1)v(𝐱_2,𝐮_2)_\mathrm{d}`$ $`=`$ $`\delta (𝐱_1𝐱_2)w(𝐮_1𝐮_2)`$ (6)
Here and further on the angular brackets with subscript $`d`$ stand for the average over disorder, and with subscript $`th`$ for the thermal average. We discuss only the case of a short-ranged random potential correlation function $`w(𝐮)`$.
In the simplest situation (which would imply, in particular, full isotropy and absence of dispersion) the first term in Eq. (4) can be chosen in the form
$$H_{\mathrm{el}}=\frac{J}{2}d^D𝐱\left(\frac{u^a}{x^\beta }\right)^2$$
(7)
with a single elastic modulus $`J`$.
The physical systems which can be described by the Hamiltonian of the form (4) include, in particular, a domain wall in a 2D or 3D Ising type ferromagnet/antiferromagnet ($`D=1,2`$; $`N=1`$); a dislocation in a crystal ($`D=1`$, $`N=2`$); a single vortex line in a large area Josephson junction ($`D=1`$, $`N=1`$) or bulk superconductor ($`D=1`$, $`N=2`$); a vortex medium in superconducting film ($`D=2`$, $`N=2`$) or bulk superconductor ($`D=3`$, $`N=2`$), a layered superconductor with in-plane field ($`D=3`$, $`N=1`$) or a large area Josephson junction with in-plane field ($`D=2`$, $`N=1`$). In all these cases the random pinning potential is automatically provided by impurities present in any solid or/and by geometrical inhomogeneities.
Note, however, that the random manifold approximation assumes the energy of the interaction with the inhomogeneities to be uncorrelated for different displacements, whereas the energy of the interaction of an ideal vortex crystal with the inhomogeneities does not change if the vortex crystal is shifted as a whole by one lattice constant. Thus, the area of applicability of the random manifold approximation for the description of vortex crystal pinning is resticted. One can use this approach if the relevant displacements do not exceed the lattice period (Larkin regime ) or when the ordering in the vortex crystal is destroyed by the presence of defects whose motion with respect to the vortex manifold is dynamically frozen in comparison with the motion of the manifold itself.
Recent experiments of the Neuchâtel group on ultrathin YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> films have demonstrated a crossover to the regime in which the contribution to resistivity associated with the motion of point-like defects is negligible in comparison with the contribution which can be ascribed to collective pinning behaviour.
## III Low energy metastable states
The low-frequency dynamics of a weakly pinned elastic manifold can be associated with the thermally activated tunneling of large domains of the manifold between different minima of the random potential . In a simple system with a discrete spectrum only the tunneling between the ground state and the first excited state is of importance at low temperatures, thus it can be reduced to a two-level system. For an infinite manifold one should take into account that such two-level systems appear at all scales and form a hierarchical structure, i. e. if some domain of the manifold can tunnel between some states one also has to consider the possibility of tunneling of smaller domains inside this area between different pairs of states.
Each of such two-level systems can be characterized by its (linear) size $`L`$, its volume $`VL^D`$ (untill specified we discuss the simplest isotropic case with a single non-dispersive elastic modulus), the typical vortex displacement between the two states $`u`$, which, for example, can be defined by
$$u^2=\frac{1}{V}d^D𝐱𝐮^2(𝐱),$$
(8)
the energy difference between the two states $`\mathrm{\Delta }`$ and the energy barrier $`U`$ which has to be overcome for moving the manifold (the vortex bundle) from one of the two states to the other.
The universality hypothesis introduced by Ioffe and Vinokur suggests that for each length scale $`L`$ there should exist only one relevant energy scale $`E(L)`$ such that, in particular, the typical values of $`\mathrm{\Delta }`$ and $`U`$ for a system of size $`L`$ are proportional to $`E(L)`$. The magnitude of $`E(L)`$ can be then estimated by estimating the elastic energy associated with the displacement of the vortex bundle of the size $`L`$:
$$E(L)JV(L)\frac{u^2(L)}{L^2}.$$
(9)
If the scale dependence of the typical displacement $`u(L)`$ is given by $`u(L)L^\zeta `$ (where $`\zeta `$ is usually called the wandering exponent) Eq. (9) leads to $`E(L)L^\chi `$ with
$$\chi =2\zeta +D2.$$
(10)
On the other hand Fisher, Fisher and Huse have suggested that the scale dependence of the typical energy barrier $`U(L)`$ can be described by another exponent $`\psi `$ not necessariliy coinsiding with $`\chi `$ ($`\psi \chi `$). For the sake of generality in the following we will keep (for a while) the separate notation for $`\psi `$, although we will assume that the typical value of $`\mathrm{\Delta }`$ (for the given length scale $`L`$) can be estimated with the help of Eq. (9).
Various properties of the manifold depend also on the form of the size distribution function $`\nu (L)`$ of the two-level systems. The hierarchical distribution implies that any two-level system can include smaller two-level systems whose size differ from that of the ”parent” two-level system by some numerical factor of the order of one. Thus the ratio of the typical ”neighboring” length scales has to be more or less constant across the whole length-scale range involved . This is compatible with a uniform distribution of the logarithms of the length scales. However, for each length scale one should also include the factor $`1/V(L)`$ proportional to the largest possible concentration of non-overlapping two-level systems of size $`L`$. Thus $`\nu (L)`$ has to be of the form
$$dL\nu (L)\frac{dL}{L}\frac{1}{V(L)}$$
(11)
Koshelev and Vinokur have introduced the first factor (imposed by the hierarchical structure) in the r.h.s. of Eq. (11) in the form $`dU/U`$ (without any explanation), which for $`U(L)`$ algebraically dependent on $`L`$ is equivalent to $`dL/L`$, whereas in Ref. the hierarchical nature of the size distribution of two-level systems has not been taken into account.
## IV Two-level systems contribution to impedance
The contribution of the two-level systems to the specific impedance of a superconductor is given by
$$z_{tl}(\omega )=n\frac{\gamma }{T}\frac{(V\overline{u})^2}{\mathrm{cosh}^2(\mathrm{\Delta }/2T)}\frac{i\omega }{1+i\tau \omega }_\mathrm{d};\gamma =\frac{B^2}{4c^2}$$
(12)
where $`n`$ is the concentration of such systems, $`B`$ is the magnetic induction and $`T`$ is the temperature, whereas $`V`$, $`\overline{u}`$, $`\mathrm{\Delta }`$ and $`\tau `$ are parameters characterizing a particular two-level system: $`V`$ is the volume in which the vortices are displaced (or the area in the case of a 2D superconductor), $`\overline{u}`$ is the average displacement of the vortices inside the bundle in the direction of current-induced force and $`\mathrm{\Delta }`$ is the difference in energy between the two states. The relaxation time $`\tau `$ describing the rate of the thermally activated (incoherent) tunneling between the two states depends exponentially:
$$\tau =\tau _0\mathrm{exp}(U/T)$$
(13)
on the barrier $`U`$ separating them.
If one splits $`z_{tl}(\omega )`$ into real and imaginary parts:
$$z_{tl}(\omega )=\rho _{tl}(\omega )+i\omega l_{tl}(\omega )$$
(14)
and assumes that the average over disorder can be estimated by taking for each length scale the typical (scale-dependent) values of all the parameters involved, the expression for the two-level systems contribution to specific impedance $`l_{tl}(\omega )`$ is reduced to
$$l_{tl}(\omega )\frac{\gamma }{T}_{L_c}^{L_\omega }𝑑L\nu (L)V^2(L)\overline{u}^2(L)\mathrm{cosh}^2(\mathrm{\Delta }/2T)_\mathrm{d}(L)$$
(15)
Due to the exponentially fast increase of $`\tau `$ with $`L`$, instead of including in Eq. (15) the factor
$$\frac{1}{1+[\tau (L)\omega ]^2}$$
(16)
the integration in it is cut off (at the upper limit) at a frequency-dependent length scale $`L_\omega `$ defined by the relation $`\tau (L_\omega )\omega 1`$. For $`U(L)ϵ(L/L_c)^\psi `$
$$L_\omega L_c\left(\frac{T}{ϵ}\mathrm{ln}\frac{1}{\tau _0|\omega |}\right)^{1/\psi }$$
(17)
The integration interval in Eq. (15) is limited from below by the (temperature dependent) collective pinning length $`L_c`$ which determines the boundary between the different regimes of fluctuations. At length scales lower than $`L_c`$ the manifold can be considered as fluctuating within one of the minima of the (thermally renormalized) random potential, whereas for larger scales only the jumps between different valleys of the potential are of importance. The contribution to $`l(\omega )`$ from length scales smaller than $`L_c`$ has a finite limit for $`\omega 0`$.
Since $`\mathrm{\Delta }`$ is the difference in energy between the two spatially separated states in the uncorrelated random potential, the distribution function $`p(\mathrm{\Delta })`$ can be expected to remain finite for $`\mathrm{\Delta }0`$ (the broad distribution assumption ). The last factor in the r.h.s. of Eq. (15) (the fraction of ”thermally active” two-level systems) can be then estimated as
$$\mathrm{cosh}^2(\mathrm{\Delta }/2T)_\mathrm{d}(L)\frac{T}{\mathrm{\Delta }(L)}$$
(18)
where $`\mathrm{\Delta }(L)`$ is the typical value of $`\mathrm{\Delta }`$ for the length scale $`L`$ \[for example the width of $`p(\mathrm{\Delta })`$\]. Note that (for any scale) only small fraction of two-level systems is assumed to be not frozen and therefore involved in linear dynamic response (or thermal fluctuations). Thus they are expected to be well separated from each other, which justifies neglecting their interaction.
Substitution of Eq. (18) into Eq. (15) leads to
$$l_{tl}(\omega )\gamma _{L_c}^{L_\omega }𝑑L\nu (L)\frac{V^2(L)\overline{u}^2(L)}{\mathrm{\Delta }(L)}$$
(19)
According to the universality hypothesis $`\mathrm{\Delta }(L)`$ has to be of the same order of magnitude as the elastic contribution to energy estimated in Eq. (9), which for $`\overline{u}u`$ gives:
$$\frac{V^2(L)\overline{u}^2(L)}{\mathrm{\Delta }(L)}V(L)\frac{L^2}{J}$$
(20)
Substitution of Eqs. (11) and (20) into Eq. (19) then leads to
$$l_{tl}(\omega )\frac{\gamma }{J}_{L_c}^{L_\omega }\frac{dL}{L}L^2\frac{\gamma }{J}L_c^2\left(\frac{T}{ϵ}\mathrm{ln}\frac{1}{\tau _0|\omega |}\right)^y,$$
(21)
where $`y=2/\psi `$.
With the same assumptions that have been used for the derivation of Eq. (21) the two-level system contribution to the resistivity is given by
$$\rho _{tl}(\omega )\frac{\gamma }{J}_{L_c}^{\mathrm{}}𝑑LL\frac{\tau (L)\omega ^2}{1+[\tau (L)\omega ]^2}$$
(22)
Alternatively $`\rho _{tl}(\omega )`$ can be restored from $`l_{tl}(\omega )`$ with the help of the simplified form of the Kramers-Kronig relation:
$$\rho _{tl}(\omega )|\omega |\frac{\pi }{2}\frac{d}{d\mathrm{ln}|\omega |}l_{tl}(\omega )$$
(23)
which is applicable for $`l_{tl}(\omega )=f(\mathrm{ln}|\omega |)`$. Both methods give
$$\rho _{tl}(\omega )\frac{\gamma }{J}\left(\frac{T}{ϵ}\right)^y|\omega |\left(\mathrm{ln}\frac{1}{\tau _0|\omega |}\right)^{y1}$$
(24)
It can be shown that $`y`$ is equal to $`2/\psi `$ not only for the simplest case of a single elastic modulus, but for the general case of non-dispersive moduli. In a bulk superconductor at large enough scales (which corresponds to low enough frequencies) all elastic moduli become local. The case of a thin film in which the strong dispersion of the compressibility modulus is unavoidable is considered in Sec. VI.
For finite values of the current density $`j`$ Eq. (12), which has been the starting point of our calculation, is applicable only if $`V\overline{u}B/c`$ is small in comparison with temperature. This determines the current dependence of the length scale $`L_j(T/j)^{\frac{1}{D+\zeta }}`$ at which the integration in Eq. (15) should be cut off if $`L_jL_\omega `$. In that case the growth of $`l_{tl}(\omega )`$ with decreasing $`\omega `$ saturates at $`l_{tl}(\omega =0,j)(T/j)^{\frac{2}{D+\zeta }}`$.
## V Comparison of two approaches to calculation of thermal fluctuations amplitude
In the present work we suggest an independent way to check the validity of the two-level system approach for the description of the linear dynamic response of a collectively pinned manifold. This can be done because in the random manifold problem the static irreducible correlation function
$`u^a(𝐱_1)u^b(𝐱_2)`$ $``$ $`[u^a(𝐱_1)u^a(𝐱_1)_{\mathrm{th}}]`$ (26)
$`\times [u^b(𝐱_2)u^b(𝐱_2)_{\mathrm{th}}]_{\mathrm{th}}_\mathrm{d}`$
(which can be associated with thermal fluctuations) according to Shultz et al should be exactly the same as in the absence of disorder:
$$u^a(𝐱_1)u^b(𝐱_2)=TG_0^{ab}(𝐱_1𝐱_2).$$
(27)
(a brief derivation can be found in Appendix). An analogous relation for the case of periodic behaviour of $`w(𝐮)`$ with respect to displacement has been suggested by Dotsenko and Feigel’man .
In the presence of disorder the long-distance behaviour of the correlation function (26) has to be mediated by the two-level systems. Therefore, the investigation of thermal fluctuations in terms of the two-level system approach and comparison of the result with well known result for the pure system allows to check the validity of different conjectures involved in the calculation of a linear dynamic response. Instead of considering the dependence of $`u^a(𝐱_1)u^b(𝐱_2)`$ on $`|𝐱_1𝐱_2|`$ one can alternatively investigate the dependence of $`u^2`$ on the size of the system $`L_0`$, which also has to be the same as in the pure case.
For a single two-level system with the energy gap $`\mathrm{\Delta }`$ the amplitude of the thermal fluctuations of the displacement is given by
$$u^2(uu_{\mathrm{th}})^2_{\mathrm{th}}=\frac{(u_1u_2)^2}{4\mathrm{cosh}^2(\mathrm{\Delta }/2T)}$$
(28)
For a collectively pinned manifold the dominant large-scale contribution to $`u^2(𝐱)`$ should come from the two-level systems which include the point $`𝐱`$ and \[on the same assumptions as have been used while calculating $`l_{tl}(\omega )`$ and $`\rho _{tl}(\omega )`$\] can be estimated as
$$u_{tl}^2T_{L_c}^{L_0}𝑑L\nu (L)\frac{V(L)u^2(L)}{\mathrm{\Delta }(L)}$$
(29)
where the upper limit of integration is now imposed by the size of the system and \[in accordance with Eq. (18)\] $`\mathrm{cosh}^2(\mathrm{\Delta }/2T)_\mathrm{d}`$ has been already replaced with $`T/\mathrm{\Delta }(L)`$.
Substitution of Eqs. (9) and (11) into Eq. (29) then gives
$$u_{tl}^2\frac{T}{J}_{L_c}^{L_0}\frac{dL}{L^{D1}}$$
(30)
which, if compared with the trivial result for the case without disorder (for which $`u^2u^2_{\mathrm{th}}`$):
$$u^2\frac{T}{J}_{|q^\alpha |>\pi /L_0}\frac{d^D𝐪}{(2\pi )^D}\frac{1}{q^2}$$
(31)
reproduces all its important features. Namely, the r.h.s. of Eq. (30) (i) contains the correct prefactor $`T/J`$; (ii) demonstrates the correct dependence on the size of the system:
$$u^2\{\begin{array}{cc}L_0^{2D}\hfill & \hfill D<2\\ \mathrm{ln}(L_0)\hfill & \hfill D=2\\ L_c^{2D}L_0^{2D}\hfill & \hfill D>2\end{array}$$
(32)
and (iii) the system-size dependent contribution to $`u_{tl}^2`$ does not depend on the unknown disorder-related parameters $`L_c`$ and $`ϵ`$.
This allows to conclude that different assumptions which have been used while calculating $`l_{tl}(\omega )`$, $`\rho _{tl}(\omega )`$ and $`u_{tl}^2`$ \[the universality hypothesis, the hierarchical distribution of two-level systems, the broad distribution assumption for $`p(\mathrm{\Delta })`$\] were indeed chosen in a reasonable way.
## VI Thin superconducting film in perpendicular magnetic field
The long-range interaction of vortices makes the compressibility modulus $`c_{11}`$ of a bulk superconductor strongly non-local for the wave-lengths smaller than the magnetic field penetration depth $`\lambda `$. In a thin superconducting film the penetration depth $`\mathrm{\Lambda }`$ is strongly increased in comparison with that of a bulk superconductor: $`\mathrm{\Lambda }=2\lambda ^2/d`$ , where $`d\lambda `$ is the thickness of the film. Therefore, in a thin film the dependence $`c_{11}(q)\overline{c}_{11}/q^2`$ (where $`\overline{c}_{11}B^2/2\pi \mathrm{\Lambda }c_{66}a^2`$) resulting from a non-screened vortex-vortex interaction holds in a much wider range of length scales than in a bulk superconductor. Here $`c_{66}\mathrm{\Phi }_0B/32\pi ^2\mathrm{\Lambda }`$ is the shear modulus of the film (which in contrast to $`c_{11}`$ is always local), $`\mathrm{\Phi }_0=hc/2e`$ is the flux quantum and $`a^2=B/\mathrm{\Phi }_0`$ is the vortex density.
If one tries to shift a vortex bundle in such a system (in search of the next potential minimum), it turns out that the optimal shape of the bundle is strongly anisotropic with the size in the direction of the displacement $`L`$ much larger than the size in the perpendicular direction $`L_{}`$. The optimal relation between $`L`$ and $`L_{}`$ can be found by minimizing the total elastic energy for the given area of a bundle $`SLL_{}`$. Minimization (for fixed $`S`$) of $`E_{\mathrm{com}}+E_{\mathrm{sh}}`$ where
$$E_{\mathrm{com}}\overline{c}_{11}S^2\left(\frac{\overline{u}}{L}\right)^2$$
(33)
and
$$E_{\mathrm{sh}}c_{66}S\left(\frac{u}{L_{}}\right)^2$$
(34)
or a simple comparison of $`E_{\mathrm{com}}`$ with $`E_{\mathrm{sh}}`$ for $`\overline{u}u`$ gives
$$L\frac{L_{}^3}{a^2}L_{}.$$
(35)
Note, however, that when the energy scale defined by Eq. (34) is used as an estimate for $`\mathrm{\Delta }`$, the factor $`Su^2/\mathrm{\Delta }`$ in the 2D version of Eq. (29) is reduced to
$$\frac{Su^2}{\mathrm{\Delta }}\frac{L_{}^2}{c_{66}}$$
(36)
for arbitrary relation between $`L`$ and $`L_{}`$.
The calculation of Sec. V has confirmed that the form of the size distribution of two-level systems can be correctly estimated by assuming that they do not overlap with each other, but can be situated inside each other forming hierarchical structures. Such estimate is consistent with the conjecture that the number of metastable states to which a particular domain of a manifold can tunnel (without moving a much larger part of the manifold) is always of the order of one . In what follows we assume that the same property holds also in presence of dispersion.
The most optimistic estimate for $`\nu (L)`$ can be then obtained by assuming that for all scales the strongly anisotropic two-level systems are arranged in the most advantageous way to cover all the area available, which corresponds to
$$dL\nu (L)\frac{dL}{L}\frac{1}{LL_{}}.$$
(37)
The more realistic estimate should probably take into account that independent anisotropic two-level systems are likely to have uncorrelated orientations, so the requirement of non-overlapping will lead to $`\nu (L)L^3`$.
After substitution of Eq. (36) into the integral \[of the form (29)\] defining $`u_{tl}^2`$ one obtains that for $`LL_{}^3`$ it is convergent at the upper limit:
$$u_{tl}^2\frac{T}{c_{66}}_{L_c}^{\mathrm{}}\frac{dL}{L}\left(\frac{L_{}}{L}\right)<\mathrm{}$$
(38)
even for the optimistic form of $`\nu (L)`$ given by Eq. (37). On the other hand, we know that in a pure system the contribution of the transverse modes (which depends only on the shear modulus which is local) leads to the logarithmic divergence of $`u^2u^2_{\mathrm{th}}`$. It follows from the results of Schulz et al that in the presence of disorder the same behaviour should be mediated by the large-scale two-level systems.
A plausible way to explain the logarithmic divergence of $`u_{tl}^2`$ consists in assuming that the film should contain not only anisotropic two-level systems with $`LL_{}`$, but also an hierarchical sequence of quasi-isotropic two-level systems in which $`L_{}`$ is of the same order as $`L`$. For such two-level systems the requirement of the balance between the different contributions to the elastic energy ($`E_{\mathrm{com}}E_{\mathrm{sh}}`$) leads to $`\overline{u}(a/L)uu`$, which means that the displacement of the vortices in quasi-isotropic bundles is mostly of rotational type.
The linear dynamic response has to be associated with the same degrees of freedom as are taken into account in the calculation of the static thermal fluctuations. However, $`z_{tl}(\omega )`$ can not be obtained by direct application of the fluctuation-dissipation theorem to $`u_{tl}^2`$, since these two quantities include different linear combinations of the degrees of freedom involved. Nonetheless, when calculating $`z_{tl}(\omega )`$ one should take into account the same set of two-level systems as for the calculation of $`u_{tl}^2`$, in contrast to the case of the non-linear creep for which the shape of the moving vortex bundles is imposed by the applied current .
Application of the expression (33) for $`E_{\mathrm{com}}`$ as an estimate for $`\mathrm{\Delta }`$ shows that the factor $`S^2\overline{u}^2/\mathrm{\Delta }`$ in the 2D version of Eq. (19) does not depend on $`L_{}`$ and can be estimanted as $`L^2/\overline{c}_{11}L^2\mathrm{\Lambda }/B^2`$. For the hierarchical sequence of quasi-isotropic two-level systems with $`\nu (L)L^3`$ this leads to the extremely weak (double logarithmic) frequency dependence of
$$l_{tl}(\omega )\frac{\gamma \mathrm{\Lambda }}{B^2}_{L_c}^{L_\omega }\frac{dL}{L}\frac{\mathrm{\Lambda }}{c^2}\mathrm{ln}\left(\frac{T}{ϵ}\mathrm{ln}\frac{1}{\tau _0|\omega |}\right),$$
(39)
which can hardly be expected to be resolvable from the background superfluid contribution $`l_0`$ in the experiments probing the low-frequency response of thin films.
However, substitution of Eq. (39) into Eq. (23) gives
$$\rho _{tl}(\omega )\frac{\mathrm{\Lambda }}{c^2}\frac{|\omega |}{\mathrm{ln}(1/\tau _0|\omega |)}$$
(40)
which, in contrast to Eq. (39), does not exhibit any specially weak dependence on $`\omega `$. Note that two unknown disorder-related parameters $`L_c`$ and $`ϵ`$ as well as the magnetic field dependence have dropped out from Eq. (40).
The presence of the hierarchical sequence of quasi-isotropic two-level systems still leaves enough place for more optimal anisotropic two-level systems with $`L_{}L`$. Although they do not contribute much to $`u_{tl}^2`$, their contributions to the components of $`z_{tl}(\omega )`$ could be of importance. However, their size distribution will be forced by the presence of hierarchical sequence of quasi-isotropic two-level systems to be of the same form $`\nu (L)L^3`$ and, therefore, their contribution to $`l_{tl}(\omega )`$ and $`\rho _{tl}(\omega )`$ will be of the same form as given by Eqs. (39)-(40).
In a thin superconducting film the compressibility modulus $`c_{11}`$ is non-local not only for $`\mathrm{\Lambda }q1`$ where $`c_{11}(q)\overline{c}_{11}/q^2`$, but also for $`\mathrm{\Lambda }q1`$ where $`c_{11}(q)B^2/2\pi q`$. An analogous calculation for such form of $`c_{11}(q)`$ produces for the components of $`z_{tl}(\omega )`$ the answers of the form (1) with $`y=1/\psi `$.
## VII Conclusion
In the present work we argue that thermal fluctuations of a collectively pinned vortex manifold are determined by the same degrees of freedom (related to thermally activated tunneling between the pairs of low-lying metastable states - two-level systems) as its low-frequency linear dynamic response. Therefore one can use the known dependence of thermal fluctuations amplitude on the size of the system (which has to be exactly the same as in absence of disorder ) for checking the consistency of the conjectures which are used in the calculation of the linear dynamic response. The set of the assumptions which are necessary to produce the correct answer for the thermal fluctuations amplitude includes, in particular, the conjecture on hierarchical distribution of two-level systems (which means that they can be situated inside each other) and also the universality hypothesis. If the same set of assumptions is used for the calculation of a vortex manifold contribution to impedance its frequency dependence (in absence of dispersion) is given by Eqs. (()I) with $`y=2/\psi `$, where $`\psi `$ is the exponent (which depends both on $`D`$ and $`N`$) describing the scale dependence of the typical energy barrier $`U(L)`$.
The same result is also applicable in the limit af small fields, when one can neglect the interaction between different vortices and treat each vortex separately as 1D manifold. In that case one should take the value of $`\psi `$ corresponding to $`D=1`$.
Note that in the framework of our analysis the universality hypothesis has been used only for the estimate of $`\mathrm{\Delta }(L)`$. If (as suggested by Ioffe and Vinokur ) it is further assumed that the same energy scale can be used for the estimate of $`U(L)`$, the value of $`\psi `$ will coinside with $`\chi `$ given by Eq. (10). Different approaches including scaling arguments , functional renormalization group and a self-consistent calculation incorporating replica symmetry breaking lead to
$$\zeta =\frac{4D}{4+\beta N}$$
(41)
with $`1/2\beta 1`$. For the case of thin superconducting film (D=2, N=2) or bulk superconductor (D=3, N=2) Eqs. (10) and (41) give the values of $`\chi `$ in the interval from $`2/3`$ to $`7/5`$, that is around $`1`$.
Koshelev and Vinokur have found the value of the exponent $`y`$ only for the particular values of $`\zeta `$ and $`\chi `$ corresponding to $`D=1`$ and $`N=2`$ and not in the general form as above. In Ref. the hierarchical nature of the size distribution (of the two-level systems) has not been taken into account and the estimate of $`\rho _{tl}(\omega )`$ has been obtained without integrating over the scales, which has led (in our notation) to $`y=2/\psi +1`$. As has been already mentioned in the Introduction the analogous calculation in Ref. has been performed using the assumption which is in contradiction with Eq. (12) and therefore can not be used for comparison.
Note that $`l_{tl}(\omega )`$ diverges in the limit of $`\omega 0`$, which corresponds to supression of superfluid density \[inversely proportional to $`lim_{\omega 0}l(\omega )`$\]. Thus the results of this work are not in agreement with the popular point of view that random manifold approach provides an exhaustive description of dynamic properties of a truely superconducting vortex glass phase (which is supposed to be formed at low temperatures due to pinning ), at least if superconductivity is understood as the ability to carry a superconducting (non-dissipative) current and not only as the vanishing of the linear resistance. Our analysis suggests that in the framework of random manifold approach the finite value of $`l_{tl}(\omega 0)`$ is incompatible with the correct scale dependence of $`u_{tl}^2`$. Therefore one has to conclude that the accurate description of a vortex glass phase which can carry a superconducting current (if such phase exists at all) requires a more sophisticated treatment than the description of vortex medium in terms of an elastic manifold interacting with a random potential. For example, it possibly should take into account that some of the defects of a vortex lattice are generated by a disorder and can not freely move with a vortex manifold.
However, both in the case of Larkin regime and in the case of dynamically frozen thermally excited defects (frozen vortex liquid regime) the frequency range of the applicability of the random manifold description of a vortex medium in a superconductor is anyway restricted from below. Moreover, the two-level system contribution to impedance $`l_{tl}(\omega )`$ is only logarithmic in $`\omega `$ and in a practical situation may be negligible in comparison with ”bare” impedance $`l_0`$ down to exponentially low frequencies.
On the other hand $`\rho _{tl}(\omega )`$ produces in the low frequency limit the dominant contribution to the resistivity in comparison with the contributions related with the normal channel conductance and with the oscillations of manifold within each minimum of a random potential, both of which at low frequencies are proportional to $`\omega ^2`$.
In thin superconducting films the compressibility modulus of vortex manifold is always non-local. This leads to the strong anisotropy of the vortex bundles participating in the non-linear creep . However our analysis has shown that the correct length-scale dependence of thermal fluctuations amplitude requires the presence of a hierarchical sequence of quasi-isotropic two-level systems. A linear dynamic response (which has to be calculated assuming that the applied current does not change the properties of the system) produced by the same set of the two-level systems corresponds to $`y=1/\psi `$ for $`c_{11}(q)1/q(\mathrm{\Lambda }q1)`$ and to $`y=0`$ \[with $`l_{tl}(\omega )`$ still diverging at $`\omega 0`$ but only as a double logarithm\] for $`c_{11}(q)1/q^2(\mathrm{\Lambda }q1)`$. The contribution from the more optimal anisotropic two-level systems can be expected to be of the same form. In both regimes \[$`c_{11}(q)1/q`$ and $`c_{11}(q)1/q^2`$\] the value of the magnetic field $`B`$ drops out from the expression for $`z_{tl}(\omega )`$.
Although in thin films the dc resistivity is always finite due to thermally activated motion of point-like defects of vortex lattice (vacancies, interstitials, dislocation pairs) , the collective pinning behaviour has been found to be accessible to experimental observation in the range of frequencies/temperatures where the activated contribution to resistivity is too small.
## ACKNOWLEDGMENTS
This work has been supported in part by the Program ”Scientific Schools of the Russian Federation” (grant No. 00-15-96747) and by the Swiss National Science Foundation. The author is grateful to G. Blatter, V. B. Geshkenbein and P. Martinoli for interesting discussions and useful comments.
##
In the case of a random manifold described by Eqs. (4)-(6) it is possible to show that the irreducible correlation function of the form (26) remains the same as in absence of disorder . To this end one can express this correlation function through the second derivative
$$u^a(𝐱_1)u^b(𝐱_2)=T\frac{^2\stackrel{~}{F}}{s_1^as_2^b}|_{s_1^a=s_2^a=0}$$
(42)
of the (disorder-averaged) free energy
$$\stackrel{~}{F}=T\mathrm{ln}\left[𝑑𝐮\mathrm{exp}\left(\frac{\stackrel{~}{H}}{T}\right)\right]_\mathrm{d}$$
(43)
with respect to the coefficients in the auxilary (source) terms added to the Hamiltonian
$$\stackrel{~}{H}=H+s_1^au^a(𝐱_1)s_2^au^a(𝐱_2).$$
(44)
In order to calculate $`\stackrel{~}{F}`$ it is convenient to shift the variables $`𝐮(𝐱)`$ (over which the integration in the partition function is performed) according to
$$u^a(𝐱)u_{}^a(𝐱)u^a(𝐱)+G_0^{ab}(𝐱𝐱_1)s_1^bG_0^{ab}(𝐱𝐱_2)s_2^b,$$
(45)
which allows to split the non-random contribution to $`\stackrel{~}{H}\{u^a\}`$ into two terms: $`H_{\mathrm{el}}\{u_{}^a\}+E(s_1^a,s_2^a)`$, the first of which does not depend on $`s_{1,2}^a`$ and the second
$`E(s_1^a,s_2^a)`$ $`=`$ $`{\displaystyle \frac{1}{2}}G_0^{ab}(0)s_1^as_1^b+G_0^{ab}(𝐱_1𝐱_2)s_1^as_2^b`$ (47)
$`{\displaystyle \frac{1}{2}}G_0^{ab}(0)s_2^as_2^b`$
does not depend on $`u_{}^a(𝐱)`$.
On the other hand the distribution function of the random potential $`v(𝐱,𝐮)`$ is not affected by the shift defined by Eq. (45). Therefore the free energy defined by Eq. (43) differs from its value for original problem (that is for $`s_1^a=s_2^a=0`$) only by addition of the term $`E(s_1^a,s_2^a)`$, differentiation of which leads to Eq. (27).
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# The X-ray Faint Early-Type Galaxy NGC 4697
## 1. Introduction
It is well-established that early-type galaxies exhibit a large range of X-ray–to–optical luminosity ratios, $`L_X/L_B`$. Although there is a strong correlation between the X-ray and optical luminosities of elliptical and S0 galaxies ($`L_XL_B^{1.73.0}`$; Canizares, Fabbiano, & Trinchieri 1987; White & Davis 1997; Brown & Bregman 1998), a large dispersion exists in this relation. Two galaxies with similar blue luminosities might have X-ray luminosities that differ by as much as a factor of 100 (Canizares et al. 1987; Fabbiano, Kim, & Trinchieri 1992; Brown & Bregman 1998).
Hot ($`10^7`$ K) gas is believed to be responsible for the bulk of the X-ray emission in high $`L_X/L_B`$ galaxies (e.g., Forman, Jones & Tucker 1985). Another possible source of X-ray emission in early-type galaxies is from discrete stellar sources, primarily low-mass X-ray binaries (LMXBs). Estimates for the stellar $`L_X/L_B`$ value are uncertain and vary by a factor of ten among various studies (Forman et al. 1985; Canizares et al. 1987). Although stellar X-ray sources are not expected to contribute significantly to the 0.1–2.0 keV X-ray emission in gas-rich, X-ray bright galaxies, this claim cannot be made with any certainty for X-ray faint (low $`L_X/L_B`$) galaxies. In these galaxies, it is possible that stellar X-ray sources are the dominant X-ray emission mechanism. For this to occur, the gas lost from stellar mass loss must be removed from the galaxy by galactic winds, by ram pressure stripping from ambient intracluster or intragroup gas, or possibly both (see, e.g., Ciotti et al. 1991; Mathews & Brighenti 1998).
Evidence is mounting that LMXB emission is detectable in nearly all early-type galaxies. Matsumoto et al. (1997) found a hard ($``$5–10 keV) spectral component in 11 of 12 early-type galaxies observed with ASCA, which scaled roughly with the optical luminosity of the galaxy (note, however, that Buote & Fabian (1998) found no formal need for a hard component in the X-ray brightest galaxies of their ASCA sample). The luminosity of this hard component was small compared to the soft, gaseous component in the X-ray brightest galaxies, but increased in importance as $`L_X/L_B`$ decreased. However, the ASCA observations did not resolve this component into discrete sources because of the insufficient spatial resolution of the instrument.
Previous studies have found that the spectra of the X-ray faintest early-type galaxies are significantly different than the spectra of gas-rich X-ray bright early-type galaxies. Einstein data revealed that the X-ray faintest galaxies exhibited a strong very soft X-ray excess (Kim et al. 1992). Subsequent ROSAT PSPC and ASCA studies found that the spectra of these galaxies were best fit with a two component model consisting of the previously mentioned hard $``$5–10 keV component (generally attributed to the integrated LMXB emission), and a very soft $``$0.3 keV component, whose origin is uncertain (Fabbiano, Kim, & Trinchieri 1994; Pellegrini 1994; Kim et al. 1996).
One recently proposed solution to the very soft X-ray excess problem in X-ray faint galaxies is that it might result from the very same collection of LMXBs responsible for the hard emission (Irwin & Sarazin 1998a,b). Little is known about the very soft X-ray properties of LMXBs since nearly all Galactic examples lie in directions of high Galactic hydrogen column densities, so their soft X-ray emission is heavily absorbed. However, two nearby Galactic LMXBs that lie in directions of low Galactic hydrogen column density (Her X-1 and MS1603+2600) exhibit significant very soft X-ray emission (Vrtilek et al. 1994; Hakala et al. 1998). But perhaps the strongest evidence comes from the bulge of M31. M31 is close enough that most ($``$75%) of its bulge X-ray emission was resolved into point sources with the ROSAT PSPC and HRI (Supper et al. 1997; Primini, Forman, & Jones 1993), a majority of which are probably LMXBs. Both individually and cumulatively, these point sources in the bulge of M31 have X-ray spectral properties very similar to the integrated emission from X-ray faint early-type galaxies in the ROSAT band (Irwin & Sarazin 1998a,b). A joint ROSAT PSPC + ASCA study found the X-ray spectrum of the bulge of M31 to be almost identical to that of the X-ray faint early-type galaxy NGC 4382 over the 0.2–10 keV band (Irwin & Bregman 1999a; Kim et al. 1996). The bulge of the Sa galaxy NGC 1291 also exhibits a very soft component and has X-ray colors identical to the bulge of M31, although it is unresolved with ROSAT because of its distance. In both bulges, the measured $`L_X/L_B`$ value is comparable to those of the X-ray faintest early-type galaxies. This suggests that in addition to providing the required spectral X-ray characteristics, LMXBs are luminous and/or numerous enough to produce the required amount of X-ray emission in the X-ray faintest early-type galaxies, as well as Sa spiral bulges. No additional very soft X-ray emission source seems to be required.
In this Paper, we analyze long ROSAT PSPC, ROSAT HRI, and ASCA GIS and SIS observations of the X-ray faint elliptical galaxy NGC 4697. At a distance of 15.9 Mpc (Faber et. al. 1989; assuming a Hubble constant of 50 km s<sup>-1</sup> Mpc<sup>-1</sup>), NGC 4697 is one of the closest normal early-type galaxies. Previous work has indicated that this galaxy has a below average but not extremely low $`L_X/L_B`$ value, and ROSAT band colors typical of those of other X-ray faint early-type galaxies (Irwin & Sarazin 1998b). Jointly fitting the ROSAT PSPC and ASCA spectra allows us to place useful constraints on both the hard and soft emission, something that is not possible using each instrument separately. In addition, NGC 4697 was one of the very few X-ray faint early-type galaxies observed at length with the ROSAT HRI. We analyze this HRI image of NGC 4697 to resolve as much of the X-ray emission as possible into discrete sources. We also present a Chandra simulation of what the X-ray emission of NGC 4697 might look like.
## 2. ROSAT PSPC and ASCA Data Reduction
From the HEASARC archive we have extracted long ROSAT PSPC (RP600262A02) and ASCA GIS and SIS (62014000) observations of NGC 4697. The PSPC data were filtered such that all data with a Master Veto Rate below 170 counts s<sup>-1</sup> were excluded from the data, yielding an observation time of 36,856 seconds. The ASCA data were screened using the standard screening criteria applied to all the archival data (Revision 2 processing). See Ohashi et al. (1996) and Makishima et al. (1996) for a description of the ASCA instruments. The total GIS and SIS exposure times were 56,698 seconds and 42,138 seconds, respectively. The SIS observation was taken in 2-CCD mode.
Spectra from a 4 radius were extracted from the PSPC, GIS, and SIS data. An extraction area of this size provided a good compromise between the minimum suggested extraction region for SIS data and achieving a reasonable signal-to-noise ratio for the PSPC data. For the PSPC we have chosen background from a source-free annular region $`30^{}40^{}`$ in extent corrected for vignetting. For the ASCA data, background was extracted from the deep blank sky data provided by the ASCA Guest Guest Observer Facility. We used the same region filter to extract the background as we did the data, so that both background and data were affected by the detector response in the same manner.
For the PSPC data, only energy channels between 0.2–2.0 keV were included in the fits, and for the ASCA data we used energy channels between 0.8–10 keV. The energy channels were regrouped so that each channel contained at least 25 counts so that the $`\chi ^2`$-test is a valid indicator of goodness of fit. For all the joint fits, we linked the normalizations of the PSPC and GIS, but let the SIS normalization vary to account for the fact that some of the emission that fell on interchip boundaries of the SIS was lost.
## 3. Spectral Modeling and Luminosities of the Global Spectrum
Using XSPEC Version 10.0, we first attempted to fit a single component thermal model to the PSPC and ASCA data separately. When the PSPC spectrum was fit alone with a MEKAL model with a variable absorption component, a good fit was obtained ($`\chi _\nu ^2=1.07/58`$ degrees of freedom) for $`kT=0.430.56`$ keV and $`Z<0.02`$ (90% confidence level), in good agreement with the analysis of the same data by Davis & White (1996). A significantly different result was found for this same model when just the ASCA spectra were analyzed; the best-fit values were $`kT=2.23.9`$ keV and $`Z<0.20`$, with $`\chi _\nu ^2=1.44/53`$ degrees of freedom. Clearly, significantly different conclusions would be drawn from the spectra if data from only one of the satellites were analyzed.
This result illustrates the necessity of using both ROSAT PSPC and ASCA data when analyzing the spectrum of X-ray faint early-type galaxies. When a one-component MEKAL model with a variable absorption component was fit to the joint PSPC+ASCA spectrum, a poor fit was obtained ($`\chi _\nu ^2=2.08/115`$ degrees of freedom) with $`kT=0.82`$ keV and $`Z<0.006`$. Following the work of Kim et al. (1996) and Irwin & Bregman (1999a), we added a bremsstrahlung model in the fit. The MEKAL+bremsstrahlung model provided an excellent fit to the data ($`\chi _\nu ^2=0.94/113`$ degrees of freedom) with the absorption value fixed at the Galactic line-of-sight value of $`2.12\times 10^{20}`$ cm<sup>-2</sup> (Stark et al. 1992). The MEKAL model had parameter values of $`kT_{MEKAL}=0.26_{0.03}^{+0.04}`$ keV and $`Z=0.07_{0.03}^{+0.05}`$ (all errors are 90% confidence levels for one interesting parameter). The best-fit bremsstrahlung temperature was $`kT_{BREM}=5.2_{1.6}^{+3.0}`$ keV. Freeing the absorption parameter led to only a slight improvement in the fit, and the 90% confidence value on $`N_H`$ was consistent with the Galactic value.
We hesitate to attach a physical significance to this MEKAL+bremsstrahlung model. The spectra of individual Galactic LMXBs are often fit by a disk-blackbody+blackbody (DBB+BB) spectrum (see, e.g., Mitsuda et al. 1984). The MEKAL+bremsstrahlung, however, provided a good fit to the X-ray faint early-type galaxy NGC 4382 (Kim et al. 1996) and the bulge of M31 (Irwin & Bregman 1999a; Trinchieri et al. 1999 showed that a Raymond-Smith+bremsstrahlung model adequately fit the BeppoSAX LECS spectrum of the bulge of M31), whereas the DBB+BB model did not provide a good fit to the ROSAT PSPC + ASCA spectrum of the bulge of M31. Since we are interested in describing the spectrum of the a collection of LMXBs in the context of an elliptical galaxy and not of LMXBs on an individual basis, we will continue to use a MEKAL+bremsstrahlung model as has been done in previous studies of X-ray emission from early-type systems. We do not attempt to justify the physical relevance of the MEKAL+bremsstrahlung model, but use it here to quantify the strength of the soft and hard components of the X-ray emission.
Assuming a distance of 15.9 Mpc, the 0.25–10 keV luminosities of the soft and and hard components were $`1.50\times 10^{40}`$ ergs s<sup>-1</sup> and $`1.83\times 10^{40}`$ ergs s<sup>-1</sup>, respectively, in the 0.25–10 keV band. In the ROSAT band (0.1–2.4 keV) the luminosities were $`2.36\times 10^{40}`$ ergs s<sup>-1</sup> and $`1.04\times 10^{40}`$ ergs s<sup>-1</sup>, respectively.
## 4. ROSAT HRI Observation of NGC 4697
NGC 4697 was one of the few X-ray faint early-type galaxies for which a long ROSAT HRI observation (RH600825A01) exists in the HEASARC archive (78,744 seconds). The observation did not contain any time intervals with excessively high background so the entire observation was used. The inner $`8^{}\times 8^{}`$ of the HRI image is shown in Figure 1. A contour map of the HRI image is shown overlaying the the Digital Sky Survey optical image in Figure 2. Several X-ray point sources are evident from the contour plot. Background-subtracted count rates were calculated for each point source detected within the same $`4^{}`$ circle used to derive the PSPC and ASCA spectra. Background was selected from an annular ring with inner and outer radius of $`4^{}`$ and $`6^{}`$, and was corrected for vignetting before being subtracted from the source. In all, 12 point sources were detected at a significance of 2 $`\sigma `$ or higher, corresponding to a detection flux limit of $`1.0\times 10^{14}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> and a limiting luminosity of $`3.0\times 10^{38}`$ ergs s<sup>-1</sup>. From the source counts catalog of a deep HRI observation by Hasinger et al. (1998), we would expect there to be $``$1 serendipitous source in a $`4^{}`$ circle field at this flux level or higher. Thus, it is likely that all or nearly all of the X-ray sources in the field belong to NGC 4697.
The positions, number of counts, significance of detection, and luminosities of the 12 point sources are shown in Table 1. The count rates were converted to luminosities using the two-component spectral model found in § 3. Two of the sources (Sources 6 and 11) are coincident within the position errors with faint optical point sources with magnitudes of approximately 18, that do not correspond to any QSO found in the Veron-Cetty & Veron (1998) catalog. The optical source near Source 11 has been identified as a bright globular cluster of NGC 4697 (Hanes 1977). The central source in Figures 1 and 2 appears elongated in the north-south direction. When the center is viewed at higher resolution, the source is split into two sources (Sources 7 and 8 in Table 1), which are separated by $`7^{\prime \prime }`$, close to the resolution limit of the HRI. Source 7 is within $`2^{\prime \prime }`$ of the optical center of NGC 4697 (R.A. = 12:48:35.71 and Dec. = -5:48:02.9), as given by Wegner et al. (1996). Given the the HRI positional uncertainty of $``$5<sup>′′</sup>, the additional uncertainty due to the crowding of Sources 7 and 8 (which might contain additional components), and the optical position uncertainty of about $`1\stackrel{}{\mathrm{.}}25`$, it is possible that Source 7 or even 8 might be coincident with the optical center of the galaxy. Thus, one of these sources might be due to an active nucleus in the galaxy. On the other hand, there is no evidence for an AGN in NGC 4697; for example, the nucleus is not a radio source (e.g., the NVSS survey, Condon et al. 1998). The positive detection of only 12 sources does not allow us to determine if the sources follow the de Vaucouleurs stellar distribution, especially considering that one if not more of the sources is associated with a globular cluster of NGC 4697, and crowding in the center will underestimate the number of detected point sources in this region.
What is not evident from Figure 2 is the presence of very faint unresolved X-ray emission within $`4^{}`$ (the contours of Figure 2 were chosen to highlight the position of the X-ray point sources). Figure 4 shows the contours from a more heavily smoothed HRI image once again overlaying the optical image. The emission appears elongated in the direction of the optical major axis of the galaxy. The surface brightness profile of the X-ray emission follows the optical surface brightness profile derived by Jedrzejewski, Davies, & Illingworth (1987) out to 4, suggesting that the X-ray emission is distributed like the stars. However, it was found that a $`\beta `$-model profile with a shallow slope ($`\beta =0.40.45`$) also fit the data adequately. The faint unresolved emission is actually the dominant source of emission in the HRI image, comprising 79% of the total emission. The nature of this unresolved emission is discussed next.
## 5. Discussion
The 12 point sources detected within $`4^{}`$ of NGC 4697 constitute 21to whether the remaining unresolved emission is the summed emission from LMXBs below the detection threshold of the HRI or is a low temperature ISM with $`kT0.25`$ keV. The absorbed flux of the hard component found from the spectral fitting of the PSPC and ASCA data sets a lower limit on the contribution of LMXBs to the unresolved emission. In the ROSAT band, the hard component contributes 40% of the total (unabsorbed) flux. Thus, even if LMXBs possess only a hard X-ray component, another $`20\%`$ of the total emission must result from LMXBs below the detection threshold of the HRI.
As mentioned in the Introduction, LMXBs might also possess a soft X-ray component. This component would have gone unnoticed in most Galactic disk or bulge LMXBs, where the high Galactic column densities associated with the Galactic plane would completely absorb any very soft emission. Although there does not appear to be a soft component towards Galactic globular cluster LMXBs that lie in directions of low column densities, this might be the effect of the low metallicities in which the LMXBs formed (see Irwin & Bregman 1999b). The bulge of M31, however, provides the nearest opportunity to study a sample of LMXBs in an environment most similar to that found in early-type galaxies. Supper et al. (1997) analyzed the ROSAT PSPC image of M31 and found 22 point sources within the inner $`5^{}`$ of the bulge with luminosities in the range $`10^{36}10^{38}`$ erg s<sup>-1</sup>. Individually, these point sources had X-ray colors that indicated the presence of a significant soft component. Seven of the 22 sources had enough counts for Supper et al. (1997) to perform spectral fitting. The best-fit bremsstrahlung temperatures ranged from 0.45–1.5 keV, very similar to the result found for NGC 4697 for the same model. Although such a simplistic model is not physically plausible, it did indicate that the spectrum of the LMXBs were significantly softer than the canonical temperature of 5–10 keV previously assumed for LMXBs.
We can search for such a soft component in the point sources suspected to be LMXBs in NGC 4697. Of the 12 point sources detected by the HRI, sources 1, 6, 10, 11, and 12 were also detected and resolved in the PSPC (sources 10 and 12 were marginally resolved from each other in the PSPC). As mentioned above, Source 6 might be associated with background/foreground objects. The counts from the other four sources were summed in three energy bands, and two X-ray colors (C21 and C32) were defined from the ratio of the three bands:
$$\mathrm{C21}=\frac{\mathrm{counts}\mathrm{in}\mathrm{PI}\mathrm{bins}5290}{\mathrm{counts}\mathrm{in}\mathrm{PI}\mathrm{bins}1141},$$
(1)
and
$$\mathrm{C32}=\frac{\mathrm{counts}\mathrm{in}\mathrm{PI}\mathrm{bins}91202}{\mathrm{counts}\mathrm{in}\mathrm{PI}\mathrm{bins}5290}.$$
(2)
Note that these colors are not corrected for absorption.
Sources 1+10+11+12 combined yielded 359.5 background-subtracted counts, or 14% of the total X-ray emission within $`4^{}`$ in the PSPC image. The colors for the sources were (C21,C32 $`=0.57\pm 0.09,0.94\pm 0.16)`$. As a comparison, the colors for all emission within $`4^{}`$ were (C21,C32 $`=0.68\pm 0.05,0.77\pm 0.05)`$. Thus, the colors of the four resolved sources are consistent with the colors for the integrated emission from the galaxy. This argues that the four LMXBs must also possess a soft component similar to that found in the global spectrum. Otherwise, the integrated emission would have had significantly different colors than the point sources. Conversely, the colors predicted from a 5.2 keV bremsstrahlung model with an absorbing column density of $`2.12\times 10^{20}`$ cm<sup>-2</sup> were (C21,C32 $`=1.204,1.893)`$. This differs from the colors of the four point sources by 7.0$`\sigma `$ and 6.0$`\sigma `$ for C21 and C32, respectively. Clearly, if LMXBs were described by only a hard component, the colors of the four point sources would be significantly harder than what they are.
We have calculated the colors of NGC 4697 in elliptical annuli (with ellipticity of 0.42 and a position angle of 67 to match the optical profile; Figure 3). C32 peaks inside of $`1\stackrel{}{\mathrm{.}}5`$ but flattens at larger radii, while C21 is constant throughout the galaxy. The peak in C32 may result from the presence of an AGN in NGC 4697. The constancy of C21 and C32 outside of $`1\stackrel{}{\mathrm{.}}5`$ implies a single emission mechanism (or two emission mechanisms with the same spatial distribution), with colors around 0.6 for both colors, in rough agreement with the colors of the four point sources. The C32 color of the diffuse emission is somewhat lower than the C32 value of the four point sources, possibly indicating the presence of an ISM component, although the discrepancy between the two colors is only at the 2$`\sigma `$ level. The colors predicted from a 0.2 keV, 20% metallicity ISM are (C21,C32 $`=0.54,0.22)`$, whereas the colors for a 0.3 keV, 20% metallicity ISM are (C21,C32 $`=1.17,0.46)`$. Thus, if any ISM is present in NGC 4697, its temperature must be below 0.3 keV and at a low metallicity, or else the C21 color of the LMXBs+ISM would be higher than the LMXBs alone, which is not observed. In conclusion, we cannot rule out the presence of some low temperature ISM in NGC 4697, although it is certain that an ISM cannot constitute a majority of the emission.
Irwin & Bregman (1999b) found that Galactic and M31 globular cluster LMXB X-ray colors were correlated with the metallicity of the globular cluster, in the sense that higher metallicity globular clusters had LMXBs with softer X-ray colors. If this correlation extended to all LMXBs, it would predict that the X-ray colors of NGC 4697 should harden with increasing radius, since metallicity decreases with radius in elliptical galaxies. Furthermore, the metallicity of NGC 4697 is rather low. Within half an effective radius, the average metallicity is only 50% solar (Trager et al. 2000). The metallicity-color relation of Irwin & Bregman (1999b) would predict a C32 color of about 1.5 within half an effective radius, and an increase with increasing radius. The colors of the four resolved sources as well as the unresolved emission are at odds with this metallicity-color relation. Apparently, if such a metallicity-color relation truly exists for LMXBs, it only applies to LMXBs that reside in globular clusters. The X-ray source in NGC 4697 associated with a globular cluster (Source 11) has a C32 color of $`1.14\pm 0.30`$, which would be consistent with the metallicity-color relation if the globular cluster has a high metallicity.
We also investigated the possibility that LMXBs below the detection threshold of the HRI could account for the unresolved emission given a reasonable LMXB luminosity distribution function for NGC 4697. We assumed a luminosity distribution function $`N(>L_X)L_X^{1.3}`$, which is consistent within the errors with the luminosity distribution function of point sources in M31 with luminosities greater than $`2\times 10^{37}`$ ergs s<sup>-1</sup> (Primini et al. 1993). The function was normalized to yield the observed X-ray luminosity of NGC 4697 when integrated over all LMXB luminosities. This model predicted 11 sources with luminosities over $`3\times 10^{38}`$ ergs s<sup>-1</sup>, which contributed 18% to the total X-ray emission from LMXBs. This agrees well with what was observed with the HRI; neglecting the point source associated with an unidentified optical counterpart, the remaining 11 detected sources comprised 19% of the total emission. Thus, if the luminosity distribution function of LMXBs of NGC 4697 is similar to that of the brighter LMXBs in M31, the integrated emission from LMXBs below the detection threshold of the HRI can account for most of the unresolved emission.
Interestingly, the detection limit of the HRI observation of $`3\times 10^{38}`$ ergs s<sup>-1</sup> lies above the Eddington luminosity limit for a $`1.4M_{}`$ neutron star. LMXBs of similar luminosities as the ones found here exist in our own Galaxy. A compilation by Christian & Swank (1997) found eight galactic LMXBs with luminosities greater than $`3\times 10^{38}`$ ergs s<sup>-1</sup> (we have converted their 0.7–4.5 keV luminosities to 0.25-10 keV luminosities using the spectral model of § 3). These high luminosities imply either that the compact object within the binary is a black hole (with $`M_{BH}6M_{}`$ for the most luminous binaries) or that the luminosities truly exceed the Eddington limit for a neutron star. The former would imply that active binaries with massive black holes are fairly common in galaxies. It should be noted that the bulge of M31 lacks the very high luminosity LMXBs that NGC 4697 has; the brightest LMXB in M31 was only $`1.8\times 10^{38}`$ ergs s<sup>-1</sup> (Supper et al. 1997). However, this is likely the result of small number statistics. Simulations of the bulge of M31 using a luminosity distribution function of the form $`N(>L_X)L_X^{1.3}`$ indicated that only 1–3 sources with luminosities exceeding $`10^{38}`$ ergs s<sup>-1</sup> should be found in M31. In five separate simulations of M31, the peak luminosity for an LMXB did not exceed $`3\times 10^{38}`$ ergs s<sup>-1</sup>, in agreement with observation.
We cannot rule out the presence of at least some interstellar medium in NGC 4697 and X-ray faint early-type galaxies in general. Using the X-ray temperature–optical velocity dispersion relation of Davis & White (1996), any ISM present in NGC 4697 would be expected to have a temperature around 0.3 keV. This would be very difficult to distinguish from the soft component from LMXBs on a spectroscopic basis alone. What is needed to separate the ISM component from the LMXB emission is the high spatial resolution that can be afforded by Chandra. Below, we present a simulation of what we expect the emission from NGC 4697 to look like in the event that the emission is composed solely of LMXBs.
## 6. Simulation of Chandra Observation of NGC 4697
We have an approved Cycle 1 40,000 s Chandra observation of NGC 4697; here we show that this observation should resolve the hard and soft X-ray emission into individual sources, assuming that the emission is from LMXBs. Conversely, the observation should cleanly separate a truly diffuse emission from that of LMXBs. Since the main goal is to resolve the issue of the very soft component, the soft X-ray sensitive backside-illuminated (BI) S3 chip of the ACIS-S array will be used for the observation. We have used the MARX (Model of AXAF Response to X-rays; Wise, Huenemoerder, & Davis 1997) Simulator to generate a synthetic image of NGC 4697. The MARX Simulator takes as input the desired spectral model and spatial distribution model of an X-ray source and creates an image of the source as it would appear once having passed through the optics of Chandra. The spectral and spatial distribution models described below were fed into MARX using the ACIS-S BI response to produce an image of NGC 4697 for a 40,000 s observation.
For the spectra of the LMXBs, we assume a model that best fit the joint ROSAT PSPC + ASCA spectrum of NGC 4697 discussed in § 3. For the spatial distribution of the X-ray emission, we assume that the LMXBs have the same spatial distribution as the stellar light. We take the optical distribution to be a de Vaucouleurs profile (de Vaucouleurs et al. 1991) with a mean half-light radius of $`72^{\prime \prime }`$, an effective semimajor axis of $`95^{\prime \prime }`$, an effective semiminor axis of $`55^{\prime \prime }`$, resulting in an ellipticity of 0.42, and elongated at a position angle of 67. (Jedrzejewski et al. 1987; Faber et al. 1989; Peletier et al. 1990). As before, we assume a luminosity distribution function for the LMXBs of the form $`N(>L_X)L_X^{1.3}`$. The luminosity distribution function contained point sources with 0.25–10.0 keV luminosities between $`10^{37}`$ and $`10^{39}`$ ergs s<sup>-1</sup>. A model galaxy was created by drawing LMXBs at random from the luminosity distribution function and the spatial distribution until the sum of the luminosities of the LMXBs totaled the X-ray luminosity of the galaxy. We did not use the observed positions and fluxes of the sources in Table 1 in the model; the sources are randomly drawn from the optical surface brightness distribution and X-ray luminosity distribution. We did not include background in this simulation. The background in the Chandra ACIS S-3 chip is highly variable (Markevitch 1999), and its level during the observation will affect the ability to detect the weakest sources.
Figure 5 shows the simulated image of NGC 4697. The image shows the inner $`4^{}`$ by $`4^{}`$ of the galaxy. The excellent spatial resolution of Chandra is evident. The dynamic range of the source brightness in this image is not clear from this greyscale representation, since all of the strong sources are nearly the same size (set by the resolution) and black. However, the sources cover a range of $``$25 in flux. For our simulation, we assumed that 16 source counts would be needed to give a $`3\sigma `$ detection considering the variable background of the ACIS S-3 chip. This assumption predicts that $``$100 sources would be detected at $`3\sigma `$ with a detection threshold of $`6\times 10^{37}`$ ergs s<sup>-1</sup>. These 100 sources comprise 50% of the total luminosity of the galaxy. In the 0.25–0.8 keV band, $``$50 sources were detected at $`3\sigma `$.
We have also simulated a Chandra image where the LMXBs have only a hard (5.2 keV) bremsstrahlung component that comprises 59% of the 0.25–10 keV emission (this was the relative contribution of the hard component to the total emission found in § 3). In this case, only about 12 LMXBs were detected in the 0.25–0.8 keV band. If LMXBs do possess a soft component, it will be immediately obvious from the 0.25–0.8 keV Chandra image of NGC 4697.
X-ray spectra of some of the brighter sources, X-ray colors of fainter sources, the cumulative X-ray spectra of individual sources and of any remaining unresolved diffuse emission can be determined, and should resolve the question of the origin of the hard and soft X-ray components in X-ray faint ellipticals.
A number of fundamental questions regarding the X-ray emission from early-type and Sa galaxies will be answered with Chandra. With the expected number of detectable LMXBs, accurate luminosity distribution functions can be determined. This is something that has been accomplished for only a handful of nearby galaxies, none of which are normal early-type galaxies. This will provide important clues to the stellar evolution of binary stars in galaxies. There is evidence that there may not be a universal stellar X-ray–to–optical luminosity ratio. A range of $`L_{X,stellar}/L_B`$ values has been suggested by comparison of the X-ray emission of the bulge of M31 to that of Cen A (Turner et al. 1997), and also to several very X-ray faint early-type galaxies in the Irwin & Sarazin (1998b) survey. If nearly all the X-ray emission in NGC 4697 turns out to be stellar in nature, it would be difficult to explain why the early-type galaxy NGC 5102, for example, has a 0.5-2.0 keV $`L_X/L_B`$ value that is 19 times lower than that of NGC 4697 (Irwin & Sarazin 1998b), if $`L_{X,stellar}/L_B`$ is constant from galaxy to galaxy. Without knowledge of the luminosity distribution function, it cannot be determined if this difference is a result of a different slope in the distribution function among galaxies, or if the slopes are the same but the normalizations of the function (relative to the optical luminosity of the galaxy) are different. The separation of the stellar and gaseous components is necessary to resolve this issue.
Since LMXB emission is expected to be present in all early-type galaxies, the magnitude of this component needs to be known accurately in order to subtract the LMXB contribution from the total X-ray emission in gas-rich early-type galaxies. This will make estimates of the mass of elliptical galaxies based on the assumption that the gas is in hydrostatic equilibrium more accurate. This is particularly important in the event that LMXBs possess a strong soft component. Although, the contribution of LMXBs to the total X-ray emission from very X-ray bright Virgo elliptical galaxies such as NGC 4472 and NGC 4636 will be small (less than 10% in the 0.1–2.4 keV band), LMXBs might contribute a significant percentage of the X-ray emission in galaxies of intermediate X-ray brightness, and will have to be dealt with accordingly if the true amount of X-ray gas is to be determined.
## 7. Conclusions
We have analyzed deep ASCA, ROSAT PSPC, and ROSAT HRI images of the X-ray faint early-type galaxy NGC 4697. Much like other X-ray faint early-type systems, the spectrum of NGC 4697 is characterized by hard (5 keV) plus very soft (0.3 keV) emission. Whereas the nature of the hard emission is generally regarded as the integrated emission from LMXBs, we have provided additional evidence that much of the soft emission also emanates from LMXBs. Four of the 12 HRI point sources were resolved by the PSPC and were found to have soft X-ray colors that were very similar to those of the galaxy as a whole. These colors were significantly softer than the colors predicted if only the hard component was attributed to LMXBs.
The 12 point sources detected by the HRI comprised 21% of the total X-ray emission within $`4^{}`$ of NGC 4697. Given a luminosity distribution function consistent with that for the brighter point sources in M31, the remaining unresolved emission could emanate solely from LMXBs below the detection threshold of the observation. However, the presence of a low temperature ISM could not be completely ruled out. Higher spatial resolution data afforded by Chandra should successfully resolve much of the LMXB emission, as we have shown in our simulations. The determination of the origin of the soft component will have important implications concerning our understanding of the fate of gas lost from stars in galaxies as well as the X-ray emission mechanism of LMXBs.
We thanks the referee, Fabrizio Brighenti, for many useful comments and suggestions concerning the manuscript. This research has made use of data obtained through the High Energy Astrophysics Science Archive Research Center Online Service, provided by the NASA/Goddard Space Flight Center. The optical image of NGC 4697 is from the Digital Sky Survey, which were produced at the Space Telescope Science Institute. The images of these surveys are based on photographic data obtained using the Oschin Schmidt Telescope on Palomar Mountain and the UK Schmidt Telescope. J. A. I. was supported by Chandra Fellowship grant PF9-10009, awarded through the Chandra Science Center. The Chandra Science Center is operated by the Smithsonian Astrophysical Observatory for NASA under contract NAS8-39073. C. L. S. was supported in part by NASA Chandra grant GO0-1019X. J. N. B. was supported by NASA grant NAG5-3247.
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# Wetting of a symmetrical binary fluid mixture on a wall
## I Introduction
Phase transitions are well known to be influenced by geometrical confinement . In practice, confinement is often imposed by rigid external constraints, for example the surfaces of porous or artificially nanostructured media. However, it can also be an inherent feature of a system, as occurs for a liquid wetting film bound to a solid substrate and in equilibrium with its vapour . In such a situation the liquid is confined between the rigid substrate and the flexible liquid-vapour interface.
The effects of confinement are particularly pronounced in the region of critical points. Under such conditions the system exhibits strong order parameter fluctuations, the correlation length of which may become comparable with the linear dimension of the confined system. When this occurs the effects of confinement are felt not just near the confining surfaces, but propagate throughout the system .
Critical fluctuation are relevant to the properties of liquid wetting films if the liquid in question possesses an additional internal degree of freedom. Then the state of the liquid is described not only by its number density, but by an additional parameter measuring the degree of internal order. Examples are binary liquids, where the additional order parameter is the relative concentration of species, and ferrofluids where it is the magnetisation. In such systems the geometrical constraint (i.e. the film thickness) can itself depend on the state of order in the liquid film. For example, simulation and experiment have recently shown that critical concentration fluctuations can change the equilibrium thickness of a wetting layer of a binary liquid—the so-called critical Casimir effect .
It transpires, however, that the interesting consequences of interplay between the degree of order in the wetting layer and film thickness are not limited to the critical point itself. To illustrate this, it is instructive to consider the following Gedanken experiment. Let us take a symmetrical binary fluid, i.e. a fluid in which particles of the same species have one strength of interaction, while interactions between dissimilar species have another strength. As elucidated in ref. , it is possible to arrange for such a system to exhibit a line of continuous demixing transitions, terminating in a critical end point on the liquid side of liquid-vapour coexistence. Suppose now that the fluid is placed in contact with a non-selective attractive substrate (wall) acting equally on both species. If the wall is sufficiently attractive, complete wetting occurs at and above the critical end point temperature $`T_{cep}`$ as liquid-vapour coexistence is approached from the vapour side. But what happens for $`T<T_{cep}`$? Far from coexistence, the wetting film is sufficiently thin that demixing will certainly be suppressed. On approaching the coexistence curve, however, the film
thickness grows and it is tempting to argue that it eventually exceeds the correlation length of composition fluctuations, whereupon the film spontaneously demixes.
Notwithstanding the appealing simplicity of this argument, it turns out to contain two flaws which render the actual situation rather more complex. First, the thickness of the mixed wetting film will not increase beyond all limits below the critical end point $`T_{cep}`$. This is because a hypothetical mixed bulk liquid would not coexist with the vapour phase at the same chemical potential $`\mu _0`$ as a demixed liquid, but rather at a chemical potential which is shifted by
$$\mu _{}\mu _0(T_{cep}T)^{2\alpha }$$
(1)
towards the liquid side in the $`\mu T`$ plane. Since the thickness of the mixed film grows as $`\mathrm{ln}(\mu _{}\mu )`$, it is bounded from above by
$$l_{}\mathrm{ln}(\mu _{}\mu _0)\mathrm{ln}(T_{cep}T).$$
(2)
The maximum thickness $`l_{}`$ of a hypothetical mixed film thus diverges logarithmically on approaching $`T_{cep}`$. In contrast, the correlation length $`\xi `$ diverges much faster, like $`\xi (T_{cep}T)^\nu `$ and will thus always exceed $`l_{}`$ sufficiently close to $`T_{cep}`$.
The second flaw in our argument is its implicit assumption that the composition or order parameter profile is confined in an effectively steplike density profile, i.e. that the interfacial width between the liquid and the vapour is much smaller than the correlation length of composition fluctuations. Although this is true in the region of the critical end point, for temperature sufficiently below $`T_{cep}`$ the correlation lengths of density and composition fluctuations can be comparable and the interplay between the two subtle.
In the present study, we deploy mean field calculation and Monte Carlo simulation to elucidate the range of possible wetting behaviour of a symmetrical binary fluid mixture at a non-selective attractive wall for temperatures below $`T_{cep}`$. For the sake of simplicity, we have chosen to ignore long range dispersion forces in the analytical calculations, instead taking the interactions to be short ranged. This allows us to base our study on a generic Ginzburg-Landau model which we solve numerically and analytically within a square gradient approximation. The latter leads to the construction of a film free energy (effective interface potential) highlighting the role of the different length scales involved in the problem. We show that the competition of length scales results in wetting phase behavior considerably more complex than has hitherto been appreciated. The analytical results are compared with (and supported by) detailed Monte Carlo simulations of a binary Lennard-Jones fluid in a semi-infinite geometry, interacting with a non-selective attractive substrate via dispersion forces.
With regard to previous related work, the sole discussion of wetting of symmetrical binary fluids at a non-selective wall (of which we are aware) is that of Dietrich and Schick who considered them in a sharp kink approximation treatment of binary fluids having long-ranged interactions. Most other work on the wetting properties of binary fluids has focused on the case of a selective substrate (favouring one component) . Although such models correspond more closely than ours to experimental conditions , they lack the aspect of simultaneous demixing/ordering and wetting which is of interest to us here. It should be stressed, however, that realisations of fluids having symmetrical internal degrees of freedom do in fact exist, notably in the form of ferrofluids , so our model is of more than purely theoretical interest.
More general studies of wetting in systems with more than one order parameter and associated length scales have been discussed by Hauge , who pointed out that wetting exponents may become nonuniversal even on the mean field level due to the competition of length scales. Later studies have often focused on this nonuniversality, e.g. in the context of wetting phenomena in superconductors , alloys and related systems .
The present paper is organised as follows. In section II we introduce our Ginzburg-Landau free energy functional and obtain its wetting behaviour in the limits of infinite and vanishing order parameter stiffness. At intermediate values of the stiffness parameter the wetting behaviour is found firstly via an analytical minimisation of the functional within a piecewise parabolic potential approximation (sec. II B), and then (in sec. II C) via a numerical minimisation of the free energy functional to obtain the density/order parameter profiles. In section III A we report the results of grand canonical Monte Carlo studies of a symmetrical binary Lennard-Jones fluid at an attractive structureless wall. The density and order parameter profiles with respect to the wall are obtained along a sub-critical isotherm for a number of different wall-fluid potential strengths. Finally we compare and discuss the mean-field and simulation results in section IV.
## II Ginzburg-Landau theory
Our theoretical studies are based on a generic Ginzburg-Landau functional for a system with two order parameters $`\varphi (\stackrel{}{r},z)`$ and $`m(\stackrel{}{r},z)`$:
$``$ $`=`$ $`{\displaystyle 𝑑\stackrel{}{r}_0^{\mathrm{}}𝑑z\left\{\frac{g}{2}(\varphi )^2+\frac{\gamma }{2}(m)^2+f(m,\varphi )\right\}}`$ (4)
$`+{\displaystyle 𝑑\stackrel{}{r}f_s(m,\varphi )}|_{z=0}`$
with the bulk free energy density
$`f(m,\varphi )`$ $`=`$ $`{\displaystyle \frac{a_\varphi }{2}}\varphi ^2+{\displaystyle \frac{b_\varphi }{4}}\varphi ^4{\displaystyle \frac{a_m}{2}}m^2+{\displaystyle \frac{b_m}{4}}m^4`$ (6)
$`+\mu \varphi \kappa m^2\varphi `$
and the bare surface free energy at the wall
$$f_s(m,\varphi )=\frac{C_\varphi }{2}\varphi ^2+H_\varphi \varphi +\frac{C_m}{2}m^2.$$
(7)
The $`z`$-axis is taken to be perpendicular to the wall and $`𝑑\stackrel{}{r}`$ integrates over the remaining spatial dimensions. In our case, the quantity $`m`$ is related to the difference between the partial densities of the two components, $`m(\rho _A\rho _B)`$, and $`\varphi `$ to the total density, $`\varphi (\rho \rho _0)`$, where the reference density $`\rho _0`$ is chosen in the liquid/vapour coexistence region such the cubic term proportional to $`(\rho \rho _0)^3`$ in (6) vanishes. Below the liquid/vapour critical point, it is convenient to set the units of $`\varphi `$, $`m`$, $`F`$ and of the length such that $`b_m=b_\varphi =a_\varphi =g=1`$, and to define $`\theta =a_m1`$. The bulk free energy density then takes the form
$`f(m,\varphi )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\varphi ^2+{\displaystyle \frac{1}{4}}\varphi ^4{\displaystyle \frac{\theta }{2}}m^2+{\displaystyle \frac{1}{4}}m^4`$ (9)
$`\mu \varphi +\kappa (1\varphi )m^2.`$
The bulk properties of this model have been discussed earlier. A $`\lambda `$-line $`\theta _\lambda (\mu )`$ of continuous transitions separates the mixed fluid from the demixed fluid at large negative $`\mu `$, corresponding to large densities $`\varphi `$. As long as $`\kappa <1`$, it is terminated by the onset of liquid/vapour coexistence at a critical end point ($`\theta _{cep}=0,\mu _{cep}=0`$). The parameter $`\mu `$ is field like and $`\theta `$ is temperature like, $`\theta (TT_{cep})`$, where $`T_{cep}`$ is the critical end point temperature. Above $`\theta _{cep}`$, liquid/vapour coexistence is encountered at $`\mu =0`$, and below $`\theta _{cep}`$, at
$$\mu _c=\frac{\theta ^2}{8(1\kappa ^2)}.$$
(10)
The coexisting liquid and gas phases are characterized by the order parameters (to linear order in $`\mu `$)
$$m_{}^{}=0\varphi _{}^{}=1\mu /2$$
(11)
in the gas phase, and
$$m_+^{}=\frac{\theta \kappa \mu }{1\kappa ^2}\varphi _+^{}=1\frac{\mu }{2}+\frac{\kappa }{2}m_+^{}^2$$
(12)
in the liquid phase. These expressions are also valid in the regime where the liquid or gas phase are metastable.
Minimizing the functional (4) yields the Euler Lagrange equations
$$g\frac{d^2\varphi }{(dz)^2}=\frac{f}{\varphi }\gamma \frac{d^2m}{(dz)^2}=\frac{f}{m}$$
(13)
with the boundary conditions
$$g\frac{d\varphi }{dz}|_{z=0}=\frac{f_s}{\varphi }\gamma \frac{dm}{d}|_{z=0}=\frac{f_s}{m}.$$
(14)
We wish to study a situation where the mixed liquid ($`m0`$) wets the wall at $`\mu =0`$ (coexistence between vapour and mixed liquid). To ensure this under all circumstances, we choose $`H_\varphi =\varphi _0C_\varphi `$ with $`\varphi _0>\varphi _+^{}`$ and take the limit $`C_\varphi \mathrm{}`$, which is equivalent to constraining the surface density at the fixed value $`\varphi (0)=\varphi _0`$. The surface coupling $`C_m`$ is taken to be positive. It accounts for weakening of the demixing tendency at the surface due to the reduced number of interacting neighbors.
One possible solution of the Euler Lagrange equations describes a mixed film at a wall. In this case, $`m(z)=0`$ everywhere and one is left with one order parameter $`\varphi `$ only. The bulk value of $`\varphi `$ in the metastable mixed phase is $`\varphi _+^{(0)}=1\mu /2`$. The standard way of solving the problem shall be sketched briefly for future reference. One begins by identifying the integration constant
$$\frac{1}{2}(\frac{d\varphi }{dz})^2f(\varphi )+f(\varphi _{}^{(0)})0.,$$
(15)
which gives an expression for $`d\varphi /dz`$ as a function of $`\varphi `$. The surface free energy can then be expressed as an integral over $`\varphi `$
$$F_{exc}^{(0)}=_{\varphi _{}^{(0)}}^{\varphi _0}𝑑\varphi \sqrt{2(f(\varphi )f(\varphi _{}^{(0)}))}$$
(16)
and the excess density $`\varphi _{exc}^{(0)}=\frac{1}{2}_0^{\mathrm{}}𝑑z\left[\varphi (z)\varphi _{}^{(0)}\right]`$ at the surface can be calculated via
$$\varphi _{exc}^{(0)}=\frac{1}{2}_{\varphi _{}^{(0)}}^{\varphi _0}\frac{d\varphi (\varphi \varphi _{}^{(0)})}{\sqrt{2(f(\varphi )f(\varphi _{}^{(0)}))}}.$$
(17)
As long as $`|f(\varphi _+^{(0)}f(\varphi _{})|(\varphi _+^{(0)}\varphi _{})^2f^{\prime \prime }(\varphi _+^{(0)}`$, which is true for $`\mu 2`$, the main contribution to this integral stems from $`\varphi `$ values around $`\varphi _+^{(0)}`$. The numerator in the integrand can then be expanded around $`\varphi _+^{(0)}`$. Carrying this to second order and assuming $`\mu (\varphi _0\varphi _+^{(0)})`$, one obtains
$`\varphi _{exc}^{(0)}`$ $``$ $`\sqrt{{\displaystyle \frac{1}{2}}}\mathrm{ln}({\displaystyle \frac{2}{\varphi _+^{(0)}\varphi _0}})(\varphi _0<\varphi _+^{(0)})`$ (18)
$`\varphi _{exc}^{(0)}`$ $``$ $`\sqrt{{\displaystyle \frac{1}{2}}}\mathrm{ln}({\displaystyle \frac{4(\varphi _0\varphi _+^{(0)})}{\mu }})(\varphi _0>\varphi _+^{(0)}).`$ (19)
Above the bulk demixing transition, the mixed film thus wets the wall at coexistence ($`\mu 0`$) for $`\varphi _0>\varphi _+^{(0)}`$, and maintains a finite thickness for $`\varphi _0<\varphi _+^{(0)}`$. We will choose $`\varphi _0>\varphi _+^{(0)}`$ hereafter. From eqn. 16, one calculates the surface free energy to leading order in $`\mu `$ and $`(\varphi _0\varphi _+^{(0)})`$.
$$F_{exc}^{(0)}=\frac{2\sqrt{2}}{3}+\frac{1}{\sqrt{2}}(\varphi _0\varphi _+^{(0)})^2$$
(20)
Below the bulk demixing transition, $`\mu =\mu _c>0`$ at coexistence and the thickness of the mixed film remains finite under all circumstances.
In the following, we shall first analyse the wetting behavior for the limiting cases where the order parameter varies varies on very short length scales ($`\gamma /m_{+}^{}{}_{}{}^{2}0`$) and on long length scales ($`\gamma /m_{+}^{}{}_{}{}^{2}\mathrm{}`$) compared to the density. Then we will discuss the general case of intermediate $`\gamma `$, first analytically in an approximation where the potential (9) is replaced by a piecewise quadratic potential, and then numerically with the full potential (9).
### A Limiting cases
We consider first the wetting behavior at ($`\gamma /m_{+}^{}{}_{}{}^{2}0`$). In this case, $`m`$ adapts locally to $`\varphi `$, and the order parameter profile $`m(z)`$ can be written as $`m(\varphi (z))`$ with $`m(\varphi )=\theta +2\kappa (\varphi 1)`$ for $`\varphi <1\theta /2\kappa `$ and $`m(\varphi )=0`$ otherwise. Hence we are left with the effective one order parameter problem of calculating the density profile $`\varphi (z)`$ in the slightly altered potential $`\widehat{f}(\varphi )=f(m(\varphi ),\varphi )`$. Since $`\widehat{f}(\varphi )`$ is a smooth function with two minima, one can proceed as sketched above for the mixed film, with the analogous result: The demixed film wets the wall at $`\varphi >\varphi _+^{}`$.
The analysis of the opposite case, ($`\gamma /m_{+}^{}{}_{}{}^{2}\mathrm{}`$), is somewhat more involved. Here $`\varphi `$ adapts locally to $`m`$; however, the bulk equation $`f/\varphi =\varphi ^3\varphi +\mu \kappa m^2=0`$ has two solutions $`\varphi _\pm (m)`$. One conveniently separates the profiles into four parts (I) – (IV) as indicated in Figure 1. The regions (I) and (III) are narrow slabs where $`\varphi (z)`$ varies rapidly and $`m`$ can be approximated by a constant, $`m=m_0`$ at the surface (I) and $`m_2`$ at the interface (III). In (I), $`\varphi `$ drops from it’s surface value $`\varphi _0`$ to the local equilibrium value $`\varphi _+(m_0)`$, and in (III), it switches from $`\varphi _+(m_2)`$ to $`\varphi _{}(m_2)`$. The other two regions, (II) and (IV), are much wider; The order parameter $`m(z)`$ varies slowly and $`\varphi (z)`$ adjusts locally to $`m(z)`$, such that $`\varphi =\varphi _+(m)`$ in (II) and $`\varphi =\varphi _{}(m)`$ in (IV).
To make the argument more quantitative, we specify the actual subdivision of the excess free energy of (4),
$$_{exc}=f_s(m_0,\varphi _0)+_I+_{II}+_{III}+_{IV}$$
(21)
with
$`_I`$ $`=`$ $`{\displaystyle _0^\delta }𝑑z\left[{\displaystyle \frac{1}{2}}({\displaystyle \frac{d\varphi }{dz}})^2+f(m_0,\varphi )f(m_0,\varphi _+(m_0))\right]`$
$`_{II}`$ $`=`$ $`{\displaystyle _0^l}𝑑z\left[{\displaystyle \frac{1}{2}}\gamma ({\displaystyle \frac{dm}{dz}})^2+f(m,\varphi _+(m))\right]`$
$`_{III}`$ $`=`$ $`{\displaystyle _{l\delta }^l}𝑑z\left[{\displaystyle \frac{1}{2}}({\displaystyle \frac{d\varphi }{dz}})^2+f(m_2,\varphi )f(m_2,\varphi _+(m_2))\right]`$
$`+`$ $`{\displaystyle _l^{l+\delta }}𝑑z\left[{\displaystyle \frac{1}{2}}({\displaystyle \frac{d\varphi }{dz}})^2+f(m_2,\varphi )f(m_2,\varphi _{}(m_2))\right]`$
$`_{IV}`$ $`=`$ $`{\displaystyle _0^l}𝑑z\left[{\displaystyle \frac{1}{2}}\gamma ({\displaystyle \frac{dm}{dz}})^2+f(m,\varphi _{}(m))\right].`$
The calculation for the regions (I), (III), and (IV) can proceed in an analogous way as sketched earlier for the mixed film: The profiles of $`\varphi (z)`$ in (I), (III), and of $`m(z)`$ in (IV) are monotonic and the integration constant (cf. 15) is known (zero). One obtains to leading order in $`\mu `$ and $`(\varphi _0\varphi _+(m_0))`$
$`F_I`$ $`=`$ $`\sqrt{{\displaystyle \frac{1}{2}}}(\varphi _0\varphi _+(m_0))^2+\mathrm{}`$ (22)
$`F_{III}`$ $`=`$ $`{\displaystyle \frac{2\sqrt{2}}{3}}+𝒪((\mu \kappa m_2{}_{}{}^{2})^2)`$ (23)
$`F_{IV}`$ $`=`$ $`{\displaystyle \frac{m_2^2}{2}}\sqrt{\gamma (4\kappa +\kappa \mu \theta )}`$ (24)
In the region (II), the integration constant is unknown,
$$\frac{1}{2}\gamma (\frac{dm}{dz})^2f(m,\varphi _+(m))=p,$$
(25)
with $`p>0`$ if the profile of $`m(z)`$ is monotonic, and $`p<0`$ if $`m(z)`$ is nonmonotonic, like in Fig. 1. A connection between $`p`$ and the width $`l`$ of the film can be established using $`l=_{m_2}^{m_0}𝑑m/|dm/dz|`$ in the first case, and
$$l=_{m_0}^{m_{max}}\frac{dm}{|dm/dz|}+_{m_2}^{m_{max}}\frac{dm}{|dm/dz|}$$
in the second case, where $`m_{max}`$ solves $`p=f(m_{max},\varphi _+(m_{max}))`$. Next we expand the function $`f(m,\varphi _+(m))`$ about it’s minimum $`m_+^{}`$, leading to
$$f(m,\varphi _+(m))(1\kappa ^2)m_+{}_{}{}^{4}((\frac{m}{m_+^{}}1)^2\frac{1}{4})$$
(26)
One deduces the characteristic length scale,
$$\lambda =\sqrt{\frac{\gamma }{2(1\kappa ^2)}}\frac{1}{m_+^{}},$$
(27)
which grows very large in the limit $`\gamma /m_{+}^{}{}_{}{}^{2}\mathrm{}`$. The result for $`_{II}`$ can therefore be expanded in powers of $`e^{l/\lambda }`$. After adding up all contributions (I)–(IV) and minimizing with respect to $`m_2`$, the total excess free energy of the demixed film takes the form $`F_{exc}=F_{surf}(m_0,\varphi _0)+F_{int}+V(l)`$ with the surface contribution
$`F_{surf}`$ $`=`$ $`f_s(m_0,\varphi _0)+\sqrt{{\displaystyle \frac{1}{2}}}(\varphi _0\varphi _+(m_0))^2`$ (29)
$`+8\lambda \mu _c(1m_0/m_+^{})^2,`$
the interface contribution
$$F_{int}=\frac{2\sqrt{2}}{3}+8\lambda \mu _c,$$
(30)
and a surface/interface interaction term
$`V(l)`$ $`=`$ $`2(\mu \mu _c)l32\lambda \mu _c(1m_0/m_+^{})e^{l/\lambda }`$ (32)
$`+16\lambda \mu _c(A(1m_0/m_+^{})^2B)e^{2l/\lambda },`$
where $`B=[\lambda \sqrt{\gamma (\kappa m_+^{}{}_{}{}^{2}(1\kappa ^2)/4}]^1`$, and $`A=3`$ or $`A=1`$, depending on whether or not the profile $`m(z)`$ is monotonic.
The result can now be discussed. At $`m_0<m_+^{}`$, the leading term $`e^{l/\lambda }`$ of the potential $`V(l)`$ is attractive, and wetting is not possible. At $`m_0>m_+^{}`$, an infinitely thick demixed film is metastable at coexistence. It’s free energy difference to the mixed film $`\mathrm{\Delta }F=F_{exc}F_{exc}^{(0)}`$ is up to third order in $`m_+^{}`$
$`\mathrm{\Delta }F`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\sqrt{\gamma (1\kappa ^2)}(\mathrm{\hspace{0.25em}1}+(1m_0/m_+^{})^2)m_{+}^{}{}_{}{}^{3}`$ (34)
$`+(C_m/\sqrt{2}\kappa (\varphi _0\varphi _+(m_0)))m_0^2)`$
The limit $`\gamma /m_{+}^{}{}_{}{}^{2}\mathrm{}`$ can be taken in two ways: either $`\gamma \mathrm{}`$ at fixed $`m_+^{}`$, or $`m_+^{}0`$ at fixed $`\gamma `$. In the first case the first term in eqn. (34) dominates and the free energy of the demixed film exceeds that of the mixed film: The film remains mixed and dewets accordingly.
The second case is more subtle. Here, the second term dominates, and the free energy of the mixed film may be less favorable, depending on the ratio of $`C_m`$ and $`(\varphi _0\varphi _+^{})`$. Note that the density enhancement at the surface, $`\sqrt{2}\kappa (\varphi _0\varphi _+^{})`$, acts as an additional surface coupling, which opposes the effect of $`C_m`$. The parameter $`C_m`$ accounts for the direct reduction of interacting neighbours right at the surface. It is counterbalanced by the fact that the density $`\varphi _0`$ close to the surface is higher than in the bulk. If the latter effect dominates, the film demixes at the surface even for $`m_+^{}0`$ or $`TT_{cep}`$.
### B Analytical results in a piecewise parabolic potential
At fixed $`m_+^{}`$, we have seen that the demixed film wets the substrate in the limit $`\gamma 0`$, where the order parameter $`m`$ varies much faster than the density $`\varphi `$, and dewets at $`\gamma \mathrm{}`$, where the density varies much faster than the order parameter. Now we consider intermediate values of $`\gamma `$, where the two characteristic length scales become comparable. Far from the critical end point, this is the usual case in a binary liquid, since the interaction ranges responsible for liquid/gas separation and demixing are comparable.
In order to carry further the analytical analysis, we approximate the free energy density $`f(\varphi ,m)`$ (9) by a piecewise quadratic form
$$f(\varphi ,m)=\frac{1}{2}(\varphi \stackrel{~}{\varphi },m\stackrel{~}{m})\genfrac{}{}{0pt}{}{}{f}\left(\begin{array}{c}\varphi \stackrel{~}{\varphi }\\ m\stackrel{~}{m}\end{array}\right)+\sigma \stackrel{~}{\mu },$$
(35)
with three pieces corresponding to the gas phase and the two liquid phases, separated by the lines
$$\varphi _{sep}(m)=\kappa (m^2+m_{+}^{}{}_{}{}^{2}|m|m_+^{})+\stackrel{~}{\mu }/2.$$
(36)
and $`m0`$ at $`\varphi >\varphi _{sep}(0)`$. Here $`\stackrel{~}{\mu }=\mu \mu _c`$, $`\sigma =1`$ for $`\varphi >\varphi _{sep}(m)`$ (gas phase), $`\sigma =+1`$ for $`\varphi <\varphi _{sep}(m)`$ (liquid phases), and the parabolae are adjusted to the leading terms in the expansion of the functional (9) about its minima,
$$\left(\begin{array}{c}\stackrel{~}{\varphi }\\ \stackrel{~}{m}\end{array}\right)=\left(\begin{array}{c}1\\ 0\end{array}\right),\genfrac{}{}{0pt}{}{}{f}=\left(\begin{array}{cc}2& 0\\ 0& 4\kappa \end{array}\right),$$
(37)
for $`\varphi <\varphi _{sep}(m)`$ (gas phase), and
$$\left(\begin{array}{c}\stackrel{~}{\varphi }\\ \stackrel{~}{m}\end{array}\right)=\left(\begin{array}{c}1+\kappa m_{+}^{}{}_{}{}^{2}/2\\ \pm m_+^{}\end{array}\right),\genfrac{}{}{0pt}{}{}{f}=\left(\begin{array}{cc}2& 2\kappa m_+^{}\\ 2\kappa m_+^{}& 2m_{+}^{}{}_{}{}^{2}\end{array}\right),$$
(38)
for $`\varphi >\varphi _{sep}(m)`$ (liquid phases), where the upper sign holds for $`m>0`$, the lower for $`m<0`$. The choice (36) of $`\varphi _{sep}`$ ensures that the potential $`f(m,\varphi )`$ is continuous.
In such a potential, profiles of demixed films correspond to paths in the $`(\varphi ,m)`$ space which can be separated into three parts: (i) moving in one of the liquid regions from $`(\varphi _0,m_0)`$ to $`(\varphi _1,m=0)`$; (ii) following the edge $`(m0)`$ between the two liquid regions from $`(\varphi _1,0)`$ to $`(\varphi _{sep}(0),0)`$; (iii) moving in the gas region from $`(\varphi _{sep}(0),0)`$ to $`(1,0)`$. On principle, a direct transition from (i) to (iii) is conceivable. For the parameters $`\varphi _0`$ of interest, however, such profiles turn out to be energetically less favorable than the profiles which have an intermediate (ii). Profiles of mixed films have two parts (ii) and (iii) only. We shall denote $`l_{(i)}l`$, $`l_{(ii)}`$, and $`l_{(iii)}`$, the length of the slab spent in region (i), (ii) or (iii), respectively.
At given slab length and boundary conditions, the free energy in each of the slabs can be calculated exactly using
$`{\displaystyle _0^l}𝑑z{\displaystyle \frac{1}{2}}\left\{({\displaystyle \frac{du}{dz}})^2+{\displaystyle \frac{u^2}{\lambda ^2}}\right\}`$ (39)
$`=`$ $`{\displaystyle \frac{1}{4}}\left\{(u(0)+u(l))^2\mathrm{tanh}{\displaystyle \frac{l}{2\lambda }}+(u(0)u(l))^2\mathrm{coth}{\displaystyle \frac{l}{2\lambda }}\right\}.`$ (40)
The calculation is straightforward in the regimes (ii) and (iii). In (i), the free energy functional has to be diagonalized first:
$$_{(i)}=\frac{1}{2}_0^l𝑑z\left\{\left[(\frac{dv}{dz})^2+\frac{v^2}{\lambda _1^2}\right]+\left[(\frac{dw}{dz})^2+\frac{w^2}{\lambda _2^2}\right]\right\}$$
with
$`\lambda _{1,2}^2`$ $`=`$ $`1+{\displaystyle \frac{m_{+}^{}{}_{}{}^{2}}{\gamma }}\sqrt{\left({\displaystyle \frac{m_{+}^{}{}_{}{}^{2}}{\gamma }}\right)^2+{\displaystyle \frac{m_{+}^{}{}_{}{}^{2}}{\gamma }}(4\kappa ^22)+1},`$ (41)
$`\left(\begin{array}{c}v\\ w\end{array}\right)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{e^\delta +e^\delta }}}\left(\begin{array}{cc}e^{\delta /2}& e^{\delta /2}\\ e^{\delta /2}& e^{\delta /2}\end{array}\right)\left(\begin{array}{c}\varphi \stackrel{~}{\varphi }\\ \sqrt{\gamma }(m\stackrel{~}{m})\end{array}\right),`$ (48)
where we have defined
$$\delta =\frac{1}{2}\mathrm{ln}(\frac{2\lambda _1^2}{\lambda _2^21}).$$
(50)
The parameter $`\delta `$ or alternatively $`\gamma /m_{+}^{}{}_{}{}^{2}`$ determines the wetting behavior. Fig. 2 shows the two length scales $`\lambda _1`$ and $`\lambda _2`$ as a function of $`\gamma /m_{+}^{}{}_{}{}^{2}`$. The length $`\lambda _1`$ is always larger than $`\lambda _2`$. At $`\gamma /m_{+}^{}{}_{}{}^{2}1`$ or $`\delta 0`$, it characterizes the spatial variations of $`m(z)`$ and grows linearly with $`\gamma /m_{+}^{}{}_{}{}^{2}`$; at $`\gamma /m_{+}^{}{}_{}{}^{2}1`$ or $`\delta 0`$, it characterizes the variations of $`\varphi (z)`$ and remains largely independent of $`\gamma /m_{+}^{}{}_{}{}^{2}`$. These are the limiting regimes discussed in the previous subsection. At $`\gamma /m_{+}^{}{}_{}{}^{2}1`$ or $`\delta 0`$, both $`\lambda _1`$ and $`\lambda _2`$ are related to linear combinations of $`\varphi (z)`$ and $`m(z)`$.
The further calculation proceeds as follows: The free energy in (iii) is given by
$$F_{(i)}=\sqrt{2}/2(1+(\varphi _{sep}(0)+1)^2).$$
(51)
In the region (ii), the result for the free energy is expanded in powers of $`e^{\sqrt{2}l_{(ii)}}`$ up to the second order and minimized with respect to $`l_{(ii)}`$. The free energy calculated in (i) is expanded up to second order in powers of $`e^{l/\lambda _1}`$ and up to first order in $`e^{l/\lambda _2}`$, where $`ll_{(i)}`$. The three contributions are then added up, and the sum is minimized with respect to $`\varphi _1`$ and $`m_0`$ at given surface coupling $`C_m`$. The solution has to be compared with the free energy of a mixed film, which is calculated analogously.
We only report the result for the case $`C_m=0`$ here. The expressions obtained for arbitrary $`C_m`$ are more complicated, but qualitatively similar. Without loss of generality, we can assume $`m>0`$ in the demixed film. As long as $`m_0>0`$, the surface order parameter $`m_0`$ and the free energy difference $`\mathrm{\Delta }F`$ between the mixed and demixed film can then be expanded as
$`\sqrt{\gamma }m_0`$ $`=`$ $`\sqrt{\gamma }m_+^{}+\iota _0+\iota _1e^{l/\lambda _1}+\iota _2e^{l/\lambda _2}+\iota _3e^{2l/\lambda _1}`$ (52)
$`\mathrm{\Delta }F(l)`$ $`=`$ $`\mathrm{\Delta }F_{(ii)}+\tau _0+\tau _2e^{l/\lambda _1}+\tau _2e^{l/\lambda _2}+\tau _3e^{2l/\lambda _1}.`$ (53)
Using the abbreviations
$$K_0=\frac{e^\delta +e^\delta }{\lambda _1\lambda _2},\text{and}K_\pm =\frac{\lambda _1\lambda _2}{e^{\pm \delta }\lambda _1+e^\delta \lambda _2},$$
the coefficients can be written as
$`\iota _0`$ $`=`$ $`(\lambda _2^1\lambda _1^1)K_{}(\varphi _0\stackrel{~}{\varphi })`$ (55)
$`\iota _{1,2}`$ $`=`$ $`2K_0K_+K_{}e^{\pm \delta }\sqrt{\gamma }m_+^{}`$ (56)
$`\iota _3`$ $`=`$ $`2K_0K_{}^2(1+8K_+e^\delta /\lambda _1)(\varphi _0\stackrel{~}{\varphi })`$ (57)
$`\tau _0`$ $`=`$ $`K_0(K_+\gamma m_{+}^{}{}_{}{}^{2}+K_{}(\varphi _0\stackrel{~}{\varphi })^2)`$ (58)
$`\tau _{1,2}`$ $`=`$ $`\pm 2K_0K_+K_{}/\lambda _{2,1}\sqrt{\gamma }m_+^{}(\varphi _0\stackrel{~}{\varphi })`$ (59)
$`\tau _3`$ $`=`$ $`K_0K_{}K_+^2e^\delta (e^\delta /\lambda _2e^\delta /\lambda _1)/\lambda _2\gamma m_{+}^{}{}_{}{}^{2}`$ (61)
$`+K_0K_{}^2e^\delta (1+8K_+e^\delta /\lambda _1)/\lambda _2(\varphi _0\stackrel{~}{\varphi })^2,`$
and with $`h=m_{+}^{}{}_{}{}^{4}(1\kappa ^2)2\stackrel{~}{\mu }`$, $`p=\varphi _0\stackrel{~}{\varphi }+\kappa m_{+}^{}{}_{}{}^{2}`$,
$$\mathrm{\Delta }F_{(ii)}=\frac{1}{\sqrt{2}}p^2\frac{h}{2\sqrt{2}}\left(\mathrm{\hspace{0.25em}1}+\mathrm{ln}\frac{4p^2}{h}\right).$$
(62)
When taking the limits $`\delta \pm \mathrm{}`$, one recovers qualitatively the behavior discussed in the previous section.
The function $`\mathrm{\Delta }F(l)`$ can be conceived as an effective interface potential for the demixed film. The parameters $`\tau _i`$ for a choice of $`\varphi _0`$ ($`\varphi _0=1.5`$) and two values of $`C_m`$, $`C_m=0`$ according to (59) and $`C_m=0.02`$, are shown in Fig. 3. One finds that $`\tau _1`$ is always positive, $`\tau _2`$ is always negative, and $`\tau _3`$ changes sign from positive to negative as $`\gamma /m_{+}^{}{}_{}{}^{2}`$ increases. The leading term of the potential $`F(l)`$ is thus positive, and one expects a first order wetting transition and a prewetting line. On the other hand, the expansions (53) are only valid as long as the surface order parameter $`m_0`$ is positive. According to eqn. (56), the coefficients $`\iota _i`$ of the expansion for $`m_0(l)`$ are negative except for the zeroth order term $`\iota _0`$. Hence $`m_0`$ decreases with film thickness and may vanish at some thickness $`l_c`$. In this case, the film mixes continuously at $`l_c`$, and the prewetting line turns into a second order demixing line sufficiently far from coexistence.
### C Numerical solution
The analytical results of the previous subsection provided insight into the competition of length scales in the binary fluid and the wetting scenarios which can be expected on a wall as a result. However, a reliable calculation of actual phase diagrams, including the details of the prewetting line, is not possible on the basis of the expansion (53). We have thus supplemented the analytical work by a numerical minimization of the functional (4) in the $`\mu \varphi _0`$ plane for selected sets of parameters $`\gamma `$ and $`C_m`$.
The problem is simplified considerably due to the fact that $`\varphi (z)`$ is a monotonic function of $`z`$, i.e., $`m(z)`$ can be expressed as a function $`m(\varphi )`$. The bulk free energy functional in (4) can thus be rewritten as
$``$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑z\left\{{\displaystyle \frac{1}{2}}(1+\gamma ({\displaystyle \frac{dm}{d\varphi }})^2)({\displaystyle \frac{d\varphi }{dz}})^2+f(m(\varphi ),\varphi )\right\}`$
$`=`$ $`{\displaystyle _\varphi _{}^{}^{\varphi _0}}𝑑\varphi \sqrt{1+\gamma ({\displaystyle \frac{dm}{d\varphi }})^2}\sqrt{f(m(\varphi ),\varphi )f(0,\varphi _{}^{})},`$
where the integration constant (15) has been identified and exploited as usual. Minimization with respect to the function $`m(\varphi )`$ leads to the Euler-Lagrange equation
$$2\gamma f(m,\varphi )\frac{d^2m}{d\varphi ^2}=(1+\gamma (\frac{dm}{d\varphi })^2)(\frac{f}{m}\gamma \frac{dm}{d\varphi }\frac{f}{\varphi }),$$
(63)
which we have solved using the Verlet algorithm.
Some results are shown in Figs. 4, 5, 6, and 7. As anticipated in the previous subsection, we find a first order wetting transition, a discontinuous prewetting line and a continuous demixing line. At surface coupling $`C_m>0`$, the demixing line joins the prewetting line in a surface critical end point (Fig. 4). The prewetting line separates a demixed thick film from a mixed thin film (see profiles
in Fig. 5) before reaching the critical end point, then two demixed films of different thickness, and finally vanishes in a critical point. On decreasing the surface coupling $`C_m`$, the critical end point and the critical point move closer to each other, until they merge in a surface tricritical point.
Fig. 6 shows two cases of phase diagrams in the $`\varphi _0\mu `$ plane for $`C_m=0`$ and two different $`\gamma `$ at fixed $`\theta `$, i. e., at fixed bulk order parameter $`m_+^{}`$. With increasing $`\gamma `$, the prewetting line shifts towards larger $`\varphi _0`$ and extends deeper into the off-coexistence region. As $`\gamma \mathrm{}`$, it moves to $`\varphi _0\mathrm{}`$, the film remains mixed and thin at all finite $`\varphi _0`$. At $`\gamma 0`$, on the other hand, the line becomes flat, approaches $`\varphi _+^{}`$, and the tricritical point where it turns into a second order line moves to $`\mu _t\mu _c`$. The numerical results thus agree with the conclusions from section II A.
Fig. 7 demonstrates what happens if instead of making $`\gamma `$ larger, one increases the characteristic length scale of order parameter fluctuations by decreasing $`\theta `$, i.e., approaching the critical end point (reducing $`m_+^{}`$). Far from liquid/vapour coexistence, the transition line still moves towards larger $`\varphi _0`$. However, the effect reverses close to coexistence, the demixing transition is now shifted to lower surface densities $`\varphi _0`$. Furthermore, the length of the prewetting line shrinks instead of growing.
## III Monte Carlo simulations
In this section we describe Monte Carlo simulation studies of the subcritical wetting behaviour of a symmetrical binary fluid at a structureless wall.
### A Model and simulation details
The system we have studied is a symmetrical binary fluid, having interparticle interactions of the Lennard-Jones (LJ) form:
$$u(r_{ij})=4ϵ_{ij}\left[\left(\frac{\sigma _{ij}}{r_{ij}}\right)^{12}\left(\frac{\sigma _{ij}}{r_{ij}}\right)^6\right]$$
(64)
We made the following choice of model parameters: $`\sigma _{11}=\sigma _{22}=\sigma _{12}=\sigma =1`$; $`ϵ_{11}=ϵ_{22}=ϵ`$; $`ϵ_{12}=0.7ϵ`$. i.e. interactions between similar species are treated identically, but those between unlike species are weakened. The inter-particle potential was truncated at a distance of $`R_c=2.5\sigma `$ and no long-range correction or potential shift was applied.
The fluid was confined within a cuboidal simulation cell having dimensions $`P_x\times P_y\times D`$, in the $`x,y`$ and $`z`$ coordinate directions respectively, with $`P_x=P_yP`$. The simulation cell was divided into cubic sub-cells (of size the cutoff $`R_c`$) in order to aid identification of particle interactions. Thus $`P=pR_c`$ and $`D=dR_c`$, with $`p`$ and $`d`$ both integers. To approximate a semi-infinite geometry, periodic boundary conditions were applied in the $`x`$ and $`y`$ directions, while hard walls were applied in the unique $`z`$ direction at $`z=0`$ and $`z=D`$. The hard wall at $`z=0`$ was made attractive, using a potential designed to mimic the long-ranged dispersion forces between the wall and the fluid :
$$V(z)=ϵ_w\left[\frac{2}{15}\left(\frac{\sigma _w}{z}\right)^9\left(\frac{\sigma _w}{z}\right)^3\right]$$
(65)
Here $`z`$ measures the perpendicular distance from the wall, $`ϵ_w`$ is a ‘well-depth’ controlling the interaction strength, and we set $`\sigma _w=1`$. No cutoff was employed and the wall potential was made to act equally on both particle species.
Monte-Carlo simulations of this system were performed using a Metropolis algorithm within the grand canonical ($`\mu ,V,T`$) ensemble . Three types of Monte-Carlo moves were employed:
1. Particle displacements
2. Particle insertions and deletions
3. Particle identity swaps: $`12`$ or $`21`$
To maintain the symmetry of the model, the chemical potentials $`\mu _1`$ and $`\mu _2`$ of the two components were set equal at all times. Thus only one free parameter, $`\mu =\mu _1=\mu _2`$, couples to the overall number density $`\rho =(N_1+N_2)/V`$. The other variables used to explore the wetting phase diagram were the reduced well depth $`ϵ/k_BT`$ and the reduced wall potential $`ϵ_w/k_BT`$. During the simulations, the observables monitored were the total particle density profile
$$\rho (z)=[N_1(z)+N_2(z)]/P^2,$$
(66)
the number difference order parameter profile,
$$n(z)=|N_1(z)N_2(z)|/P^2$$
(67)
These profiles was accumulated in the form of a histogram. Other observables monitored were the total interparticle energy and the wall interaction energy.
The choice of system size was, as ever, a compromise between minimising finite-size effects and maximising computational throughput. Tests showed the profiles to be largely insensitive to the size of the wall area and hence $`p=7`$ was used, this being the largest computationally tractable size consistent with the necessary choice of the slit width $`d`$. The latter must clearly be considerably larger than the film thicknesses of interest in order to prevent the liquid film directly interacting with the hard wall at $`z=D`$. In the results presented below, the typical slit width used was $`d=16`$, corresponding to some $`40`$ molecular diameters. For thin films a narrower slit of width $`d=8`$ was used.
### B Wetting behaviour along a subcritical isotherm
Accurate knowledge of bulk phase behaviour is an essential prerequisite for detailed studies of near-coexistence wetting properties. In the present model, the locus of the liquid vapour coexistence curve and location
of the critical end point are already known to high precision from a previous MC simulation study . The phase diagram in the $`\mu `$-$`T`$ plane (in standard Lennard-Jones reduced units ) is shown in fig. 8. The critical end point is located at $`T_{cep}=0.958(3),\mu _{cep}=3.017(3)`$ . We note that although the locus of the coexistence curve is known to five significant figures, the position of the CEP along this tightly determined line is known only to three significant figures.
To determine the wetting properties at temperatures below $`T_{cep}`$, the number density profile $`\rho (z)`$ was studied along the isotherm $`T=0.9467`$ as coexistence was approached from the vapour side. To achieve this, the chemical potential was incremented up to its coexistence value $`\mu _{cx}(T)`$ in a sequence of $`6`$-$`10`$ steps of constant size $`\mathrm{\Delta }\mu =0.0025`$. This procedure was repeated for a number of different values of the wall-fluid potential strength $`ϵ_w`$, allowing the influence of this parameter on the wetting behaviour to be ascertained. In all, six values of the $`ϵ_w`$ were studied ($`ϵ_w=1.0,1.7,1.75,2.0,3.0,4.0`$). We describe the wetting behaviour for each in turn.
For $`ϵ_w=1.0`$, fig. 9 shows that although the film thickness grows very slightly as coexistence is approached, it never exceed two molecular diameters. At no point in the profile does the density attain that of the liquid phase ($`\rho 0.6`$). The presence of a thin wetting layer right up to coexistence implies incomplete (partial) wetting.
Increasing the wall potential to $`ϵ_w=1.70`$ \[fig. 10\], results in considerably more structure near the wall compared to $`ϵ_w=1.0`$, with clear density oscillations arising from excluded-volume ‘packing effects’ . The profile is much more responsive to changes in the chemical potential and reaches a thickness of 4-5 molecular diameters close to coexistence.
For $`ϵ_w=1.75`$, however, the situation changes qualitatively, as shown in fig. 11. On increasing the chemical potential, a clear jump is observed in both the thickness of the film, and the value of its density. In the thick film, the density of a significant portion of the film is that of
the bulk liquid. This thin-thick jump constitutes a prewetting transition, as previously observed in simulation studies of lattice gas models , Lennard-Jones fluids as well as experimentally .
As the wall potential is increased to $`ϵ_w=2.0`$ (fig. 12), the sharp prewetting transition is lost and instead the film thickness increases smoothly as $`\mu `$ approaches its coexistence value. This suggests that here the system is above the prewetting critical point ().
On increasing $`ϵ_w`$ to $`3.0`$, a new feature emerges (fig. 13(a)). As the chemical potential increases, the thickness of the film initially increases smoothly with increasing $`\mu `$. However, once the thickness reaches some $`10`$ molecular diameters, a large jump occurs to a thickness of about $`15`$ molecular diameters. Concomitant with this jump is a demixing of the film as a whole, as seen in the order parameter profile fig. 13(b). The size of the jump in the layer thickness appears to reduce as the wall strength is increased to $`ϵ_w=4.0`$ (fig. 14), suggesting a weakening of the transition.
## IV Discussion
The Monte Carlo simulation results at subcritical temperatures provide evidence that the mean field calculations correctly identify the qualitative wetting behaviour. They show that depending on the fluid-wall interaction strength $`ϵ_w`$, a number of different wetting scenarios occur as liquid-vapour coexistence is approached from the vapour side. At small $`ϵ_w`$, only a very thin film builds up on the wall. For intermediate values of $`ϵ_w`$, a first prewetting transition is observed from a thin mixed film to a thick liquidlike mixed layer. Further increasing $`ϵ_w`$ induces a second prewetting transition between a mixed liquidlike layer and a thicker demixed film, the situation being very similar to that shown in figure 5. The abrupt, first order, character of this latter transition appears to weaken on further increasing $`ϵ_w`$, in accord with the theoretical predictions.
We will now attempt to set our results within the context of the bulk phase diagram of the binary liquid. To this end, we discuss the possible wetting scenarios in the vicinity of the critical end point $`T_{cep}`$. As previously argued in the introduction, for temperatures $`T<T_{cep}`$ sufficiently close to $`T_{cep}`$, the bulk correlation length $`\xi `$ of the demixed liquid is larger than the thickness $`l_{}`$ of a mixed liquid layer at the wall. The state of order of the film thus depends strongly on the boundary conditions of the two interfaces confining the liquid layer. The nonselective liquid-vapor interface always favors mixing due to the reduced number of interacting neighbors in the interfacial region. The liquid-substrate interface, on the other hand, can either favour mixing or demixing depending on the strength of the fluid-wall potential. For a weakly attractive wall potential, mixing is favoured because the particle density at the wall is low and the presence of the wall reduces the number of interacting neighbours. For a strongly attractive wall, however, the high density at the wall can counteract the missing neighbour effect leading to an overall demixing tendency.
If the net effect favours mixing at the wall, a continuous demixing of the layer as coexistence is approached can be excluded. A first order transition involving a discontinuous increase of the film thickness upon demixing is still conceivable. However, we have shown in section II A, that (at the mean field level, at least) the demixed wetting film has a higher free energy than the corresponding mixed film provided the correlation length of composition fluctuations is sufficiently large.
At walls which suppress demixing, the film is thus always mixed close to the critical end point, and its thickness $`l_{}`$ below the critical end point is finite. Hence the critical end point is automatically a critical wetting point. The resulting phase diagram is shown schematically in fig. 15 (a). Note that the wetting transition here is pinned by a bulk phase transition, a situation somewhat reminiscent of triple-point wetting .
The situation changes if the substrate favours demixing. In this situation, one component segregates to the surface of the film already slightly above $`T_{cep}`$, and the order propagates continuously into the bulk of the film at $`T_{cep}`$. The film remains wet at $`T_{cep}`$. From the results of section II C (in particular, Fig. 7), one can deduce two possible scenarios. The film may still exhibit a first order wetting transition to a nonwet state at a temperature below $`T_{cep}`$ (e.g. in Fig. 7) at $`\varphi _0=1.14`$). The discontinuous phase transition at liquid/vapour coexistence then spawns a prewetting line which eventually switches into a second order demixing line and loops around the critical end point as suggested in Fig. 15 (b). If the wall is strongly attractive (e.g. at $`\varphi _0=1.27`$) in Fig. 7), the wall wets at all temperatures, the prewetting line detaches from the coexistence line and is continued by second order demixing lines both at the high and low temperature side as sketched in Fig. 15 (c).
### Acknowledgements
NBW thanks the Royal Society of Edinburgh, the EPSRC (grant no. GR/L91412) and the British Council for financial support.
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# From the solution of the Tsarev system to the solution of the Whitham equations
## 1 Introduction
The Whitham equations are a collection of quasi-linear hyperbolic systems of the form ,,
$$\frac{u_i}{t}\lambda _i(u_1,u_2,\mathrm{},u_{2g+1})\frac{u_i}{x}=0,x,t,u_i\mathrm{I}\mathrm{R},i=1,\mathrm{},2g+1,g=0,1,2,\mathrm{},$$
(1.1)
with the ordering $`u_1>u_2>\mathrm{}>u_{2g+1}`$. For a given $`g`$ the system (1.1) is called $`g`$-phase Whitham equations. Because of the diagonal form of systems (1.1) the dependent variables $`u_1>u_2>\mathrm{}>u_{2g+1}`$ are called Riemann invariants. For $`g>0`$ the speeds $`\lambda _i(u_1,u_2,\mathrm{},u_{2g+1})`$, $`i=1,2,\mathrm{},2g+1`$, depend through $`u_1,\mathrm{},u_{2g+1}`$ on complete hyperelliptic integrals of genus $`g`$. For this reason the $`g`$-phase system is also called genus $`g`$ system. The zero-phase Whitham equation has the form
$$\frac{u}{t}6u\frac{u}{x}=0,$$
(1.2)
where we use the notation $`u_1=u`$. In this paper we study the initial value problem of the Whitham equations for monotone increasing smooth ($`C^{\mathrm{}}`$) initial data $`u(x,t=0)=u_0(x)`$, where the range of $`u_0(x)`$, $`x\mathrm{I}\mathrm{R}`$, is the interval $`(a,b)`$ and $`\mathrm{}a<b+\mathrm{}`$.
The initial value problem consists of the following. We consider the evolution on the $`xu`$ plane of the initial curve $`u(x,t=0)=u_0(x)`$ according to the zero-phase equation (1.2). The solution $`u(x,t)`$ of (1.2) with the initial data $`u_0(x)`$ is given by the characteristic equation
$$x=6tu+f(u)$$
(1.3)
where $`f(u)`$ is the inverse function of $`u_0(x)`$. The solution $`u(x,t)`$ in (1.3) is globally well defined only for $`0t<t_0`$, where $`t_0=\frac{1}{6}\mathrm{min}_{u\mathrm{I}\mathrm{R}}[f^{}(u)]`$ is the time of gradient catastrophe of (1.3). Near the point of gradient catastrophe and for a short time $`t>t_0`$, the evolving curve is given by a multivalued function with three branches $`b>u_1(x,t)>u_2(x,t)>u_3(x,t)>a`$, which evolve according to the one-phase Whitham equations.
Outside the multivalued region the solution is given by the zero-phase solution $`u(x,t)`$ defined in (1.3). On the phase transition boundary the zero-phase solution and the one-phase solution are attached $`C^1`$-smoothly (see Fig1.1).
Since the Whitham equations are hyperbolic, other points of gradient catastrophe can appear in the branches $`u_1(x,t)>u_2(x,t)>u_3(x,t)`$ themselves or in $`u(x,t)`$.
In general, for $`t>t_0`$, the evolving curve is given by a multivalued function with and odd number of branches $`b>u_1(x,t)>u_2(x,t)>\mathrm{}>u_{2g+1}(x,t)>a`$, $`g0`$. These branches evolve according to the $`g`$-phase Whitham equations. The $`g`$-phase solutions for different $`g`$ must be glued together in order to produce a $`C^1`$-smooth curve in the $`(x,u)`$ plane evolving smoothly with $`t`$ (see Fig. 1.2). The initial value problem of the Whitham equations is to determine, for almost all $`t>0`$ and $`x`$, the phase $`g(x,t)0`$ and the corresponding branches $`u_1(x,t)>u_2(x,t)>\mathrm{}>u_{2g+1}(x,t)`$ from the initial data $`x=f(u)|_{t=0}`$.
For example for the initial data $`x=x_0+6t_0u+(uu_0)^3`$, $`t_00`$, the solution of the Whitham equations is of genus at most equal to one . It is of genus one inside the cusp $`12\sqrt{3}(tt_0)^{\frac{3}{2}}<xx_0<\frac{4}{3}\sqrt{\frac{5}{3}}(tt_0)^{\frac{3}{2}}`$ ; it is of genus zero outside this cusp. The point $`x=x_0,t=t_0`$ is the point of gradient catastrophe of the zero-phase solution. The curve $`x^{}(t)=x_012\sqrt{3}(tt_0)^{\frac{3}{2}}`$, $`t>t_0`$, describes the locus of points where $`u_1(x,t)=u_2(x,t)`$. The curve $`x^+(t)=x_0+\frac{4}{3}\sqrt{\frac{5}{3}}(tt_0)^{\frac{3}{2}}`$, $`t>t_0`$, describes the locus of points where $`u_2(x,t)=u_3(x,t)`$. For generic initial data it is not known whether the genus of the solution of the Whitham equations is bounded. We say that the Cauchy problem for the Whitham equations has a global solution if the genus $`g<\mathrm{}`$ for all $`x`$ and $`t0`$.
Using the geometric-Hamiltonian structure of the Whitham equations, Tsarev showed that these equations can be locally integrated by a generalization of the method of characteristic. Namely he proved that if the functions $`w_i=w_i(u_1,u_2,\mathrm{},u_{2g+1})`$, $`i=1,\mathrm{},2g+1`$, solve the linear over-determined system
$$\frac{w_i}{u_j}=\frac{1}{\lambda _i\lambda _j}\frac{\lambda _i}{u_j}[w_iw_j],i,j=1,2,\mathrm{},2g+1,ij,$$
(1.4)
where $`\lambda _i=\lambda _i(u_1,u_2,\mathrm{},u_{2g+1}),i=1,\mathrm{},2g+1`$, are the speeds in (1.1), then the solution $`\stackrel{}{u}(x,t)=(u_1(x,t),u_2(x,t),\mathrm{},u_{2g+1}(x,t))`$ of the so called hodograph transform
$$x=\lambda _i(\stackrel{}{u})t+w_i(\stackrel{}{u})i=1,\mathrm{},2g+1,$$
(1.5)
satisfies system (1.1). Conversely, any solution $`(u_1(x,t),u_2(x,t),\mathrm{},u_{2g+1}(x,t))`$ of (1.1) can be obtained in this way.
Tsarev theorem relies on two factors:
a) the existence of a solution of the linear over-determined system (1.4);
b) the existence of a real solutions $`u_1(x,t)>u_2(x,t)>\mathrm{}>u_{2g+1}(x,t)`$ of the hodograph transform (1.5).
In this paper we solve completely problem a), namely we build for any smooth monotone increasing initial data $`x=f(u)|_{t=0}`$ the solution of system (1.4) for any $`g0`$. The solution $`w_i(\stackrel{}{u})`$, $`i=1,\mathrm{},2g+1`$, of genus $`g`$ satisfies some natural boundary conditions which guarantee its uniqueness.
Regarding problem b) we characterize initial data such that the solution of the Whitham equations exists only for $`gN`$ where $`N`$ is a positive integer. Namely we show that if the initial data satisfies the condition
$$\frac{d^{2N+1}}{du^{2N+1}}f(u):=f^{(2N+1)}(u)>0,1N\mathrm{I}\mathrm{N}$$
(1.6)
for all real $`u`$ belonging to the domain of $`f`$ except one point, then the solution of the Whitham equations has genus at most $`N`$ for any $`x`$ and $`t0`$. For $`N=1`$ this result has already been proved by Tian .
The investigation of the initial value problem of the Whitham equations was initiated by Gurevich and Pitaevskii . In the case $`g1`$ they solved the initial value problem of system (1.1) for step-like initial data and studied numerically the case of cubic initial data.
The initial value problem of the Whitham equations was deeply studied by Tian. In he constructed the general solution of the Tsarev system (1.4) for $`0g1`$ and for smooth monotone increasing initial data. He also proved the solvability of the hodograph transform (1.5) for $`0g1`$. In he obtained the solution of the Tsarev system for $`g>0`$ for polynomial initial data.
The equations (1.1) were found by Whitham in the single phase case $`g=1`$ and more generally by Flaschka, Forest and McLaughlin in the multi-phase case. The Whitham equations were also found in when studying the zero dispersion limit of the Korteweg de Vries equation. The hyperbolic nature of the equations was found by Levermore .
This paper is organized as follows.
In Sec. 2 we give some background about Abelian differentials on hyperelliptic Riemann surfaces.
In Sec. 3 we describe the Whitham equations.
In Sec. 4 we build the solution of the Tsarev system (1.4) for smooth monotone increasing initial data.
In Sec. 5 we show that under the hypothesis (1.6) the hodograph transform is solvable for $`gN`$.
In Sec. 6 we draw the conclusions.
## 2 Riemann surfaces and Abelian differentials: notations and definitions
Let
$$𝒮_g:=\left\{P=(r,\mu ),\mu ^2=\underset{j=1}{\overset{2g+1}{}}(ru_j)\right\},u_1>u_2>\mathrm{}>u_{2g+1},u_i\mathrm{I}\mathrm{R},$$
(2.1)
be the hyperelliptic Riemann surface of genus $`g0`$. We shall use the standard representation of $`𝒮_g`$ as a two-sheeted covering of $`C\mathrm{I}\mathrm{P}^1`$ with cuts along the intervals
$$[u_{2k},u_{2k1}],k=1,\mathrm{},g+1,u_{2g+2}=\mathrm{}.$$
We choose the basis $`\{\alpha _j,\beta _j\}_{j=1}^g`$ of the homology group $`H_1(\mathrm{\Gamma }_g)`$ so that $`\alpha _j`$ lies fully on the upper sheet and encircles clockwise the interval $`[u_{2j},u_{2j1}]`$, $`j=1,\mathrm{},g`$, while $`\beta _j`$ emerges on the upper sheet on the cut $`[u_{2j},u_{2j1}]`$, passes anti-clockwise to the lower sheet trough the cut $`(\mathrm{},u_{2g+1}]`$ and return to the initial point through the lower sheet.
The one-forms that are analytic on the closed Riemann surface $`𝒮_g`$ except for a finite number of points are called Abelian differentials.
We define on $`𝒮_g`$ the following differentials :
1) The canonical basis of holomorphic one-forms or Abelian differentials of the first kind $`\varphi _1,\varphi _2\mathrm{}\varphi _g`$:
$$\varphi _k(r)=\frac{r^{g1}\gamma _1^k+r^{g2}\gamma _2^k+\mathrm{}+\gamma _g^k}{\mu (r)}dr,k=1,\mathrm{},g.$$
(2.2)
The constants $`\gamma _i^k`$ are uniquely determined by the normalization conditions
$$_{\alpha _j}\varphi _k=\delta _{jk},j,k=1,\mathrm{},g.$$
(2.3)
We remark that an holomorphic differential having all its $`\alpha `$-periods equal to zero is identically zero .
2) The set $`\sigma _k^g`$, $`k0`$, $`g0`$, of Abelian differentials of the second kind with a pole of order $`2k+2`$ at infinity, with asymptotic behavior
$$\sigma _k^g(r)=\left[r^{k\frac{1}{2}}+O(r^{\frac{3}{2}})\right]dr\text{for large}r$$
(2.4)
and normalized by the condition
$$_{\alpha _j}\sigma _k^g=0,j=1,\mathrm{},g.$$
(2.5)
We use the notation
$$\sigma _0^g(r)=dp^g(r),12\sigma _1^g(r)=dq^g(r)g0.$$
(2.6)
In literature the differentials $`dp^g(r)`$ and $`dq^g(r)`$ are called quasi-momentum and quasi-energy respectively . The explicit formula for the differentials $`\sigma _k^g`$, $`k0`$, is given by the expression
$`\sigma _k^g(r)={\displaystyle \frac{P_k^g(r)}{\mu (r)}}dr,P_k^g(r)=r^{g+k}+a_1^kr^{g+k1}+a_2^kr^{g+k2}\mathrm{}+a_{g+k}^k,`$ (2.7)
where the coefficients $`a_i^k=a_i^k(\stackrel{}{u})`$, $`\stackrel{}{u}=(u_1,u_2,\mathrm{},u_{2g+1})`$, $`i=1,\mathrm{},g+k,`$ are uniquely determined by (2.4) and (2.5).
3) The Abelian differential of the third kind $`\omega _{qq_0}(r)`$ with first order poles at the points $`Q=(q,\mu (q))`$ and $`Q_0=(q_0,\mu (q_0))`$ with residues $`\pm 1`$ respectively. Its periods are normalized by the relation
$$_{\alpha _j}\omega _{qq_0}(r)=0,j=1,\mathrm{},g.$$
(2.8)
### 2.1 Riemann bilinear relations
Let $`\omega _1`$ and $`\omega _2`$ be two Abelian differentials on the Riemann surface $`𝒮_g`$. If all the residues of $`\omega _1`$ and $`\omega _2`$ are equal to zero, then the integrals $`d^1\omega _1`$ and $`d^1\omega _2`$ do not have logarithm singularities on $`𝒮_g`$. If the differential $`\omega _1`$ has non zero residue, then its integral has logarithm singularities. Let $`s`$ be the path connecting the singular points of $`d^1\omega _1`$. We have the following relation.
$$\underset{j=1}{\overset{g}{}}\left[_{\alpha _j}\omega _1_{\beta _j}\omega _2_{\alpha _j}\omega _2_{\beta _j}\omega _1\right]+_s\mathrm{\Delta }(d^1\omega _1)\omega _2=2\pi i\underset{𝒮_gs}{}Res[(d^1\omega _1)\omega _2],$$
(2.9)
where $`\mathrm{\Delta }(d^1\omega _1)`$ is the difference of the values of $`d^1\omega _1`$ on the two sides of the cut $`s`$ and the quantity $`_{𝒮_gs}Res[(d^1\omega _1)\omega _2]`$ is the sum of the residues of the differential $`(d^1\omega _1)\omega _2`$ on the cut surface $`𝒮_gs`$. This formula is known as the Riemann bilinear period relation .
Assuming $`\omega _1=\omega _{qq_0}`$ and $`\omega _2=\varphi _k`$ in (2.9) we obtain
$$_{\beta _k}\omega _{qq_0}=2\pi i_{q_0}^q\varphi _k,k=1,\mathrm{},g.$$
(2.10)
Assuming $`\omega _1=\omega _{qq_0}`$ and $`\omega _2=\omega _{pp_0}`$ in (2.9) we obtain
$$_{p_0}^p\omega _{qq_0}=_{q_0}^q\omega _{pp_0}.$$
(2.11)
Differentiating with respect to $`p`$ and $`q`$ the above expression we obtain the identity
$$d_q[\omega _{qq_0}(p)]=d_p[\omega _{pp_0}(q)],$$
(2.12)
where $`d_q`$ and $`d_p`$ denote differentiation with respect to $`q`$ and $`p`$ respectively. From the expression (2.11) it follows that $`\omega _{qq_0}(r)`$ is a many-value analytic function of the variable $`q`$. The many-value character of $`\omega _{qq_0}(r)`$ as a function of $`q`$ can be described by the equations
$$_{\alpha _k}d_q[\omega _{qq_0}(r)]=0,_{\beta _k}d_q[\omega _{qq_0}(r)]=2\pi i\varphi _k(r),k=1,\mathrm{},g,$$
(2.13)
In the following we mainly use the normalized differential $`\omega _z^g(r)`$ which has simple poles at the points $`Q^\pm (z)=(z,\pm \mu (z))`$ with residue $`\pm 1`$ respectively.
The differential $`\omega _z^g(r)`$ is explicitly given by the expression
$$\omega _z^g(r)=\frac{dr}{\mu (r)}\frac{\mu (z)}{rz}\underset{k=1}{\overset{g}{}}\varphi _k(r)_{\alpha _k}\frac{dt}{\mu (t)}\frac{\mu (z)}{tz},$$
(2.14)
where $`\varphi _k(r)`$, $`k=1,\mathrm{},g`$, is the normalized basis of holomorphic differentials. $`\omega _z^g(r)`$ as a function of $`z`$, is an Abelian integral having poles of first order at the points $`Q^\pm (r)=(r,\pm \mu (r))`$. The periods of this integral are obtained from the relations (2.13)
$$_{\alpha _j}d_z[\omega _z^g(r)]=0,_{\beta _j}d_z[\omega _z^g(r)]=4\pi i\varphi _k(r),j=1,\mathrm{}g.$$
(2.15)
We apply the Riemann bilinear relation (2.9) to the differentials $`\sigma _m^g(r)`$ and $`\omega _z^g(r)`$ getting
$$\begin{array}{cc}\hfill _{Q^{}(z)}^{Q^+(z)}\sigma _m^g(\xi )=& \underset{[}{Res}r=\mathrm{}][\omega _z^g(r)d^1\sigma _m^g(r)],m=0,\mathrm{},g,\hfill \\ \hfill =& \frac{4}{2m+1}\left(\mu (z)ϵ_{mg}+\underset{k=1}{\overset{g}{}}\underset{j=1}{\overset{g}{}}\gamma _j^k\mathrm{\Gamma }_{m+1j}_{\alpha _k}\frac{dt}{\mu (t)}\frac{\mu (z)}{tz}\right).\hfill \end{array}$$
(2.16)
In the above formula $`ϵ_{mg}=1`$ for $`m=g`$ and zero otherwise, the coefficients $`\gamma _j^k`$ have been defined in (2.2) and the $`\mathrm{\Gamma }_l`$’s are the coefficients of the expansion for $`\xi \mathrm{}`$ of
$$\frac{1}{\mu (\xi )}=\xi ^{g\frac{1}{2}}(\mathrm{\Gamma }_0+\frac{\mathrm{\Gamma }_1}{\xi }+\frac{\mathrm{\Gamma }_2}{\xi ^2}+\mathrm{}+\frac{\mathrm{\Gamma }_l}{\xi ^l}+\mathrm{}).$$
(2.17)
We define $`\mathrm{\Gamma }_k=0`$ for $`k<0`$.
Inverting (2.16) and introducing the quantities $`N_j(z,\stackrel{}{u})=_{k=1}^g\gamma _j^k_{\alpha _k}{\displaystyle \frac{dt}{\mu (t)}}{\displaystyle \frac{\mu (z)}{tz}}`$ we obtain
$$\left(\begin{array}{c}N_1(z,\stackrel{}{u})\\ N_2(z,\stackrel{}{u})\\ \mathrm{}\\ N_g(z,\stackrel{}{u})\\ \mu (z)\end{array}\right)=\left(\begin{array}{ccccc}\stackrel{~}{\mathrm{\Gamma }}_0& 0& 0& \mathrm{}& 0\\ \stackrel{~}{\mathrm{\Gamma }}_1& \stackrel{~}{\mathrm{\Gamma }}_0& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \stackrel{~}{\mathrm{\Gamma }}_{g1}& \stackrel{~}{\mathrm{\Gamma }}_{g2}& \mathrm{}& \stackrel{~}{\mathrm{\Gamma }}_0& 0\\ \stackrel{~}{\mathrm{\Gamma }}_g& \stackrel{~}{\mathrm{\Gamma }}_{g1}& \mathrm{}& \stackrel{~}{\mathrm{\Gamma }}_1& \stackrel{~}{\mathrm{\Gamma }}_0\end{array}\right)\left(\begin{array}{c}\frac{1}{4}_{Q^{}(z)}^{Q^+(z)}\sigma _0^g(\xi )\\ \frac{3}{4}_{Q^{}(z)}^{Q^+(z)}\sigma _1^g(\xi )\\ \mathrm{}\\ \frac{2g1}{4}_{Q^{}(z)}^{Q^+(z)}\sigma _{g1}^g(\xi )\\ \frac{2g+1}{4}_{Q^{}(z)}^{Q^+(z)}\sigma _g^g(\xi )\end{array}\right)$$
(2.18)
where the $`\stackrel{~}{\mathrm{\Gamma }}_k`$’s are the coefficients of the expansion for $`\xi \mathrm{}`$ of
$$\mu (\xi )=\xi ^{g+\frac{1}{2}}(\stackrel{~}{\mathrm{\Gamma }}_0+\frac{\stackrel{~}{\mathrm{\Gamma }}_1}{\xi }+\frac{\stackrel{~}{\mathrm{\Gamma }}_2}{\xi ^2}+\mathrm{}+\frac{\stackrel{~}{\mathrm{\Gamma }}_l}{\xi ^l}+\mathrm{}).$$
(2.19)
Using (2.18) the differential $`\omega _z^g(r)`$ turns out to be given by the relation
$$\omega _z^g(r)=\frac{dr}{\mu (r)}\frac{\mu (z)}{rz}\underset{k=1}{\overset{g}{}}N_k(z,\stackrel{}{u})\frac{r^{gk}}{\mu (r)}dr.$$
(2.20)
We remark that $`\omega _z^g(r)`$ is a multivalued function of $`z`$ and it is regular at infinity.
From the relation (2.18) we obtain the identity which will be useful later
$$\mu (z)=\frac{1}{4}\underset{k=1}{\overset{g+1}{}}(2k1)\stackrel{~}{\mathrm{\Gamma }}_{g+1k}_{Q^{}(z)}^{Q^+(z)}\sigma _{k1}^g(\xi ).$$
(2.21)
The next proposition is also important for our subsequent considerations.
###### Proposition 2.1
The Abelian differentials of the second kind $`\sigma _k^g(r)`$, $`k0`$, defined in (2.4) satisfy the relations
$$\sigma _k^g(r)=\frac{1}{2}\underset{z=\mathrm{}}{Res}\left[\omega _z^g(r)z^{k\frac{1}{2}}dz\right]=\frac{1}{2k+1}d_r\underset{z=\mathrm{}}{Res}\left[\omega _r^g(z)z^{k+\frac{1}{2}}\right],$$
(2.22)
where $`\omega _z^g(r)`$ has been defined in (2.14), $`\omega _r^g(z)`$ is the normalized Abelian differential of the third kind with simple poles at the points $`Q^\pm (r)=(r,\pm \mu (r))`$ with residue $`\pm 1`$ respectively and $`d_r`$ denotes differentiation with respect to $`r`$.
Proof: the differential $`\underset{z=\mathrm{}}{Res}\left[\omega _z^g(r)z^{k\frac{1}{2}}dz\right]`$ is normalized because $`\omega _z^g(r)`$ is a normalized differential. Furthermore
$$\underset{z=\mathrm{}}{Res}\left[\omega _z^g(r)z^{k\frac{1}{2}}dz\right]=r^{k\frac{1}{2}}dr+O(r^{\frac{3}{2}})dr\text{for}r\mathrm{}.$$
Therefore $`\underset{z=\mathrm{}}{Res}\left[\omega _z^g(r)z^{k\frac{1}{2}}dz\right]`$ coincides with the normalized Abelian differential of the second kind $`\sigma _k^g(r)`$. For proving the second equality in (2.22) we consider the integral in the $`z`$ variable
$$0=_C_{\mathrm{}}d_z(\omega _z^g(r)z^{k+\frac{1}{2}})=_C_{\mathrm{}}z^{k+\frac{1}{2}}(d_z\omega _z^g(r))+_C_{\mathrm{}}(k\frac{1}{2})(\omega _z^g(r)z^{k+\frac{1}{2}}),$$
where $`C_{\mathrm{}}`$ is a close contour around the point at infinity. Substituting the identity $`d_z\omega _z^g(r)=d_r\omega _r^g(z)`$ in the right hand side of the above relation we obtain the second relation in (2.22). $`\mathrm{}`$
## 3 Preliminaries on the theory of the Whitham equations
The speeds $`\lambda _i(u_1,u_2,\mathrm{},u_{2g+1})`$ of the $`g`$-phase Whitham equations (1.1) are given by the ratio ,:
$$\lambda _i(\stackrel{}{u})=\frac{dq^g(r)}{dp^g(r)}|_{r=u_i},i=1,2,\mathrm{},2g+1,$$
(3.1)
where $`dp^g(r)`$ and $`dq^g(r)`$ have been defined in (2.6). In the case $`g=0`$
$$dp^0(r)=\frac{dr}{\sqrt{ru}},dq^0(r)=\frac{12r6u}{\sqrt{ru}}dr,$$
(3.2)
so that one obtains the zero-phase Whitham equation (1.2).
For monotone increasing smooth initial data $`x=f(u)|_{t=0}`$, the solution of the zero-phase equation (1.2) is obtained by the method of characteristic and is given by the expression
$$x=6tu+f(u).$$
(3.3)
The zero-phase solution is globally well defined only for $`0t<t_0`$ where $`t_0=\frac{1}{6}\mathrm{min}_{u\mathrm{I}\mathrm{R}}[f^{}(u)]`$ is the time of gradient catastrophe of the solution (3.3). The breaking is cause by an inflection point in the initial data. For $`tt_0`$ we expect to have single, double and higher phase solutions. For higher genus the Whitham equations can be locally integrated using a generalization of the characteristic equation (3.3). We have the following theorem of Tsarev
###### Theorem 3.1
If $`w_i(\stackrel{}{u})`$, $`\stackrel{}{u}=(u_1,u_2,\mathrm{},u_{2g+1})`$, solves the linear over-determined system
$$\frac{w_i}{u_j}=a_{ij}(\stackrel{}{u})[w_iw_j],i,j=1,2,\mathrm{},2g+1,ij,$$
(3.4)
where
$$a_{ij}=\frac{1}{\lambda _i\lambda _j}\frac{\lambda _i}{u_j}i,j=1,2,\mathrm{},2g+1,ij,$$
(3.5)
then the solution $`(u_1(x,t),u_2(x,t),\mathrm{},u_{2g+1}(x,t))`$ of the hodograph transformation
$$x=\lambda _i(\stackrel{}{u})t+w_i(\stackrel{}{u})i=1,\mathrm{},2g+1,$$
(3.6)
satisfies system (1.1). Conversely, any solution $`(u_1,u_2,\mathrm{},u_{2g+1})`$ of (1.1) can be obtained in this way.
To guarantee that the $`g`$phase solutions for different $`g`$ are attached continuously, the following natural boundary conditions must be imposed on $`w_i(u_1,u_2,\mathrm{},u_{2g+1})`$, $`i=1,\mathrm{},2g+1`$.
When $`u_l=u_{l+1}`$, $`\mathrm{\hspace{0.33em}1}l2g`$,
$$w_l^g(u_1,\mathrm{},u_l,u_l,\mathrm{},u_{2g+1})=w_{l+1}^g(u_1,\mathrm{},u_l,u_l,\mathrm{},u_{2g+1})$$
(3.7)
and for $`1i2g+1,il,l+1`$
$`w_i^g(u_1,\mathrm{},u_l,u_l,\mathrm{},u_{2g+1})=w_i^{g1}(u_1,\mathrm{},\widehat{u}_l,\widehat{u}_l,\mathrm{},u_{2g+1}).`$ (3.8)
The sup-scripts $`g`$ and $`g1`$ in the $`w_i`$’s specify the corresponding genus and the hat denotes the variable that have been dropped. When $`g=1`$ and $`u_2=u_3`$ we have that
$`\begin{array}{ccc}& w_1(u_1,u_3,u_3)=f(u_1)\hfill & \\ & w_2(u_1,u_3,u_3)=w_3(u_1,u_3,u_3),\hfill & \end{array}`$ (3.11)
where $`f(u)`$ is the initial data. Similar conditions hold true when $`u_1=u_2`$, namely
$`\begin{array}{ccc}& w_3(u_1,u_1,u_3)=f(u_3)\hfill & \\ & w_1(u_1,u_1,u_3)=w_2(u_1,u_1,u_3).\hfill & \end{array}`$ (3.14)
We remark that the $`\lambda _i(\stackrel{}{u})`$’s satisfy the boundary conditions (3.7-3.8) and for $`g=1`$ we have
$$\lambda _1(u_1,u_3,u_3)=6u_1,\lambda _2(u_1,u_3,u_3)=\lambda _3(u_1,u_3,u_3),$$
and
$$\lambda _3(u_1,u_1,u_3)=6u_3,\lambda _1(u_1,u_1,u_3)=\lambda _2(u_1,u_1,u_3).$$
The solution of the boundary value problem (3.4), (3.7-3.14) has been obtained in for monotone increasing analytic initial data of the form
$$x=f_a(u)=c_0+c_1u+\mathrm{}+c_ku^k+\mathrm{}$$
(3.15)
where we assume that only a finite number of $`c_k`$ is different from zero. For such initial data the $`w_i(\stackrel{}{u})`$’s which satisfy (3.4) and the boundary conditions (3.7-3.14) are given by the expression
$$w_i(\stackrel{}{u})=\frac{ds^g(r)}{dp^g(r)}|_{r=u_i},i=1,\mathrm{},2g+1.$$
(3.16)
The differential $`ds^g`$ in (3.16) is given by
$$ds^g(r)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{2^kk!}{(2k1)!!}c_k\sigma _k^g(r),$$
(3.17)
and the differentials $`\sigma _k^g(r)`$, $`k0`$ have been defined in (2.7).
The solution (3.6) of the $`g`$-phase Whitham equations can also be written in an equivalent algebro-geometric form namely
$$(xdp^g(r)tdq^g(r)+ds^g(r))_{r=u_i}=0,i=1,2\mathrm{},2g+1,$$
(3.18)
where $`dp^g`$, $`dq^g`$ and $`ds^g`$ have been defined in (2.6) and (3.17) respectively.
The solution $`u_1>u_2>\mathrm{}>u_{2g+1}`$ of the $`g`$-phase Whitham equations (1.1) is implicitly defined as a function of $`x`$ and $`t`$ by the equations (3.6) or (3.18). The solution is uniquely defined only for those $`x`$ and $`t`$ such that the functions $`u_i(x,t)`$ are real and $`_xu_i(x,t)`$, $`i=1,\mathrm{},2g+1`$, are not vanishing.
One of the problems in the theory of the Whitham equations is to determine when (3.6) or (3.18) are solvable for real $`u_1,\mathrm{},u_{2g+1}`$ as functions of $`x`$ and $`t`$. This problem has been solved by Tian for $`g1`$.
###### Theorem 3.2
Consider a monotone increasing initial data $`x=f(u)|_{t=0}`$. Suppose that $`u^{}`$ is the inflection point of $`f(u)`$ that causes the breaking of the zero-phase solution (1.2) at $`x=x_0,t=t_0`$. Let be $`f^{\prime \prime \prime }(u)>0`$ in a deleted neighborhood of $`u=u^{}`$. Then the one-phase Whitham equations has a solution $`u_1>u_2>u_3`$ within a cusp in the $`xt`$ plane for a short time after the breaking time of the zero-phase solution. Furthermore this solution satisfies the boundary conditions (3.11) and (3.14) on the cusp. If the initial data satisfies the condition $`f^{\prime \prime \prime }(u)>0`$ for all $`u`$ except $`u=u^{}`$, then the solution of the Whitham equations exists for all $`t>0`$. The solution is of genus one inside the cusp $`x^{}(t)<x<x^+(t)`$, $`t>t_0`$, where $`x^{}(t)<x^+(t)`$ are two real functions satisfying the condition $`x^{}(t_0)=x^+(t_0)=x_0`$. The solution is of genus zero outside the cusp $`x^{}(t)<x<x^+(t)`$, $`t>t_0`$.
## 4 Solution of Tsarev system
In this section we build the solution of the boundary value problem (3.4), (3.7-3.14) for monotone smooth initial data and we show that the solution obtained is unique. We consider initial data of the form $`x=f(u)|_{t=0}`$ where $`f(u)`$ is a monotone increasing function. The domain of $`f`$ is the interval $`(a,b)`$ where $`\mathrm{}a<b+\mathrm{}`$, and the range of $`f`$ is the real line $`(\mathrm{},+\mathrm{})`$.
In order to obtain the solution of the boundary value problem (3.4), (3.7-3.14) we need the following technical lemma.
###### Lemma 4.1
The differential $`ds^g(r)`$ defined in (3.17) can be written in the form
$$ds^g(r)=2\mu (r)\left(_r\mathrm{\Psi }^g(r;\stackrel{}{u})+\underset{k=1}{\overset{2g+1}{}}_{u_k}\mathrm{\Psi }^g(r;\stackrel{}{u})\right)dr+\frac{R^g(r)}{\mu (r)}dr,$$
(4.1)
where
$$\mathrm{\Psi }^g(r;\stackrel{}{u})=\underset{z=\mathrm{}}{Res}\left[\frac{(z)dz}{2\mu (z)(zr)}\right],q_k(\stackrel{}{u})=\underset{z=\mathrm{}}{Res}\left[\frac{z^{gk}(z)dz}{2\mu (z)}\right],k=1,\mathrm{}g,$$
(4.2)
$$(z)=_0^z\frac{f_a(\xi )}{\sqrt{z\xi }}𝑑\xi ,$$
(4.3)
$$R^g(r)=2\underset{k=1}{\overset{2g+1}{}}_{u_k}q_g(\stackrel{}{u})\underset{l=1,lk}{\overset{2g+1}{}}(ru_l)+\underset{k=1}{\overset{g}{}}q_k(\stackrel{}{u})\underset{l=1}{\overset{k}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{kl}P_{l1}^g(r),$$
(4.4)
and the polynomials $`P_l^g(r)`$, $`l0`$, have been defined in (2.7), the $`\stackrel{~}{\mathrm{\Gamma }}_k`$’s have been defined in (2.19).
Proof: using the second identity in (2.22), we rewrite the differential $`ds^g(r)`$ defined in (3.17) in the form
$$ds^g(r)=d_r\left(\underset{z=\mathrm{}}{Res}[\omega _r^g(z)(z)]\right),$$
(4.5)
where $`\omega _r^g(z)`$ has been defined in (2.14) and $`(z)`$ is the Abel transform defined in (4.3) of the analytic initial data (3.15). The identity (4.5) can be checked straightforward. Using the explicit expression of $`\omega _r^g(z)`$ in (2.20) we obtain
$$ds^g(r)=2d_r(\mu (r)\mathrm{\Psi }^g(r;\stackrel{}{u}))+\underset{k=1}{\overset{g}{}}q_k(\stackrel{}{u})\underset{l=1}{\overset{k}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{kl}\sigma _{l1}^g(r),$$
(4.6)
where $`\mathrm{\Psi }^g(r;\stackrel{}{u})`$ and $`q_k(\stackrel{}{u})`$ have been defined in (4.2).
From (4.2) we get the relations
$$\frac{\mathrm{\Psi }^g(r;\stackrel{}{u})}{ru_i}\frac{\mathrm{\Psi }^g(u_i;\stackrel{}{u})}{ru_i}=2_{u_i}\mathrm{\Psi }^g(r;\stackrel{}{u}),\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}2}_{u_i}q_g(\stackrel{}{u})=\mathrm{\Psi }^g(u_i;\stackrel{}{u})$$
(4.7)
and for $`g=0`$ we define $`u_1=u`$ and
$$2_uq_0(u):=\mathrm{\Psi }^0(u;u)=f_a(u).$$
Using (4.7) we transform the expression for $`ds^g(r)`$ in (4.6) to the form (4.1). $`\mathrm{}`$
The relation (4.1) enables us to write the quantities $`w_i(\stackrel{}{u})=\frac{ds^g(r)}{dp^g(r)}|_{r=u_i},i=1,\mathrm{},2g+1,`$ in (3.16) in the form
$$w_i(\stackrel{}{u})=\frac{1}{P_0^g(u_i)}\left[2_{u_i}q_g(\stackrel{}{u})\underset{l=1,li}{\overset{2g+1}{}}(u_iu_l)+\underset{k=1}{\overset{g}{}}q_k(\stackrel{}{u})\underset{l=1}{\overset{k}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{kl}P_{l1}^g(u_i)\right].$$
(4.8)
Observe that in the formula (4.8) all the information on the initial data is contained in the functions $`q_k(\stackrel{}{u})`$. The functions $`q_k=q_k(\stackrel{}{u})`$, $`k=1,\mathrm{},g`$, solve the linear over-determined system
$`\begin{array}{ccc}\hfill 2(u_iu_j){\displaystyle \frac{^2q_k(\stackrel{}{u})}{u_iu_j}}& =& {\displaystyle \frac{q_k(\stackrel{}{u})}{u_i}}{\displaystyle \frac{q_k(\stackrel{}{u})}{u_j}},i,j=1,\mathrm{},2g+1,\hfill \\ & & \\ \hfill q_k(\underset{2g+1}{\underset{}{u,u,\mathrm{},u}})& =& {\displaystyle \frac{2^{g1}}{(2g1)!!}}u^{k+\frac{1}{2}}{\displaystyle \frac{d^{gk}}{du}}\left(u^{g\frac{1}{2}}f_a^{(k1)}(u)\right),\hfill \end{array}`$ (4.12)
where $`f_a^{(k1)}(u)`$ is the $`(k1)th`$ derivative of the polynomial initial data $`f_a(u)`$. The function $`\mathrm{\Psi }^g(r;\stackrel{}{u})`$ in (4.2) satisfies a similar linear over-determined system.
###### Theorem 4.2 (First Main Theorem)
Let be $`f(u)`$ a smooth monotone increasing function with domain $`(a,b)`$, $`\mathrm{}a<b+\mathrm{}`$. If $`q_k=q_k(u_1,u_2,\mathrm{},u_{2g+1})`$, $`1kg`$, is the symmetric solution of the linear over-determined system
$`\{\begin{array}{ccc}2(u_iu_j){\displaystyle \frac{^2q_k(\stackrel{}{u})}{u_iu_j}}={\displaystyle \frac{q_k(\stackrel{}{u})}{u_i}}{\displaystyle \frac{q_k(\stackrel{}{u})}{u_j}},ij,i,j=1,\mathrm{},2g+1,g>0\hfill & & \\ & & \\ q_k(\underset{2g+1}{\underset{}{u,u,\mathrm{},u}})=F_k(u)\hfill & & \\ & & \\ F_k(u)={\displaystyle \frac{2^{(g1)}}{(2g1)!!}}u^{k+\frac{1}{2}}{\displaystyle \frac{d^{gk}}{du^{gk}}}\left(u^{g\frac{1}{2}}f^{(k1)}(u)\right),\hfill & & \end{array}`$ (4.18)
with the ordering $`u_1>u_2>\mathrm{}>u_{2g+1}`$, then
$$w_i(\stackrel{}{u})=\frac{1}{P_0^g(u_i)}\left[2_{u_i}q_g(\stackrel{}{u})\underset{l=1,li}{\overset{2g+1}{}}(u_iu_l)+\underset{k=1}{\overset{g}{}}q_k(\stackrel{}{u})\underset{l=1}{\overset{k}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{kl}P_{l1}^g(u_i)\right],i=1,\mathrm{},2g+1,$$
(4.19)
solves the boundary value problem (3.4), (3.7-3.14). Conversely every solution of (3.4), (3.7-3.14) can be obtained in this way.
Before proving the theorem we show how to obtain the solution of system (4.18) for generic smooth initial data. We follow the procedure in . We start with the following lemma.
###### Lemma 4.3
The system
$`\{\begin{array}{ccc}2(zy)p_{zy}& =\hfill & p_z\rho p_y,\rho >0\hfill \\ p(z,z)& =\hfill & s(z)\hfill \end{array}`$ (4.22)
has, for any smooth initial data $`s(z)`$ one and only one solution. Moreover, the solution can be written explicitly
$$p(z,y)=\frac{1}{C_\rho }_1^1\frac{s(\frac{1+\mu }{2}z+\frac{1\mu }{2}y)}{\sqrt{1\mu }}(1+\mu )^{\frac{\rho }{2}1}𝑑\mu $$
(4.23)
where
$$C_\rho =_1^1\frac{(1+\mu )^{\frac{\rho }{2}1}}{\sqrt{1\mu }}𝑑\mu .$$
(4.24)
Using the above lemma, the linear over-determined systems (4.18) can be integrated for any smooth initial data in the following way. Suppose that $`q_k(u_1,u_2,\mathrm{},u_{2g+1})`$ is a solution of (4.18).
Clearly $`A_k(u_1,u_{2g+1})=q_k(\underset{2g}{\underset{}{u_1,u_1,\mathrm{},u_1}},u_{2g+1})`$ satisfies
$`\begin{array}{ccc}\hfill 2(u_1u_{2g+1}){\displaystyle \frac{^2A_k}{u_1u_{2g+1}}}& =& {\displaystyle \frac{A_k}{u_1}}2g{\displaystyle \frac{A_k}{u_{2g+1}}}\hfill \\ & & \\ \hfill A_k(u,u)& =& F_k(u)\hfill \end{array}`$ (4.28)
which by lemma 4.3 implies that
$$A_k(u_1,u_{2g+1})=\frac{1}{C_{2g}}_1^1\frac{F_k(\frac{1+\xi _{2g}}{2}u_1+\frac{1\xi _{2g}}{2}u_{2g+1})}{\sqrt{1\xi _{2g}}}(1+\xi _{2g})^{g1}𝑑\xi _{2g}.$$
(4.29)
For each fixed $`u_{2g+1}`$ the function $`B_k(u_1,u_{2g},u_{2g+1})=q_k(\underset{2g1}{\underset{}{u_1,\mathrm{},u_1}},u_{2g},u_{2g+1})`$ satisfies
$`\begin{array}{ccc}2(u_1u_{2g}){\displaystyle \frac{^2B_k}{u_1u_{2g}}}& =\hfill & {\displaystyle \frac{B_k}{u_1}}(2g1){\displaystyle \frac{B_k}{u_{2g}}}\hfill \\ & & \\ B_k(u,u,u_{2g+1})& =\hfill & A_k(u,u_{2g+1})\hfill \end{array}`$ (4.33)
Using again lemma 4.3 we obtain
$$\begin{array}{cc}\hfill B_k(u_1,u_{2g},u_{2g+1})=& \frac{1}{C_{2g}C_{2g1}}_1^1_1^1d\xi _{2g}d\xi _{2g1}(1+\xi _{2g})^{g1}(1+\xi _{2g1})^{g\frac{3}{2}}\times \hfill \\ & \\ & \frac{F_k(\frac{1+\xi _{2g}}{2}(\frac{1+\xi _{2g1}}{2}u_1+\frac{1+\xi _{2g1}}{2}u_{2g})\frac{1\xi _{2g}}{2}u_{2g+1})}{\sqrt{1\xi _{2g}}\sqrt{1\xi _{2g1}}}\hfill \end{array}$$
(4.34)
Going on in the process of integration we obtain the solution $`q_k(\stackrel{}{u})=q_k(u_1,u_2,\mathrm{},u_{2g+1})`$ of the boundary value problem (4.18) namely
$$\begin{array}{cc}\hfill q_k(\stackrel{}{u})=& \frac{1}{C}_1^1_1^1\mathrm{}_1^1d\xi _1d\xi _2\mathrm{}d\xi _{2g}(1+\xi _{2g})^{g1}(1+\xi _{2g1})^{g\frac{3}{2}}\mathrm{}(1+\xi _3)^{\frac{1}{2}}(1+\xi _1)^{\frac{1}{2}}\times \hfill \\ & \\ & \frac{F_k(\frac{1+\xi _{2g}}{2}(\mathrm{}(\frac{1+\xi _2}{2}(\frac{1+\xi _1}{2}u_1+\frac{1\xi _1}{2}u_2)+\frac{1\xi _2}{2}u_3)+\mathrm{})+\frac{1\xi _{2g}}{2}u_{2g+1})}{\sqrt{(1\xi _1)(1\xi _2)\mathrm{}(1\xi _{2g})}}\hfill \end{array}$$
(4.35)
where $`C=_{j=1}^{2g}C_j`$ and $`C_j`$ has been defined in (4.24). When the initial data is of the form (3.15), the expression (4.35) for the $`q_k(\stackrel{}{u})`$’s is equivalent to (4.2).
###### Theorem 4.4
The boundary value problem (4.18) has one and only one solution. The solution is symmetric and is given by (4.35).
Proof: Uniqueness follows from lemma 4.3 and the argument previous to (4.35). The boundary condition (4.18) is clearly satisfied. The symmetry follows from the construction. Indeed in the process of integration we can interchange the role of any of the variable $`u_i`$. $`\mathrm{}`$ We have the following relations.
###### Lemma 4.5
The functions $`F_k(u)`$ and the solutions $`q_k(\stackrel{}{u})`$, $`k=1,\mathrm{},g`$, of the boundary value problem (4.18) satisfy the following relations.
$`\begin{array}{ccc}_uF_k(u)& =\hfill & {\displaystyle \frac{2g+1}{2}}F_{k+1}(u)+u_uF_{k+1}(u),k=1,\mathrm{},g1,g>0\hfill \\ & & \\ _{u_i}q_k(\stackrel{}{u})& =\hfill & {\displaystyle \frac{1}{2}}q_{k+1}(\stackrel{}{u})+u_i_{u_i}q_{k+1}(\stackrel{}{u})i=1.\mathrm{},2g+1,k=1,\mathrm{},g1g>0.\hfill \end{array}`$ (4.39)
Proof of Theorem 4.2 (First Main Theorem).
We consider the non trivial case where $`q_k(\stackrel{}{u})0,k=1,\mathrm{}g,`$ and $`_{u_j}q_g(\stackrel{}{u})0`$, $`j=1,\mathrm{},2g+1`$.
The proof consists of three parts.
a)The $`w_i(\stackrel{}{u})`$’s defined in (4.19) satisfy (3.4).
Using the definition of $`w_i(\stackrel{}{u})`$ in (4.19) we have the following relation
$$\begin{array}{cc}\hfill _{u_j}w_i(\stackrel{}{u})=& 2\frac{_{\stackrel{l=1}{li}}^{2g+1}(u_iu_l)}{P_0^g(u_i)}_{u_j}_{u_i}q_g(\stackrel{}{u})2\frac{_{\stackrel{l=1}{li}}^{2g+1}(u_iu_l)}{P_0^g(u_i)}_{u_i}q_g(\stackrel{}{u})\left(\frac{_{u_j}P_0^g(u_i)}{P_0^g(u_i)}+\frac{1}{u_iu_j}\right)\hfill \\ \hfill +& _{u_j}\left(\underset{l=1}{\overset{g}{}}(2l1)\frac{P_{l1}^g(u_i)}{P_0^g(u_i)}\underset{k=l}{\overset{g}{}}q_k(\stackrel{}{u})\stackrel{~}{\mathrm{\Gamma }}_{kl}\right),ij,i,j=1,\mathrm{},2g+1.\hfill \end{array}$$
(4.40)
The following identities hold
$$\frac{}{_{u_j}}\frac{P_k^g(u_i)}{P_0^g(u_i)}=\frac{1}{\lambda _i\lambda _j}\frac{\lambda _i}{u_j}\left(\frac{P_k^g(u_i)}{P_0^g(u_i)}\frac{P_k^g(u_j)}{P_0^g(u_j)}\right),ij,i,j=1,\mathrm{},2g+1,k1$$
(4.41)
$$\frac{1}{\lambda _i\lambda _j}\frac{\lambda _i}{u_j}=\frac{_{u_j}P_0^g(u_i)}{P_0^g(u_i)}\frac{1}{2}\frac{1}{u_iu_j},ij,i,j=1,\mathrm{},2g+1,$$
(4.42)
where $`\lambda _i(\stackrel{}{u})`$ has been defined in (3.1) and $`P_k^g(r)`$ has been defined in (2.4).
Using the definition of $`\stackrel{~}{\mathrm{\Gamma }}_k`$ in (2.19) it is easy to verify that
$$\frac{\stackrel{~}{\mathrm{\Gamma }}_k}{u_j}=\frac{1}{2}\underset{m=1}{\overset{k1}{}}\stackrel{~}{\mathrm{\Gamma }}_{km}u_j^{m1}.$$
(4.43)
Applying repeatedly the relations (4.39) we obtain the following expression for $`_{u_j}q_k(\stackrel{}{u})`$:
$$_{u_j}q_k(\stackrel{}{u})=\frac{1}{2}\underset{m=1}{\overset{gk}{}}q_{m+k}(\stackrel{}{u})u_j^{m1}+u_j^{gk}_{u_j}q_g(\stackrel{}{u}),k=1,\mathrm{},g1.$$
(4.44)
From (4.18) and (4.41-4.44) we can write $`_{u_j}w_i(\stackrel{}{u})`$ in (4.40) in the form
$$\begin{array}{cc}\hfill _{u_j}w_i(\stackrel{}{u})=& \frac{_{l=1,li}^{2g+1}(u_iu_l)}{P_0^g(u_i)}\frac{_{u_i}q_g(\stackrel{}{u})_{u_j}q_g(\stackrel{}{u})}{u_iu_j}\hfill \\ \hfill & 2\frac{_{l=1,li}^{2g+1}(u_iu_l)}{P_0^g(u_i)}_{u_i}q_g(\stackrel{}{u})\left(\frac{_{u_j}P_0^g(u_i)}{P_0^g(u_i)}+\frac{1}{u_iu_j}\right)\hfill \\ \hfill +& \underset{l=1}{\overset{g}{}}(2l1)\frac{P_{l1}^g(u_i)}{P_0^g(u_i)}\left(\frac{1}{2}\underset{k=l}{\overset{g1}{}}\stackrel{~}{\mathrm{\Gamma }}_{kl}\underset{m=1}{\overset{gk}{}}q_{m+k}(\stackrel{}{u})u_j^{m1}+\underset{k=l}{\overset{g}{}}\stackrel{~}{\mathrm{\Gamma }}_{kl}u_j^{gk}_{u_j}q_g(\stackrel{}{u})\right)\hfill \\ \hfill +& \underset{l=1}{\overset{g}{}}(2l1)\frac{P_{l1}^g(u_i)}{P_0^g(u_i)}\underset{k=l}{\overset{g}{}}q_k(\stackrel{}{u})\left(\frac{1}{2}\underset{m=1}{\overset{kl}{}}\stackrel{~}{\mathrm{\Gamma }}_{klm}u_j^{m1}\right)\hfill \\ \hfill +& \underset{l=1}{\overset{g}{}}(2l1)\frac{1}{\lambda _i\lambda _j}\frac{\lambda _i}{u_j}\left(\frac{P_{l1}^g(u_i)}{P_0^g(u_i)}\frac{P_{l1}^g(u_j)}{P_0^g(u_j)}\right)\underset{k=l}{\overset{g}{}}q_k(\stackrel{}{u})\stackrel{~}{\mathrm{\Gamma }}_{kl}ij,i,j=1,\mathrm{},2g+1.\hfill \end{array}$$
Simplifying we obtain
$$\begin{array}{cc}\hfill _{u_j}w_i(\stackrel{}{u})=& \frac{_{l=1,li}^{2g+1}(u_iu_l)}{(u_iu_j)P_0^g(u_i)}_{u_j}q_g(\stackrel{}{u})+2\frac{_{l=1,li}^{2g+1}(u_iu_l)}{P_0^g(u_i)}_{u_i}q_g(\stackrel{}{u})\left(\frac{1}{\lambda _i\lambda _j}\frac{\lambda _i}{u_j}\right)\hfill \\ \hfill +& \underset{l=1}{\overset{g}{}}(2l1)\frac{P_{l1}^g(u_i)}{P_0^g(u_i)}\underset{k=l}{\overset{g}{}}\stackrel{~}{\mathrm{\Gamma }}_{kl}u_j^{gk}_{u_j}q_g(\stackrel{}{u})\hfill \\ \hfill +& \frac{1}{\lambda _i\lambda _j}\frac{\lambda _i}{u_j}\underset{l=1}{\overset{g}{}}(2l1)\left(\frac{P_{l1}^g(u_i)}{P_0^g(u_i)}\frac{P_{l1}^g(u_j)}{P_0^g(u_j)}\right)\underset{k=l}{\overset{g}{}}q_k(\stackrel{}{u})\stackrel{~}{\mathrm{\Gamma }}_{kl}ij,i,j=1,\mathrm{},2g+1.\hfill \end{array}$$
(4.45)
Adding and subtracting the quantity $`{\displaystyle \frac{1}{\lambda _i\lambda _j}}{\displaystyle \frac{\lambda _i}{u_j}}w_j`$ to (4.45), we can reduce it to the form
$$\begin{array}{cc}\hfill _{u_j}w_i& \frac{1}{\lambda _i\lambda _j}\frac{\lambda _i}{u_j}[w_iw_j]=_{u_j}q_g(\stackrel{}{u})(\frac{2}{\lambda _i\lambda _j}\frac{\lambda _i}{u_j}\frac{_{l=1,lj}^{2g+1}(u_ju_l)}{P_0^g(u_j)}\hfill \\ \hfill & \frac{_{l=1,li}^{2g+1}(u_iu_l)}{(u_iu_j)P_0^g(u_i)}+\underset{l=1}{\overset{g}{}}(2l1)\frac{P_{l1}^g(u_i)}{P_0^g(u_i)}\underset{k=l}{\overset{g}{}}\stackrel{~}{\mathrm{\Gamma }}_{kl}u_j^{gk})\hfill \end{array}$$
(4.46)
The term in parenthesis in the right hand side of (4.46) does not depend on the initial data $`f(u)`$. It is identically zero for the analytic initial data (3.15) because in such case the $`w_i`$’s satisfy (3.4) . Therefore we can conclude that
$$_{u_j}w_i\frac{1}{\lambda _i\lambda _j}\frac{\lambda _i}{u_j}[w_iw_j]=0$$
(4.47)
for any smooth monotone increasing initial data $`x=f(u)`$. b)The $`w_i(\stackrel{}{u})`$’s satisfy the boundary conditions (3.7-3.14).
In the following we use the sup-script $`g`$ to denote the corresponding genus of the quantities we are referring to. We have the following relations. When $`u_l=u_{l+1}=v`$ for $`l=1,\mathrm{},2g`$, the $`\stackrel{~}{\mathrm{\Gamma }}_k`$’s defined in (2.19) satisfy
$$\begin{array}{cc}\hfill \stackrel{~}{\mathrm{\Gamma }}_k^g(u_1,\mathrm{},u_{l1},v,v,u_{l+2},\mathrm{},u_{2g+1})=& \stackrel{~}{\mathrm{\Gamma }}_k^{g1}(u_1,\mathrm{},u_{l1},u_{l+2},\mathrm{},u_{2g+1})\hfill \\ \hfill & v\stackrel{~}{\mathrm{\Gamma }}_{k1}^{g1}(u_1,\mathrm{},u_{l1},u_{l+2},\mathrm{},u_{2g+1}),k1,g>1\hfill \end{array}$$
(4.48)
and the $`q_k(\stackrel{}{u})`$’s defined in (4.35) satisfy
$$\begin{array}{cc}\hfill q_k^g(u_1,\mathrm{},u_{l1},v,v,u_{l+2},\mathrm{},u_{2g+1})& vq_{k+1}^g(u_1,\mathrm{},u_{l1},v,v,u_{l+2},\mathrm{},u_{2g+1})\hfill \\ & =q_k^{g1}(u_1,\mathrm{},u_{l1},u_{l+2},\mathrm{},u_{2g+1}),k=1,\mathrm{},g1,g>1.\hfill \end{array}$$
(4.49)
For $`u_iu_l=u_{l+1}=v`$ we have that
$$\begin{array}{cc}\hfill _{u_i}q_{g1}^{g1}(u_1,\mathrm{},u_{l1},u_{l+2},\mathrm{},u_{2g+1})& \frac{1}{2}q_g^g(u_1,\mathrm{},u_{l1},v,v,u_{l+2},\mathrm{},u_{2g+1})=\hfill \\ & (u_iv)_{u_i}q_g^g(u_1,\mathrm{},u_{l1},v,v,u_{l+2},\mathrm{},u_{2g+1}),\hfill \end{array}$$
(4.50)
which follows from (4.39). When $`u_l=u_{l+1}=v`$ the polynomials $`P_k^g(r)`$’s defined in (2.7) satisfy the relation
$$P_k^g(r)=(rv)P_k^{g1}(r),k0.$$
(4.51)
Using the relations (4.48-4.51) we have that for $`il,l+1`$, $`i=1,\mathrm{},2g+1`$,
$$\begin{array}{cc}\hfill w_i^g(u_1,\mathrm{},u_{l1},v,v,u_{l+2},& \mathrm{},u_{2g+1})=2\frac{_{k=1,ki,l,l+1}^{2g+1}(u_iu_k)}{P_0^{g1}(u_i)}_{u_i}q^{g1}_{g1}\hfill \\ \hfill +& \underset{m=1}{\overset{g1}{}}(2m1)\frac{P_{m1}^{g1}(u_i)}{P_0^{g1}(u_i)}\underset{k=m}{\overset{g1}{}}q_k^{g1}\stackrel{~}{\mathrm{\Gamma }}_{km}^{g1}\hfill \\ \hfill +& \left(\underset{m=1}{\overset{g}{}}(2m1)\frac{P_{m1}^{g1}(u_i)}{P_0^{g1}(u_i)}\stackrel{~}{\mathrm{\Gamma }}_{gm}^{g1}\frac{_{k=1,ki,l,l+1}^{2g+1}(u_iu_k)}{P_0^{g1}(u_i)}\right)q_g^g,\hfill \end{array}$$
(4.52)
which reduces to the form
$$\begin{array}{cc}\hfill w_i^g(u_1,\mathrm{},u_{l1},v,v,u_{l+2},& \mathrm{},u_{2g+1})=w_i^{g1}(u_1,\mathrm{},u_{l1},u_{l+2},\mathrm{},u_{2g+1})\hfill \\ \hfill +& \left(\underset{m=1}{\overset{g}{}}(2m1)\frac{P_{m1}^{g1}(u_i)}{P_0^{g1}(u_i)}\stackrel{~}{\mathrm{\Gamma }}_{gm}^{g1}\frac{_{k=1,ki,l,l+1}^{2g+1}(u_iu_k)}{P_0^{g1}(u_i)}\right)q_g^g.\hfill \end{array}$$
(4.53)
Using (2.21) the term in parenthesis in the right hand side of (4.53) turns out to be identically zero. Therefore the boundary conditions (3.8) are satisfied for any smooth monotone increasing initial data.
Since the functions $`q_k(\stackrel{}{u})`$’s in (4.35) and the $`\stackrel{~}{\mathrm{\Gamma }}_k(\stackrel{}{u})`$’s defined in (2.19) are symmetric with respect to $`u_1,u_2,\mathrm{},u_{2g+1}`$ we immediately deduce from (4.51) that
$$w_l^g(u_1,\mathrm{},u_{l1},v,v,u_{l+2},\mathrm{},u_{2g+1})=w_{l+1}^g(u_1,\mathrm{},u_{l1},v,v,u_{l+2},\mathrm{},u_{2g+1})$$
(4.54)
so that the boundary conditions (3.7) are satisfied (for a more detailed analysis of this limit see section 5.1). When $`g=1`$ we deduce from (4.54)
$$w_1(u_1,u_1,u_3)=w_2(u_1,u_1,u_3)$$
and from (4.50-4.51)
$$w_3(u_1,u_1,u_3)=2(u_3u_1)_{u_3}q_1(u_1,u_1,u_3)+q_1(u_1,u_1,u_3).$$
From (4.18) and (4.35) we get the relation
$$q_1(u_1,u_1,u_3)=f(u_3)+2(u_1u_3)_{u_3}q_1(u_1,u_1,u_3)$$
so that
$$w_3(u_1,u_1,u_3)=f(u_3).$$
An analogous result can be obtain when $`u_2=u_3`$, so that the boundary conditions (3.11-3.14) are satisfied. c) Uniqueness. We prove that when $`f(u)0`$ then $`w_i^g(\stackrel{}{u})0`$ for all $`b>u_1>u_2>\mathrm{}>u_{2g+1}>a`$, for $`1i2g+1`$ and for any $`g0`$. The proof is by induction on $`g`$. For $`g=0`$ the statement is satisfied.
For $`g=1`$ we repeat the arguments of . We fix $`u_2`$ and we consider the equation (3.4) with boundary condition (3.11-3.14), namely
$`\begin{array}{ccc}\hfill {\displaystyle \frac{w_1}{u_3}}& =& a_{13}[w_1w_3]\hfill \\ & & \\ \hfill {\displaystyle \frac{w_3}{u_1}}& =& a_{31}[w_3w_1]\hfill \\ \hfill w_1(u_1,u_2,u_2)& =& f(u_1)0\hfill \\ \hfill w_3(u_2,u_2,u_3)& =& f(u_3)0.\hfill \end{array}`$
We can regard the above equations as a first-order linear ordinary differential equations with non homogeneous term. Integrating them we obtain a couple integral equation. By standard contraction mapping method it can be shown that when $`f(u)0`$ this system has only zero solution, i.e. $`w_1=w_30`$ for $`(u_1,u_3)`$ satisfying $`b>u_1>u_2>u_3>a`$. Because of the arbitrariness of $`u_2`$, $`w_1`$ and $`w_3`$ vanish as a function of $`(u_1,u_2,u_3)`$ and therefore, by (3.4) so does $`w_2(\stackrel{}{u})`$. Now we suppose the theorem true for genus $`g1`$ and we proof it for genus $`g`$. We fix $`b>u_2>u_3>\mathrm{}>u_{2g}>a`$ and we consider the equation (3.4) for $`w_1^g`$ and $`w_{2g+1}^g`$ with boundary condition (3.7-3.8), namely
$`\begin{array}{ccc}\hfill {\displaystyle \frac{}{u_{2g+1}}}w_1^g& =\hfill & a_{1(2g+1)}[w_1^gw_{2g+1}^g]\hfill \\ \hfill {\displaystyle \frac{}{u_1}}w_{2g+1}^g& =\hfill & a_{(2g+1)1}[w_{2g+1}^gw_1^g]\hfill \\ \hfill w_1^g(u_1,u_2,\mathrm{},u_{2g},u_{2g})& =\hfill & w_1^{g1}(u_1,u_2,\mathrm{},u_{2g1},\widehat{u}_{2g},\widehat{u}_{2g})0\hfill \\ \hfill w_{2g+1}^g(u_2,u_2,\mathrm{},u_{2g},u_{2g+1})& =\hfill & w_{2g+1}^{g1}(\widehat{u}_2,\widehat{u}_2,u_3,\mathrm{},u_{2g},u_{2g+1})0.\hfill \end{array}`$ (4.60)
Repeating the arguments developed for genus $`g=1`$ we may conclude that $`w_1^g(\stackrel{}{u})=w_{2g+1}^g(\stackrel{}{u})0`$, for arbitrary $`b>u_1>u_2>\mathrm{}>u_{2g+1}>a`$. We then repeat the above argument fixing $`b>u_1>u_3>\mathrm{}>u_{2g1}>u_{2g+1}>a`$ and considering the equations (3.4) for $`w_2^g(\stackrel{}{u})`$ and $`w_{2g}^g(\stackrel{}{u})`$, namely
$`\begin{array}{ccc}\hfill {\displaystyle \frac{}{u_{2g}}}w_2^g& =\hfill & a_{2(2g)}[w_2^gw_{2g}^g]\hfill \\ \hfill {\displaystyle \frac{}{u_2}}w_{2g}^g& =\hfill & a_{(2g)2}[w_{2g}^gw_2^g]\hfill \\ \hfill w_2^g(u_1,u_2,\mathrm{},u_{2g1},u_{2g+1},u_{2g+1})& =\hfill & w_2^{g1}(u_1,u_2,\mathrm{},u_{2g1},\widehat{u}_{2g+1},\widehat{u}_{2g+1})0\hfill \\ \hfill w_{2g}^g(u_1,u_1,u_3,\mathrm{},u_{2g},u_{2g+1})& =\hfill & w_{2g}^{g1}(\widehat{u}_1,\widehat{u}_1,u_3,\mathrm{},u_{2g},u_{2g+1})0.\hfill \end{array}`$ (4.65)
It can be easily shown that also $`w_2^g(\stackrel{}{u})=w_{2g}^g(\stackrel{}{u})0`$ for arbitrary $`b>u_1>u_2>\mathrm{}>u_{2g+1}>a`$. Repeating these arguments other $`g2`$ times we conclude that $`w_i^g(\stackrel{}{u})0`$ for $`1ig,g+2i2g+1`$ and for arbitrary $`b>u_1>u_2>\mathrm{}>u_{2g+1}>a`$. Applying (3.4) and the boundary conditions (3.7)-(3.8) we can prove that also $`w_{g+1}^g(\stackrel{}{u})`$ is identically zero. The theorem is then proved. $`\mathrm{}`$ In the next section we consider the problem of reality of the solution of the hodograph transform (3.6).
###### Remark 4.6
The boundary conditions (3.7-3.14) guarantee the $`C^1`$-smoothness of the solution of the Whitham equations. Indeed it can be proved that the $`x`$ derivatives $`_xu_i(x,t)`$, $`i=1,\mathrm{}2g+1`$, are continuous on the phase transition boundaries.
## 5 An upper bound to the genus of the solution
In this section we give an upper bound to the genus of the solution of the Whitham equations for the monotone increasing smooth initial data $`x=f(u)|_{t=0}`$.
###### Theorem 5.1 (Second Main Theorem)
If the monotone increasing smooth initial data $`f(u)`$ satisfies the condition
$$\frac{d^{2N+1}}{du^{2N+1}}f(u)>0,1N\mathrm{I}\mathrm{N},$$
(5.1)
for all $`u(a,b)`$ except at one point, then the solution of the Whitham equations (1.1) has genus at most $`N`$.
Proof: the solution of the Whitham equations (1.1) for different $`g0`$ determines a decomposition of the $`xt`$ plane, $`t0`$, into a number of domains $`D_g`$ with g=0,1,2…, (see Fig. 5.3).
To the inner part of each domain $`D_g`$ it corresponds the $`g`$-phase solution $`b>u_1(x,t)>u_2(x,t)>\mathrm{}>u_{2g+1}(x,t)>a`$ of the Whitham equations (1.1). The common boundaries of the domains $`D_g`$, $`g0`$, and $`D_{g+n}`$, $`n1`$, are the points of phase transition between the $`g`$-phase solution and the $`(g+n)`$-phase solution.
We will show that each domain $`D_g`$, $`gN`$, does not have a common boundary with any of the domains $`D_m`$, $`m>N`$. Since the set of domains $`\{D_g\}_{gN}`$ is not empty because $`D_0\mathrm{}`$ the set $`\{D_g\}_{g>N}`$ must be empty. Indeed on the contrary the $`xt`$ plane, $`t0`$, which is a connected set, would be split into a number of domains whose union forms a disconnected set.
Before determining the boundaries of the domains $`D_g`$, $`g0`$, we need to study more in detail the hodograph transform (3.6).
###### Proposition 5.2
Let us consider the polynomial
$$Z^g(r):=xP_0^g(r)12tP_1^g(r)+R^g(r),$$
(5.2)
where $`R^0(r)=f(u)`$ and $`R^g(r)`$, $`g>0`$, is given by the expression
$$R^g(r)=2\underset{k=1}{\overset{2g+1}{}}_{u_k}q_g(\stackrel{}{u})\underset{l=1,lk}{\overset{2g+1}{}}(ru_l)+\underset{k=1}{\overset{g}{}}q_k(\stackrel{}{u})\underset{l=1}{\overset{k}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{kl}P_{l1}^g(r),$$
(5.3)
with the polynomials $`P_l^g(r)`$, $`l0`$, defined in (2.7) and the functions $`q_k(\stackrel{}{u})`$, $`k=1,\mathrm{},g`$, defined in (4.18). Then the hodograph transform (3.6) is equivalent, for $`g>0`$, to the equation
$$Z^g(r)0,g>0.$$
(5.4)
Proof: We observe that the $`w_i(\stackrel{}{u})`$’s defined in (4.19) are given by the ratio $`w_i(\stackrel{}{u})={\displaystyle \frac{R^g(u_i)}{P_0^g(u_i)}}`$, $`i=1,\mathrm{},2g+1`$, where $`R^g(r)`$ is the polynomial defined in (5.3). Hence we can write the hodograph transform (3.6) in the form
$$[xP_0^g(r)12tP_1^g(r)+R^g(r)]_{r=u_i}=0,i=1,\mathrm{},2g+1.$$
(5.5)
For $`g>0`$, $`Z^g(r)`$ is a polynomial of degree $`2g`$ and because of (5.5) it must have at least $`2g+1`$ real zeros. Therefore it is identically zero. Hence for $`g>0,`$ (5.4) is equivalent to (5.5) and (3.6). $`\mathrm{}`$ In the following analysis we give some conditions for the existence of a real solution $`b>u_1(x,t)>u_2(x,t)>\mathrm{},u_{2g+1}(x,t)>a`$ of the hodograph transform (3.6). For the purpose let us consider the function $`\mathrm{\Psi }^g(r;\stackrel{}{u})`$ which solves the boundary value problem
$`\{\begin{array}{ccc}& & {\displaystyle \frac{}{u_i}}\mathrm{\Psi }^g(r;\stackrel{}{u}){\displaystyle \frac{}{u_j}}\mathrm{\Psi }^g(r;\stackrel{}{u})=2(u_iu_j){\displaystyle \frac{^2}{u_iu_j}}\mathrm{\Psi }^g(r;\stackrel{}{u}),ij,i,j=1,\mathrm{}2g+1\hfill \\ & & {\displaystyle \frac{}{r}}\mathrm{\Psi }^g(r;\stackrel{}{u})2{\displaystyle \frac{}{u_j}}\mathrm{\Psi }^g(r;\stackrel{}{u})=2(ru_j){\displaystyle \frac{^2}{ru_j}}\mathrm{\Psi }^g(r;\stackrel{}{u}),j=1,\mathrm{}2g+1\hfill \\ & & \mathrm{\Psi }^g(r;\underset{2g+1}{\underset{}{r,\mathrm{},r}})={\displaystyle \frac{2^g}{(2g+1)!!}}f^{(g)}(r)\hfill \end{array}`$ (5.9)
where $`f^{(g)}(r)`$ is the $`g`$th derivative of the smooth monotone increasing initial data $`f(u)`$. From the results of Sec. 4 we are able to integrate (5.9) obtaining
$$\begin{array}{cc}\hfill \mathrm{\Psi }^g(r;\stackrel{}{u})=& \frac{1}{K}_1^1_1^1\mathrm{}_1^1d\xi _1d\xi _2\mathrm{}d\xi _{2g+1}(1+\xi _{2g+1})^g(1+\xi _{2g})^{g\frac{1}{2}}\mathrm{}(1+\xi _2)^{\frac{1}{2}}\times \hfill \\ & \\ & \frac{f^{(g)}(\frac{1+\xi _{2g+1}}{2}(\mathrm{}(\frac{1+\xi _2}{2}(\frac{1+\xi _1}{2}r+\frac{1\xi _1}{2}u_1)+\frac{1\xi _2}{2}u_2)+\mathrm{})+\frac{1\xi _{2g+1}}{2}u_{2g+1})}{\sqrt{(1\xi _1)(1\xi _2)\mathrm{}(1\xi _{2g+1})}}\hfill \end{array}$$
(5.10)
where $`K={\displaystyle \frac{2^g}{(2g+1)!!}}_{j=2}^{2g+2}C_j`$ and the $`C_j`$’s have been defined in (4.24). The function $`\mathrm{\Psi }^g(r;\stackrel{}{u})`$ is symmetric with respect to the variables $`b>u_1>u_2>\mathrm{}>u_{2g+1}>a`$. For the initial data (3.15), the expression of $`\mathrm{\Psi }^g(r;\stackrel{}{u})`$ in (5.10) is equivalent to (4.2).
###### Proposition 5.3
Let us consider the function
$$\mathrm{\Phi }^g(r;\stackrel{}{u})=_r\mathrm{\Psi }^g(r;\stackrel{}{u})+\underset{i=1}{\overset{2g+1}{}}_{u_i}\mathrm{\Psi }^g(r;\stackrel{}{u}),$$
(5.11)
where $`\mathrm{\Psi }^g(r;\stackrel{}{u})`$ has been defined in (5.10). If $`u_1(x,t)>u_2(x,t)>\mathrm{}>u_{2g+1}(x,t)`$ satisfy the Whitham equations (1.1) then the function $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ has, in the $`r`$ variable, at least one real zero in each of the intervals $`(u_{2k},u_{2k1}),k=1,\mathrm{},g`$, $`g>0`$.
Proof: If the functions $`u_1(x,t)>u_2(x,t)>\mathrm{}>u_{2g+1}(x,t)`$ satisfy the Whitham equations then by proposition 5.2 the polynomial $`Z^g(r)0`$. Therefore
$$\begin{array}{cc}\hfill 0& _{\alpha _k}\frac{Z^g(r)}{\mu (r)}𝑑r=_{\alpha _k}\frac{xP_0^g(r)12tP_1^g(r)+R^g(r)}{\mu (r)}𝑑r\hfill \\ \hfill =& _{\alpha _k}2\frac{_{i=1}^{2g+1}_{u_i}q_g(\stackrel{}{u})_{j=1,ji}^{2g+1}(ru_j)}{\mu (r)}𝑑r,k=1,\mathrm{},g.\hfill \end{array}$$
(5.12)
In the third equality of (5.12) we have used the fact that
$$_{\alpha _k}\frac{P_l^g(r)}{\mu (r)}𝑑r=_{\alpha _k}\sigma _l^g(r)=0,l0,k=1,\mathrm{},g,$$
because of the normalization conditions (2.5). The function $`\mathrm{\Psi }^g(r;\stackrel{}{u})`$ satisfies the relations
$$\frac{\mathrm{\Psi }^g(r;\stackrel{}{u})}{ru_i}\frac{\mathrm{\Psi }^g(u_i;\stackrel{}{u})}{ru_i}=2_{u_i}\mathrm{\Psi }^g(r;\stackrel{}{u}),\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}2}_{u_i}q_g(\stackrel{}{u})=\mathrm{\Psi }^g(u_i;\stackrel{}{u})$$
(5.13)
which can be easily obtained from (5.10). Using (5.13) we can rewrite the last term in (5.12) in the form
$$\begin{array}{cc}\hfill 0=& _{\alpha _k}2\frac{_{i=1}^{2g+1}_{u_i}q_g(\stackrel{}{u})_{j=1,ji}^{2g+1}(ru_j)}{\mu (r)}𝑑r,k=1,\mathrm{},g,\hfill \\ \hfill =& 2_{u_{2k}}^{u_{2k1}}\mu (r)\left(\underset{i=1}{\overset{g}{}}\frac{\mathrm{\Psi }^g(r;\stackrel{}{u})}{ru_i}2_{u_i}\mathrm{\Psi }^g(r;\stackrel{}{u})\right)𝑑r\hfill \\ \hfill =& 4_{u_{2k}}^{u_{2k1}}\mu (r)\left(_r\mathrm{\Psi }^g(r;\stackrel{}{u})+\underset{i=1}{\overset{g}{}}_{u_i}\mathrm{\Psi }^g(r;\stackrel{}{u})\right)𝑑r,k=1,\mathrm{},g,\hfill \end{array}$$
where the last equality has been obtained integrating by parts. Using the definition of $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ in (5.11) we rewrite the above relation in the form
$$0=4_{u_{2k}}^{u_{2k1}}\mu (r)\mathrm{\Phi }^g(r;\stackrel{}{u})𝑑r,k=1,\mathrm{},g.$$
(5.14)
Relation (5.14) is satisfied only if the function $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ changes sign at least once in each of the intervals $`(u_{2k},u_{2k1})`$, $`k=1,\mathrm{},g`$. $`\mathrm{}`$
In the following we are going to determine the equations which describe, on the $`xt0`$ plane, the boundary between the domains $`D_g`$ and $`D_{g+1}`$.
The boundary between the domains $`D_g`$ and $`D_{g+1}`$ represents a singular behavior in the solution of the $`g`$ or $`(g+1)`$-phase equations. As we have shown in the example on fig 1.1, the boundary between the domains $`D_0`$ and $`D_1`$ is described by the curves $`x^\pm (t)`$ where $`u_2(x,t)=u_3(x,t)`$ and $`u_1(x,t)=u_2(x,t)`$ respectively and by the point $`x_0`$, $`t_0`$ of gradient catastrophe of the zero phase solution.
In a similar way the generic boundary between the domains $`D_g`$ and $`D_{g+1}`$ is described by
a) the curves $`x_g^\pm (t)`$ where two Riemann invariants of the $`(g+1)`$-phase solution coalesce;
b) the point of gradient catastrophe $`x_c`$, $`t_c>0`$ of the $`g`$-phase solution, namely the point where one of the $`2g+1`$ Riemann invariants has a vertical inflection point. The equations determining the point of gradient catastrophe of the $`g`$-phase solution can also be obtained considering the limit of the $`g+1`$ phase solution when three Riemann invariants coalesce.
To treat case a) we consider the Riemann surface $`𝒮_{g+1}`$ of genus $`g+1`$ given by the equations
$$\stackrel{~}{\mu }^2=(rv\sqrt{\delta })(rv+\sqrt{\delta })\mu ^2,v\mathrm{I}\mathrm{R},$$
$$\mu ^2=\underset{j=1}{\overset{2g+1}{}}(ru_j),b>u_1>u_2>\mathrm{}>u_{2g+1}>a,$$
where $`vu_j,j=1,\mathrm{}2g+1`$, and $`0<\delta 1`$. The Riemann invariants are the $`2g+3`$ variables $`\stackrel{~}{u}_1=v+\sqrt{\delta },\stackrel{~}{u}_2=v\sqrt{\delta }`$, $`u_1>u_2>\mathrm{}>u_{2g+1}`$. We suppose $`\stackrel{~}{u}_1,\stackrel{~}{u}_2u_j,j=1,\mathrm{},2g+1`$. The hodograph transform (3.6) for these $`2g+3`$ variables has two different behavior for $`\delta 0`$ when $`v`$ belongs to one of the bands
$$(u_{2g+1},u_{2g})(u_{2g1},u_{2g2})\mathrm{}(u_3,u_2)(u_1,b)$$
(5.15)
or gaps
$$(a,u_{2g+1})(u_{2g},u_{2g1})\mathrm{}(u_4,u_3)(u_2,u_1).$$
(5.16)
We call leading edge of the phase transition boundary the case in which $`v`$ belongs to the bands. We call trailing edge of the phase transition boundary the case in which $`v`$ belongs to the gaps (5.16).
###### Theorem 5.4
The leading edge of the phase transition boundary between the $`g`$-phase solution and the $`(g+1)`$-phase solution is described by the system
$`\{\begin{array}{ccc}\mathrm{\Phi }^g(v;\stackrel{}{u})6tϵ_{g0}=0\hfill & & \\ _v\mathrm{\Phi }^g(v;\stackrel{}{u})=0\hfill & & \\ x=\left[12t{\displaystyle \frac{P_1^g(r)}{P_0^g(r)}}+{\displaystyle \frac{R^g(r)}{P_0^g(r)}}\right]_{r=u_i}\hfill & & ,i=1,\mathrm{},2g+1,g0\hfill \end{array}`$ (5.20)
where $`v(u_{2j+1},u_{2j}),\mathrm{\hspace{0.33em}0}jg,u_0=+b`$, the function $`ϵ_{g0}=1`$ for $`g=0`$ and zero otherwise, the function $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ has been defined in (5.11) and the polynomial $`R^g(r)`$ has been defined in (5.3). We assume $`(_v)^2\mathrm{\Phi }^g(v;\stackrel{}{u})0`$ and $`\mathrm{\Phi }^g(u_i;\stackrel{}{u})0`$, $`i=1,\mathrm{},2g+1`$, on the solution of (5.20).
###### Remark 5.5
System (5.20) is a system of $`2g+3`$ equations in $`2g+4`$ unknowns $`x,t,v`$ and $`u_1>u_2>\mathrm{}>u_{2g+1}`$. If system (5.20) is uniquely solvable for real $`x,v`$ and $`u_1>u_2>\mathrm{}>u_{2g+1}`$ as a function of $`t0`$, then a phase transition between the $`g`$-phase solution and the $`(g+1)`$-phase solution occurs. The curve $`x_g^{}=x_g^{}(t)`$ describes on the $`xt0`$ plane the boundary between the domains $`D_g`$ and $`D_{g+1}`$ associated to the leading edge. The conditions $`(_v)^2\mathrm{\Phi }^g(v;\stackrel{}{u})0`$ and $`\mathrm{\Phi }^g(u_i;\stackrel{}{u})0`$, $`i=1,\mathrm{},2g+1`$, on the solution of (5.20) exclude higher order degeneracy in the transition. Indeed in such case it can be proved that in a neighborhood of the solution $`x(t)`$, $`v(t)`$, $`u_1(t)>u_2(t)>\mathrm{}>u_{2g+1}(t)`$ of (5.20) the $`(g+1)`$-phase solution is uniquely defined.
###### Theorem 5.6
The trailing edge of the phase transition boundary between the $`g`$-phase solution and the $`(g+1)`$-phase solution is described by the solution of the system
$`\{\begin{array}{ccc}\mathrm{\Phi }^g(v;\stackrel{}{u})6tϵ_{g0}=0\hfill & & \\ & & \\ {\displaystyle _v^{u_{2j1}}}(\mathrm{\Phi }^g(r;\stackrel{}{u})6tϵ_{g0})\mu (r)𝑑r=0\hfill & & \\ & & \\ x=\left[t{\displaystyle \frac{P_1^g(r)}{P_0^g(r)}}+{\displaystyle \frac{R^g(r)}{P_0^g(r)}}\right]_{r=u_i},i=1,\mathrm{},2g+1,g>0\hfill & & \end{array}`$ (5.26)
where $`v(u_{2j},u_{2j1}),\mathrm{\hspace{0.33em}1}jg+1,u_{2g+2}=a`$ and the function $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ has been defined in (5.11). We assume $`_v\mathrm{\Phi }^g(v;\stackrel{}{u})0`$ and $`\mathrm{\Phi }^g(u_i;\stackrel{}{u})0`$, $`i=1,\mathrm{},2g+1`$, on the solution of (5.26).
If system (5.26) is uniquely solvable for real $`x,v`$ and $`u_1>u_2>\mathrm{}>u_{2g+1}`$ as a function of $`t0`$, then a phase transition between the $`g`$-phase solution and the $`(g+1)`$-phase solution occurs. The curve $`x_g^+=x_g^+(t)`$ describes on the $`xt0`$ plane the boundary between the domains $`D_g`$ and $`D_{g+1}`$ associated to the trailing edge.
The following theorem enables one to determine a point of gradient catastrophe of the $`g`$-phase solution. This point can be obtained either as a limit of the $`(g+1)`$ phase solution when three Riemann invariants coalesce or imposing a vertical inflection point on the $`g`$-phase solution.
###### Theorem 5.7
Let us consider the system
$`\{\begin{array}{ccc}_r\mathrm{\Phi }^g(r;\stackrel{}{u})|_{r=u_l}=0\hfill & & \\ \mathrm{\Phi }^g(u_l;\stackrel{}{u})6tϵ_{g0}=0\hfill & & \\ x=\left[12t{\displaystyle \frac{P_1^g(r)}{P_0^g(r)}}+{\displaystyle \frac{R^g(r)}{P_0^g(r)}}\right]_{r=u_i},i=1,\mathrm{},2g+1,g0\hfill & & \end{array}`$ (5.30)
where $`ϵ_{g0}=1`$ for $`g=0`$ and zero otherwise, the function $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ has been defined in (5.11). When system (5.30) is uniquely solvable for $`x_c`$, $`t_c0`$ and $`b>u_1(x_c,t_c)>u_2(x_c,t_c)>\mathrm{}>u_{2g+1}(x_c,t_c)>a`$, then the $`g`$-phase solution has a point of gradient catastrophe on the $`u_l`$ branch, $`1l2g+1`$.
To avoid higher order degeneracies in the transition we impose the condition
$$(_r)^2\mathrm{\Phi }^g(r;\stackrel{}{u}(x_c,t_c))|_{r=u_l(x_c,t_c)}0.$$
(5.31)
Indeed we can prove that the above condition guarantees that the genus of solution of the Whitham equations increases at most by one in the neighborhood of the point of gradient catastrophe. Therefore it is legitimate to consider the point of gradient catastrophe that solve (5.30) and satisfies (5.31) as a point of the boundary between the domains $`D_g`$ and $`D_{g+1}`$. The condition (5.31) in not essential in the genus $`g=0`$ case as illustrated in Theorem 3.2.
###### Remark 5.8
We observe that both systems (5.20) and (5.26) in the limit $`vu_l`$, $`l=1,\mathrm{},2g+1`$, coincide with system (5.30).
Theorems 5.4, 5.6 and 5.7 characterize all types of boundaries between the $`g`$-phase solution and the $`(g+1)`$-phase solution, namely the leading edge, the trailing edge and the points of gradient catastrophe. We will prove these theorems in the next section.
Example 3.1 For $`x=u^k`$, $`k=3,5,7,\mathrm{}`$, the solution of the Whitham equations has genus at most equal to one ,. On the $`xt`$, $`t0`$, plane we have only the domains $`D_0`$ and $`D_1`$. The one phase solution is defined within the cusp $`x_{}(k)t^{\frac{k}{k1}}<x<x_+(k)t^{\frac{k}{k1}}`$, where $`x_{}(k)<x_+(k)`$ are two real constants and $`t>0`$. The point $`x=0,t=0`$ is the point of gradient catastrophe of the zero-phase solution. The constants $`x_{}(k)`$ and $`x_+(k)`$ can be computed explicitly. From (5.10) and (5.11) we obtain
$$\mathrm{\Phi }_0(r,u)=\frac{1}{2\sqrt{ru}}_u^r\frac{f^{}(\xi )d\xi }{\sqrt{r\xi }}.$$
(5.32)
On the leading edge where $`u_1=u_2=v`$ and $`u_3=u`$, $`v>u`$, (5.20) is equivalent to the equations
$$_u^v\frac{(\xi u)f^{\prime \prime }(\xi )}{\sqrt{v\xi }}𝑑\xi =0,_u^v\frac{f^{}(\xi )6t}{\sqrt{v\xi }}𝑑\xi =0,x=6tu+f(u).$$
(5.33)
On the trailing edge where $`u_1=u`$, $`u_2=u_3=v`$, $`u>v`$, system (5.26) is equivalent, for $`g=0`$ to
$$_v^u(f^{}(\xi )6t)\sqrt{v\xi }𝑑\xi =0,_u^v\frac{f^{}(\xi )6t}{\sqrt{v\xi }}𝑑\xi =0,x=6tu+f(u).$$
(5.34)
Equations (5.33) and (5.34) have already been obtained in . Solving system (5.33) for $`t`$ when $`f(u)=u^k`$, $`k=3,5,7,\mathrm{}`$ we obtain $`x_{}(k)`$ ,
$`x_{}(k)=6{\displaystyle \frac{k1}{k}}(2z_{}(k)1)\left[{\displaystyle \frac{6}{k}}(1+2(k1)z_{}(k))\right]^{\frac{1}{k1}},k=3,5,7,\mathrm{},`$ (5.35)
where $`z_{}(k)>1`$ is the unique real solution of $`F(k+2,2,\frac{5}{2};z)=0`$. Here $`F(a,b,c;z)`$ is the hypergeometric series. The quantity $`x_+(k)`$ is obtained from (5.34)
$`x_+(k)=2{\displaystyle \frac{k1}{k}}(2z_+(k)3)\left[{\displaystyle \frac{2}{k}}(3+2(k1)z_+(k))\right]^{\frac{1}{k1}},k=3,5,7,\mathrm{},`$ (5.36)
where the number $`z_+(k)>1`$ is the unique real solution of the equation $`F(k+2,2,\frac{7}{2};z)=0`$. We give some numerical values: $`x_{}(3)=12\sqrt{3},x_{}(5)=16.85,x_{}(7)=16.21,x_{}(9)=16.09,x_+(3)=\frac{4}{3}\sqrt{\frac{5}{3}},x_+(5)=1.58,x_+(7)=1.61,x_+(9)=1.72`$.
Theorems 5.4, 5.6 and 5.7 can be generalized to the case of multiple phase transitions.
###### Theorem 5.9
The transition boundary between the domains $`D_g`$, $`g0`$ and $`D_{g+n}`$, $`n>1`$, having $`n_1`$ leading edges and $`n_2`$ trailing edges and $`n_3`$ points of gradient catastrophes, $`n_1+n_2+n_3=n`$, $`n_1,n_2,n_30`$, is described by the solution of the system
$`\{\begin{array}{ccc}_{v_k}\mathrm{\Phi }^g(v_k;\stackrel{}{u})=0,k=1,\mathrm{},n_1\hfill & & \\ \mathrm{\Phi }^g(v_k;\stackrel{}{u})6tϵ_{g0}=0,k=1,\mathrm{},n_1\hfill & & \\ \mathrm{\Phi }^g(y_l;\stackrel{}{u})6tϵ_{g0}=0,l=1,\mathrm{},n_2\hfill & & \\ _{y_l}^{u_{2j_l1}}(\mathrm{\Phi }^g(r;\stackrel{}{u})6tϵ_{g0})\mu (r)𝑑r=0,l=1,\mathrm{},n_2\hfill & & \\ _r\mathrm{\Phi }^g(r;\stackrel{}{u})|_{r=u_{j_m}}=0,m=1,\mathrm{},n_3\hfill & & \\ \mathrm{\Phi }^g(u_{j_m};\stackrel{}{u})6tϵ_{g0}=0m=1,\mathrm{},n_3\hfill & & \\ & & \\ x=\left[12t{\displaystyle \frac{P_1^g(r)}{P_0^g(r)}}+{\displaystyle \frac{R^g(r)}{P_0^g(r)}}\right]_{r=u_i},i=1,\mathrm{},2g+1,\hfill & & \end{array}`$ (5.45)
where $`v_k(u_{2j_k},u_{2j_k+1})`$, $`0j_kg`$, $`k=1,\mathrm{},n_1`$ and $`y_l(u_{2j_l1},u_{2j_l})`$, $`1j_lg+1`$, $`l=1,\mathrm{},n_2`$.
###### Remark 5.10
The transition between the $`g`$-phase solution and the $`(g+n)`$phase solution, $`n>1`$, is highly non generic (cfr. Fig 5.45). Indeed system (5.45) is a systems of $`2n+2g+1`$ equations in $`n_1+n_2+2g+3`$ unknowns $`v_1,\mathrm{},v_{n_1}`$, $`y_1,\mathrm{},y_{n_2}`$, $`u_1,\mathrm{},u_{2g+1}`$, $`x`$ and $`t0`$. Therefore if system (5.45) admits a real solution, $`n+n_32`$ variables are functions of all the others.
Theorem 5.9 includes all the degenerate cases that have been excluded in Theorems 5.4, 5.6 and 5.7. For example let us consider a transition with a double-leading edge described by the variables $`v_1`$, $`v_2`$ and $`u_1>u_2>\mathrm{}>u_{2g+1}`$ which solve (5.45) with $`n_1=2`$, $`n_2=0,n_3=0`$. In the limit $`v_1v_2=v`$ we get a single degenerate leading edge. The variables $`v`$, $`u_1>u_2>\mathrm{}>u_{2g+1}`$ satisfy (5.4) and the equations $`(_v)^2\mathrm{\Phi }^g(v;\stackrel{}{u})=0`$ and $`(_v)^3\mathrm{\Phi }^g(v;\stackrel{}{u})=0`$. Therefore we regard such degenerate single leading edge as a point of the boundary of the domains $`D_g`$ and $`D_{g+2}`$.
###### Proposition 5.11
If the domains $`D_g`$ and $`D_{g+1}`$ have a common boundary, then the function
$$\mathrm{\Phi }^g(r;\stackrel{}{u})6tϵ_{g0}$$
has at least $`g+2`$ real zeros in the $`r`$ variable for $`u_1(x,t)>u_2(x,t)>\mathrm{}>u_{2g+1}(x,t)`$ satisfying (5.4) and for $`t>0`$.
Proof: The domains $`D_g`$ and the $`D_{g+1}`$ have a common boundary if one of the systems (5.20), (5.26) or (5.30) is solvable for some $`t>0`$. We first consider the leading edge. For $`g=0`$ the statement is obvious from system (5.20). For $`g>0`$ system (5.20) can have a solution if the function $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ has a double zero at $`r=v`$ when $`v`$ belongs to the bands (5.15). Combining this observation with proposition 5.3 we immediately obtain the statement. As regarding the trailing edge, the theorem is obvious for $`g=0`$. For $`g>0`$ and $`v(a,u_{2g+1})`$ system (5.26) can be satisfied if the function $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ has at least two zeros in the interval $`(a,u_{2g+1})`$. Combining this observation with proposition 5.3 we obtain the statement. When $`v(u_{2j},u_{2j1})`$, $`1jg`$, (5.14) and (5.26) are satisfied if the function $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ has at least three real zeros in the interval $`(u_{2j},u_{2j1})`$, $`1jg`$, and changes sign at least once in each of the intervals $`(u_{2k},u_{2k1})`$, $`k=1,\mathrm{},g`$, and $`kj`$. Therefore $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ has at least $`g+2`$ real zeros at the trailing edge. If the point of the boundary between the domains $`D_g`$ and the $`D_{g+1}`$ corresponds to a point of gradient catastrophe of the $`g`$-phase solution the statement is obvious from system (5.30) and proposition 5.3. Therefore, on the phase transition boundary between the domains $`D_g`$ and $`D_{g+1}`$ the function $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ has at least $`g+2`$ real zeros in the $`r`$ variable, $`b>r>a`$, for $`b>u_1(x,t)>u_2(x,t)>\mathrm{}>u_{2g+1}(x,t)>a`$ and for $`t>0`$. . $`\mathrm{}`$ Proposition 5.11 can be generalized to multiple phase transitions.
###### Proposition 5.12
On the phase transition boundary between the domains $`D_g`$ and the $`D_{g+n}`$ the function
$$\mathrm{\Phi }^g(r;\stackrel{}{u})6tϵ_{g0}$$
has at least $`g+2n`$ real zeros in the $`r`$ variable for $`b>u_1(x,t)>u_2(x,t)>\mathrm{}>u_{2g+1}(x,t)>a`$ and for $`t>0`$.
The proof is analogous to that of proposition 5.11.
###### Lemma 5.13
If the smooth initial data $`x=f(u)|_{t=0}`$ satisfies (5.1), then the function
$$\mathrm{\Phi }^g(r;\stackrel{}{u})$$
has at most $`2Ng`$ real zeros (counting multiplicity) in the $`r`$ variable for $`0<g2N`$ and for any real $`x`$, $`t0`$ and $`b>u_1>u_2>\mathrm{}>u_{2g+1}>a`$.
Proof: for proving the lemma we need the following elementary result. If the real smooth function $`\xi (r)`$ satisfies the condition
$$\frac{d^m}{dr^m}\xi (r)>0,\mathrm{\hspace{0.33em}0}m\mathrm{I}\mathrm{N},$$
for all $`r`$ belonging to the domain of $`\xi `$, then $`\xi (r)`$ has at most $`m`$ real zero.
Using (5.1), (5.9) and (5.11) we obtain
$$\begin{array}{cc}& \frac{^{2Ng}}{r^{2Ng}}\mathrm{\Phi }^g(r;\stackrel{}{u})=_1^1\mathrm{}_1^1𝑑\xi _1\mathrm{}𝑑\xi _{2g+1}(1+\xi _{2g+1})^{2N}(1+\xi _{2g})^{2N\frac{1}{2}}\mathrm{}(1+\xi _2)^{2Ng+\frac{1}{2}}(1+\xi _1)^{2Ng}\hfill \\ & \\ & \times \frac{f^{(2N+1)}(\frac{1+\xi _{2g+1}}{2}(\mathrm{}(\frac{1+\xi _2}{2}(\frac{1+\xi _1}{2}r+\frac{1\xi _1}{2}u_1)+\frac{1\xi _2}{2}u_2)+\mathrm{})+\frac{1\xi _{2g+1}}{2}u_{2g+1})}{2^{(2g+1)(2Ng)}K\sqrt{(1\xi _2)(1\xi _2)\mathrm{}(1\xi _{2g+1})}}>0\hfill \end{array}$$
(5.46)
for $`g>0`$, for any real $`r`$ and for any fixed real $`u_1>u_2>\mathrm{}>u_{2g+1}`$ belonging to the interval $`(a,b)`$. Therefore $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ has at most $`2Ng`$ real zeros in the $`r`$ variable for $`0<g2N`$ and for any fixed real $`b>u_1>u_2>\mathrm{}>u_{2g+1}>a`$. $`\mathrm{}`$
###### Lemma 5.14
If the smooth monotone increasing initial data $`x=f(u)|_{t=0}`$ satisfies (5.1), then the function
$$\mathrm{\Phi }^0(r;u)6t$$
has at most $`2N`$ real zeros (counting multiplicity) in the $`r`$ variable, for any $`t0`$ and $`uu^{}`$, where $`u^{}`$ is the unique solution of the equation $`f^{(2N+1)}(u^{})=0`$. For what $`\mathrm{\Phi }^0(r;u^{})6t`$ is concerned we have two possibilities: either it has at most $`2N`$ real zeros (counting multiplicity) in the $`r`$ variable for any $`t0`$ and some of these zeros are distinct, or it has at most two real zeros in the $`r`$ variable for any $`t>0`$ except for $`t=t^{}0`$ where $`r=u^{}`$ is a zero of multiplicity higher than $`2N`$.
Proof: Since $`f^{}(u)0`$ and because of (5.1) we have two possibilities: either $`f^{}(u)`$ has at most $`N`$ simple minima and some of these minima are distinct, or $`f^{}(u)`$ is a non negative convex function with a minimum at $`u=u^{}`$, where $`u^{}`$ is the unique solution of the equation $`f^{(2N+1)}(u^{})=0`$. From the above considerations the lemma easily follows. $`\mathrm{}`$ We remark that when $`f^{}(u)`$ is a convex function, we can apply Tian’s Theorem 3.2. Indeed in such case we have $`f^{\prime \prime \prime }(u)>0`$ for all $`uu^{}`$ and therefore the solution of the Whitham equations has genus at most one.
We continue the proof of the second main theorem. In the following we exclude the case in which $`f^{}(u)`$ is a convex function. From proposition 5.12 we deduce that when the domains $`D_g`$ and $`D_m`$, $`m>g0`$, have a common boundary the function $`\mathrm{\Phi }^g(r;\stackrel{}{u})6tϵ_{g0}`$ defined in (5.11) has at least $`2mg`$ real zeros for some values of $`x`$ and $`t`$. From lemma 5.13 the function $`\mathrm{\Phi }^g(r;\stackrel{}{u})6tϵ_{g0}`$ defined in (5.11) has at most $`2Ng`$ real zeros for $`g2N`$ when the initial data $`f(u)`$ satisfies (5.1). Therefore
$$2Ng2mg\text{or}mN.$$
(5.47)
This shows that the set of domains $`\{D_g\}_{0gN}`$ does not have common boundaries with any of the domains in the set $`\{D_m\}_{m>N}`$. The second main theorem is proved. $`\mathrm{}`$
### 5.1 Phase transitions
In this subsection we prove theorems 5.4, 5.6, 5.7 and (5.9).
In the following we denote with a sup-script the genus of the quantity we are referring to. Whenever we omit the sup-script we are referring to genus $`g`$ quantity.
We denote with $`\sigma _k^{g+1}=\sigma _k^{g+1}(r,\delta ,v)`$ the normalized Abelian differential of the second kind with pole at infinity of order $`2k+2`$ defined on the surface $`𝒮_{g+1}`$ of genus $`g+1`$
$$\stackrel{~}{\mu }^2=(rv\sqrt{\delta })(rv+\sqrt{\delta })\mu ^2,v\mathrm{I}\mathrm{R},$$
$$\mu ^2=\underset{j=1}{\overset{2g+1}{}}(ru_j),u_1>u_2>\mathrm{}>u_{2g+1},$$
where $`vu_j,j=1,\mathrm{}2g+1`$, and $`0<\delta 1`$. We define the polynomials
$$P_k^{g+1}(r,\delta ,v)=\stackrel{~}{\mu }(r)\frac{\sigma _k^{g+1}(r,\delta ,v)}{dr}.$$
The expression of the polynomial in (5.3) for the Riemann invariants $`\stackrel{~}{u}_1=v+\sqrt{\delta }`$, $`\stackrel{~}{u}_2=v\sqrt{\delta }`$ and $`u_1>u_2>\mathrm{}>u_{2g+1}`$ reads
$$\begin{array}{cc}\hfill R^{g+1}(r,\delta ,v)& =2\underset{i=1}{\overset{2g+1}{}}_{u_i}q_{g+1}^{g+1}\underset{j=1,ji}{\overset{2g+1}{}}(ru_j)(r\stackrel{~}{u}_1)(r\stackrel{~}{u}_2)+2\mu ^2(r)(r\stackrel{~}{u}_1)_{\stackrel{~}{u}_2}q_{g+1}^{g+1}\hfill \\ \hfill +& 2\mu ^2(r)(r\stackrel{~}{u}_2)_{\stackrel{~}{u}_1}q_{g+1}^{g+1}+\underset{k=1}{\overset{g+1}{}}q_k^{g+1}\underset{l=1}{\overset{k}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{kl}^{g+1}P_{l1}^{g+1}(r,\stackrel{}{u},\delta )\hfill \end{array}$$
(5.48)
where $`q_k^{g+1}=q_k^{g+1}(\stackrel{~}{u}_1,\stackrel{~}{u}_2,u_1,\mathrm{},u_{2g+1})`$, $`k=1,\mathrm{},g+1`$ and $`\stackrel{~}{\mathrm{\Gamma }}_k^{g+1}=\stackrel{~}{\mathrm{\Gamma }}_k^{g+1}(\stackrel{~}{u}_1,\stackrel{~}{u}_2,u_1,\mathrm{},u_{2g+1})`$, $`k0`$. We will sometimes omit the explicit dependence of $`\sigma _k^{g+1}(r,\delta ,v)`$, $`P_k^{g+1}(r,\delta ,v)`$ and $`R^{g+1}(r,v,\delta )`$ on the parameters $`v`$ and $`\delta `$.
Proof of Theorem 5.4.
We write the $`(g+1)`$-phase solution (3.6) for the variables $`\stackrel{~}{u}_1=v+\sqrt{\delta }`$, $`\stackrel{~}{u}_2=v\sqrt{\delta }`$, $`u_1>u_2>\mathrm{}>u_{2g+1}`$ in the form
$`\{\begin{array}{ccc}0\hfill & =\hfill & {\displaystyle \frac{1}{\stackrel{~}{u}_1\stackrel{~}{u}_2}}\left(\left[12t{\displaystyle \frac{P_1^{g+1}(\stackrel{~}{u}_1,\delta )}{P_0^{g+1}(\stackrel{~}{u}_1,\delta )}}+{\displaystyle \frac{R^{g+1}(\stackrel{~}{u}_1,\delta )}{P_0^{g+1}(\stackrel{~}{u}_1,\delta )}}\right]\left[12t{\displaystyle \frac{P_1^{g+1}(\stackrel{~}{u}_2,\delta )}{P_0^{g+1}(\stackrel{~}{u}_2,\delta )}}+{\displaystyle \frac{R^{g+1}(\stackrel{~}{u}_2,\delta )}{P_0^{g+1}(\stackrel{~}{u}_2,\delta )}}\right]\right),\hfill \\ & & \\ x\hfill & =\hfill & \left[12t{\displaystyle \frac{P_1^{g+1}(r,\delta )}{P_0^{g+1}(r,\delta )}}+{\displaystyle \frac{R^{g+1}(r,\delta )}{P_0^{g+1}(r,\delta )}}\right]_{r=\stackrel{~}{u}_2},\hfill \\ & & \\ x\hfill & =\hfill & \left[12t{\displaystyle \frac{P_1^{g+1}(r,\delta )}{P_0^{g+1}(r,\delta )}}+{\displaystyle \frac{R^{g+1}(r,\delta )}{P_0^{g+1}(r,\delta )}}\right]_{r=u_i},i=1,\mathrm{},2g+1\hfill \end{array}`$ (5.54)
The reason to write the hodograph transform (3.6) in the form (5.54) is that system (5.54) is non degenerate at the phase transition $`\stackrel{~}{u}_1=\stackrel{~}{u}_2`$.
We show that system (5.54) reduces to system (5.20) when $`\stackrel{~}{u}_1=\stackrel{~}{u}_2`$ or $`\delta =0`$.
It is clear from Theorem 4.2 and proposition 5.2 that the last $`2g+1`$ equations in (5.54) in the limit $`\stackrel{~}{u}_1=\stackrel{~}{u}_2`$ reduce to those of system (5.20).
For computing the limit $`\stackrel{~}{u}_1\stackrel{~}{u}_2`$ of the first two equations in (5.54) we first need to study the behavior of the differentials $`\sigma _k^{g+1}(r,\delta ,v)`$, $`k0`$, when $`\delta 0`$. When $`v`$ belongs to one of the bands (5.15), the expansion of $`\sigma _k^{g+1}(r,\delta ,v)`$ for $`\delta 0`$ reads
$$\sigma _k^{g+1}(r,\delta ,v)=\sigma _k^g(r)+\frac{\delta }{2}\sigma _k^g(v)_v\omega _v^g(r)+O(\delta ^2)$$
(5.55)
where $`\sigma _k^g(r)`$ is the normalized differential of the second kind defined on $`𝒮_g`$ with pole at infinity of order $`2k+2`$, $`\sigma _k^g(v)=\frac{\sigma _k^g(r)}{dr}|_{r=v}`$ and $`_v={\displaystyle \frac{}{v}}`$. The differential $`\omega _v^g(r)`$ is the normalized Abelian differential of the third kind with poles at the points $`Q^\pm (v)=(v,\pm \mu (v))`$ with residue $`\pm 1`$ respectively. The explicit expression of $`\omega _v^g(r)`$ has been given in (2.14) The differential $`O(\delta ^2)/\delta ^2`$ has a pole at $`r=v`$ of order $`4`$ and zero residue.
From (5.55) we can get the expansion of the polynomial $`P_k^{g+1}(r,\delta ,v)=\stackrel{~}{\mu }(r){\displaystyle \frac{\sigma _k^{g+1}(r,\delta ,v)}{dr}}`$, namely
$$P_k^{g+1}(r,\delta ,v)=(rv)P_k^g(r)+\frac{\delta }{2}\sigma _k^g(v)\left(\mu ^{}(v)(rv)\underset{k=1}{\overset{g}{}}r^{gk}N_k^{}(v)\right)+O(\delta ^2),$$
(5.56)
where $`P_k^g(r)`$ has been defined in (2.7) and the $`N_k(v)`$’s have been defined in (2.18). From (5.56) we can evaluate the following
$$\frac{P_k^{g+1}(r,\delta )}{P_0^{g+1}(r,\delta )}|_{r=v\pm \sqrt{\delta }}=\frac{P_k^g(v)}{P_0^g(v)}\pm \sqrt{\delta }_v\left(\frac{P_k^g(v)}{P_0^g(v)}\right)+O(\delta ).$$
(5.57)
To evaluate the first two equations in (5.54) at the point of phase transition we need also the following relations
$$\begin{array}{cc}\hfill q_{g+1}^{g+1}(v,v,\stackrel{}{u})& =\mathrm{\Psi }^g(v;\stackrel{}{u})\hfill \\ \hfill \frac{}{\stackrel{~}{u}_i}q_{g+1}^{g+1}(\stackrel{~}{u}_1,\stackrel{~}{u}_2,\stackrel{}{u})|_{\stackrel{~}{u}_1=\stackrel{~}{u}_2}& =\frac{1}{2}\frac{}{v}\mathrm{\Psi }^g(v;\stackrel{}{u}),i=1,2,\hfill \\ \hfill \frac{^2}{\stackrel{~}{u}_1\stackrel{~}{u}_2}q_{g+1}^{g+1}(\stackrel{~}{u}_1,\stackrel{~}{u}_2,\stackrel{}{u})|_{\stackrel{~}{u}_1=\stackrel{~}{u}_2}& =\frac{1}{4}\frac{^2}{v^2}\mathrm{\Psi }^g(v;\stackrel{}{u}),\hfill \end{array}$$
(5.58)
where the function $`\mathrm{\Psi }^g(r;\stackrel{}{u})`$ has been defined in (5.10). From (4.48-4.49), (5.57) and (5.58) we obtain
$$\begin{array}{cc}\hfill \frac{R^{g+1}(\stackrel{~}{u}_i,\delta )}{P_0^{g+1}(\stackrel{~}{u}_i,\delta )}|_{\stackrel{~}{u}_1=\stackrel{~}{u}_2}& =\frac{2\mu ^2(v)_v\mathrm{\Psi }^g(v;\stackrel{}{u})}{P_0^g(v)}+\mathrm{\Psi }^g(v;\stackrel{}{u})\underset{l=1}{\overset{g+1}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{g+1l}^g\frac{P_{l1}^g(v)}{P_0^g(v)}\hfill \\ & +\underset{k=1}{\overset{g}{}}q_k^g(\stackrel{}{u})\underset{l=1}{\overset{k}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{kl}^g\frac{P_{l1}^g(v)}{P_0^g(v)},i=1,2.\hfill \end{array}$$
(5.59)
From (2.21), (5.13) and using the definition of $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ in (5.11) the above relation can be written in the form
$$\frac{R^{g+1}(\stackrel{~}{u}_i,\delta )}{P_0^{g+1}(\stackrel{~}{u}_i,\delta )}|_{\stackrel{~}{u}_1=\stackrel{~}{u}_2}=\frac{2\mu ^2(v)\mathrm{\Phi }^g(v;\stackrel{}{u})+R^g(v,\stackrel{}{u})}{P_0^g(v)},i=1,2,$$
(5.60)
where $`R^g(r,\stackrel{}{u})`$ has been defined in (5.3). We need to consider also the quantity
$$\begin{array}{cc}& \left[\frac{1}{\stackrel{~}{u}_1\stackrel{~}{u}_2}\left(\frac{R^{g+1}(\stackrel{~}{u}_1,\delta )}{P_0^{g+1}(\stackrel{~}{u}_1,\delta )}\frac{R^{g+1}(\stackrel{~}{u}_2,\delta )}{P_0^{g+1}(\stackrel{~}{u}_2,\delta )}\right)\right]_{\stackrel{~}{u}_1=\stackrel{~}{u}_2}=8\left(\sqrt{\delta }\frac{\mu ^2(v+\sqrt{\delta })}{P_0^{g+1}(v+\sqrt{\delta },\delta )}_{\stackrel{~}{u}_1}_{\stackrel{~}{u}_2}q_{g+1}^{g+1}(\stackrel{~}{u}_1,\stackrel{~}{u}_2,\stackrel{}{u})\right)|_{\delta =0}\hfill \\ \hfill +& 2\left(\left(\frac{\mu ^2(v+\sqrt{\delta })}{P_0^{g+1}(v+\sqrt{\delta },\delta )}\frac{\mu ^2(v\sqrt{\delta })}{P_0^{g+1}(v\sqrt{\delta },\delta )}\right)_{\stackrel{~}{u}_2}q_{g+1}^{g+1}(\stackrel{~}{u}_1,\stackrel{~}{u}_2,\stackrel{}{u})\right)|_{\delta =0}\hfill \\ \hfill +& \underset{k=1}{\overset{g}{}}q_k^g(\stackrel{}{u})\underset{l=1}{\overset{k}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{kl}^g(\stackrel{}{u})\left(\frac{1}{2\sqrt{\delta }}\left(\frac{P_{l1}^{g+1}(v+\sqrt{\delta },\delta )}{P_0^{g+1}(v+\sqrt{\delta },\delta )}\frac{P_{l1}^{g+1}(v\sqrt{\delta },\delta )}{P_0^{g+1}(v\sqrt{\delta },\delta )}\right)\right)|_{\delta =0}\hfill \\ \hfill +& q_{g+1}^{g+1}(v,v,\stackrel{}{u})\underset{l=1}{\overset{g+1}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{g+1l}^g(\stackrel{}{u})\left(\frac{1}{2\sqrt{\delta }}\left(\frac{P_{l1}^{g+1}(v+\sqrt{\delta },\delta )}{P_0^{g+1}(v+\sqrt{\delta },\delta )}\frac{P_{l1}^{g+1}(v\sqrt{\delta },\delta )}{P_0^{g+1}(v\sqrt{\delta },\delta )}\right)\right)|_{\delta =0}\hfill \end{array}$$
where we have used (4.48-4.49) to obtain the right hand side. Using (2.21), (5.13), (5.57) and (5.58) the above reduces to the form
$$\left[\frac{1}{\stackrel{~}{u}_1\stackrel{~}{u}_2}\left(\frac{R^{g+1}(\stackrel{~}{u}_1,\delta )}{P_0^{g+1}(\stackrel{~}{u}_1,\delta )}\frac{R^{g+1}(\stackrel{~}{u}_2,\delta )}{P_0^{g+1}(\stackrel{~}{u}_2,\delta )}\right)\right]_{\stackrel{~}{u}_1=\stackrel{~}{u}_2}=\frac{}{_v}\left(\frac{2\mu ^2(v)\mathrm{\Phi }^g(v)+R^g(v)}{P_0^g(v)}\right).$$
(5.61)
From (5.57), (5.60) and (5.61), the system (5.54) on the point of phase transition $`\stackrel{~}{u}_1=\stackrel{~}{u}_2=v`$ reads
$`\{\begin{array}{ccc}0\hfill & =\hfill & {\displaystyle \frac{}{v}}\left({\displaystyle \frac{12tP_1^g(v)+R^g(v,\stackrel{}{u})+2\mu ^2(v)\mathrm{\Phi }^g(v;\stackrel{}{u})}{P_0^g(v,\stackrel{}{u})}}\right)\hfill \\ & & \\ x\hfill & =\hfill & {\displaystyle \frac{12tP_1^g(v,\stackrel{}{u})+R^g(v,\stackrel{}{u})+2\mu ^2(v)\mathrm{\Phi }^g(v;\stackrel{}{u})}{P_0^g(v,\stackrel{}{u})}}\hfill \\ & & \\ x\hfill & =\hfill & \left[12t{\displaystyle \frac{P_1^g(r,\stackrel{}{u})}{P_0^g(r,\stackrel{}{u})}}+{\displaystyle \frac{R^g(r)}{P_0^g(r,\stackrel{}{u})}}\right]_{r=u_i},i=1,\mathrm{},2g+1\hfill \end{array}`$ (5.67)
From proposition 3.3 the last $`2g+1`$ equations in (5.67) are equivalent to the condition
$$xP_0^g(r)12tP_1^g(r)+R^g(r)0,g>0.$$
Therefore system (5.67) is equivalent to system (5.20) for $`g>0`$. For $`g=0`$ substituting (3.2) in (5.67) we obtain
$`\{\begin{array}{ccc}0=_v(6tu+f(u)+2(vu)(\mathrm{\Phi }_0(v,u)6t))\hfill & & \\ & & \\ x=6tu+f(u)+2(vu)(\mathrm{\Phi }_0(v,u)6t)\hfill & & \\ & & \\ x=6tu+f(u)\hfill & & \end{array}`$ (5.73)
where $`u_1=u`$ and $`f(u)`$ is the initial data. It is clear that (5.73) is equivalent to (5.20) for $`g=0`$.
$`\mathrm{}`$
Proof of Theorem 5.6.
For getting the equations determining the trailing edge we have to repeat all the above calculation in a slightly different way. We write the $`(g+1)`$-phase solution for the variables $`\stackrel{~}{u}_1=v+\sqrt{\delta }`$, $`\stackrel{~}{u}_2=v\sqrt{\delta }`$, $`u_1>u_2>\mathrm{}>u_{2g+1}`$ in the form
$`\{\begin{array}{ccc}0\hfill & =\hfill & {\displaystyle \frac{1}{\mathrm{log}(\frac{\stackrel{~}{u}_1\stackrel{~}{u}_2}{2})^2(\stackrel{~}{u}_1\stackrel{~}{u}_2)}}\left(12t\left({\displaystyle \frac{P_1^{g+1}(\stackrel{~}{u}_2,\delta )}{P_0^{g+1}(\stackrel{~}{u}_2,\delta )}}{\displaystyle \frac{P_1^{g+1}(\stackrel{~}{u}_1,\delta )}{P_0^{g+1}(\stackrel{~}{u}_1,\delta )}}\right){\displaystyle \frac{R^{g+1}(\stackrel{~}{u}_1,\delta )}{P_0^{g+1}(\stackrel{~}{u}_1,\delta )}}+{\displaystyle \frac{R^{g+1}(\stackrel{~}{u}_2,\delta )}{P_0^{g+1}(\stackrel{~}{u}_2,\delta )}}\right),\hfill \\ & & \\ x\hfill & =\hfill & \left[12t{\displaystyle \frac{P_1^{g+1}(r,\delta )}{P_0^{g+1}(r,\delta )}}+{\displaystyle \frac{R^{g+1}(r,\delta )}{P_0^{g+1}(r,\delta )}}\right]_{r=\stackrel{~}{u}_2},\hfill \\ & & \\ x\hfill & =\hfill & \left[12t{\displaystyle \frac{P_1^{g+1}(r,\delta )}{P_0^{g+1}(r,\delta )}}+{\displaystyle \frac{R^{g+1}(r,\delta )}{P_0^{g+1}(r,\delta )}}\right]_{r=u_i},i=1,\mathrm{},2g+1\hfill \end{array}`$ (5.79)
The reason to write the hodograph transform in the form (5.79) is that system (5.79) is non degenerate at the phase transition $`\stackrel{~}{u}_1=\stackrel{~}{u}_2=v`$ when $`v(u_{2j},u_{2j1}),\mathrm{\hspace{0.33em}\hspace{0.33em}1}jg+1,u_{g+2}=a`$.
From Theorem 4.2 and proposition 5.2 we deduce that the last $`2g+1`$ equations in (5.79) in the limit $`\stackrel{~}{u}_1=\stackrel{~}{u}_2`$ reduce to those of system (5.26).
Next we investigate the behavior of the Abelian differentials of the second kind $`\sigma _k^{g+1}(r,\delta ,v)`$, $`k0`$, in the limit $`\delta 0`$, when $`v`$ belongs to one of the gaps (5.16). We have that
$$\sigma _k^{g+1}(r,\delta ,v)\sigma _k^g(r)\frac{1}{\mathrm{log}\delta }\omega _v^g(r)_{Q^{}(v)}^{Q^+(v)}\sigma _k^g(\xi ),$$
(5.80)
where $`\omega _v^g(r)`$ has been defined in (2.14) and $`Q^\pm (v)=(v,\pm \mu (v))`$. From (5.80) we can obtain the expansion for the polynomial $`P_k^{g+1}(r,\delta ,v)=\stackrel{~}{\mu }(r){\displaystyle \frac{\sigma _k^{g+1}(r,\delta ,v)}{dr}}`$, namely
$$P_k^{g+1}(r,\delta ,v)(rv)P_k^g(r)\frac{1}{\mathrm{log}\delta }(rv)\mu (r)\omega _v^g(r)_{Q^{}(v)}^{Q^+(v)}\sigma _k^g(\xi ),$$
(5.81)
so that
$$\frac{P_k^{g+1}(r,\delta )}{P_0^{g+1}(r,\delta )}|_{r=v\pm \sqrt{\delta }}\frac{_{Q^{}(v)}^{Q^+(v)}\sigma _k^g(\xi )}{_{Q^{}(v)}^{Q^+(v)}\sigma _0^g(\xi )}\left(1\sqrt{\delta }\mathrm{log}\delta \left(\frac{\sigma _k^g(v)}{_{Q^{}(v)}^{Q^+(v)}\sigma _k^g(\xi )}\frac{\sigma _0^g(v)}{_{Q^{}(v)}^{Q^+(v)}\sigma _0^g(\xi )}\right)\right)$$
(5.82)
Using (5.3), (4.48-4.49) and (5.82) we have that
$$\begin{array}{cc}\hfill \frac{R^{g+1}(\stackrel{~}{u}_i,\delta )}{P_0^{g+1}(\stackrel{~}{u}_i,\delta )}|_{\stackrel{~}{u}_1=\stackrel{~}{u}_2}=& \underset{k=1}{\overset{g}{}}q_k^g(\stackrel{}{u})\underset{l=1}{\overset{k}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{kl}^g(\stackrel{}{u})\frac{_{Q^{}(v)}^{Q^+(v)}\sigma _k^g(\xi )}{_{Q^{}(v)}^{Q^+(v)}\sigma _0^g(\xi )}\hfill \\ \hfill +& q_{g+1}^{g+1}(v,v,\stackrel{}{u})\underset{l=1}{\overset{g+1}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{g+1l}^g(\stackrel{}{u})\frac{_{Q^{}(v)}^{Q^+(v)}\sigma _k^g(\xi )}{_{Q^{}(v)}^{Q^+(v)}\sigma _0^g(\xi )},i=1,2.\hfill \end{array}$$
(5.83)
Adding and subtracting
$$\frac{2}{_{Q^{}(v)}^{Q^+(v)}\sigma _0^g(\xi )}\underset{i=1}{\overset{2g+1}{}}_{Q^{}(v)}^{Q^+(v)}\frac{\mu (\xi )d\xi }{\xi u_i}_{u_i}q_g(\stackrel{}{u}),$$
to (5.83) and using (2.21) and (5.13) we obtain
$$\frac{R^{g+1}(\stackrel{~}{u}_i,\delta )}{P_0^{g+1}(\stackrel{~}{u}_i,\delta )}|_{\stackrel{~}{u}_1=\stackrel{~}{u}_2}=\frac{1}{_{Q^{}(v)}^{Q^+(v)}\sigma _0^g(\xi )}_{Q^{}(v)}^{Q^+(v)}\frac{2\mu ^2(\xi )\mathrm{\Phi }^g(\xi ;\stackrel{}{u})+R^g(\xi )}{\mu (\xi )}𝑑\xi ,i=1,2$$
(5.84)
where $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ and $`R^g(r)`$ have been defined in (5.11) and (5.3) respectively. We observe that
$$_{Q^{}(v)}^{Q^+(v)}\sigma _0^g(\xi )0$$
for all $`v`$ belonging to the gaps (5.16).
We need to consider also the quantity in the first equation of (5.79) namely
$$\begin{array}{cc}& \left(\frac{1}{\mathrm{log}(\frac{\stackrel{~}{u}_1\stackrel{~}{u}_2}{2})^2(\stackrel{~}{u}_1\stackrel{~}{u}_2)}\left(\frac{R^{g+1}(\stackrel{~}{u}_1,\delta )}{P_0^{g+1}(\stackrel{~}{u}_1,\delta )}\frac{R^{g+1}(\stackrel{~}{u}_2,\delta )}{P_0^{g+1}(\stackrel{~}{u}_2,\delta )}\right)\right)|_{\stackrel{~}{u}_1=\stackrel{~}{u}_2}=\hfill \\ \hfill & 8\left(\frac{\sqrt{\delta }}{\mathrm{log}\delta }\frac{\mu ^2(v+\sqrt{\delta })}{P_0^{g+1}(v+\sqrt{\delta },\delta )}_{\stackrel{~}{u}_1}_{\stackrel{~}{u}_2}q_{g+1}^{g+1}(\stackrel{~}{u}_1,\stackrel{~}{u}_2,\stackrel{}{u})\right)|_{\delta =0}\hfill \\ \hfill & 2\left(\left(\frac{\mu ^2(v+\sqrt{\delta })}{\mathrm{log}\delta P_0^{g+1}(v+\sqrt{\delta },\delta )}+\frac{\mu ^2(v\sqrt{\delta })}{\mathrm{log}\delta P_0^{g+1}(v\sqrt{\delta },\delta )}\right)_{\stackrel{~}{u}_2}q_{g+1}^{g+1}(\stackrel{~}{u}_1,\stackrel{~}{u}_2,\stackrel{}{u})\right)|_{\delta =0}\hfill \\ \hfill & \underset{k=1}{\overset{g}{}}q_k^g(\stackrel{}{u})\underset{l=1}{\overset{k}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{kl}^g(\stackrel{}{u})\left(\frac{1}{2\sqrt{\delta }\mathrm{log}\delta }\left(\frac{P_{l1}^{g+1}(v+\sqrt{\delta },\delta )}{P_0^{g+1}(v+\sqrt{\delta },\delta )}\frac{P_{l1}^{g+1}(v\sqrt{\delta },\delta )}{P_0^{g+1}(v\sqrt{\delta },\delta )}\right)\right)|_{\delta =0}\hfill \\ \hfill & q_{g+1}^{g+1}(v,v,\stackrel{}{u})\underset{l=1}{\overset{g+1}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{g+1l}^g(\stackrel{}{u})\left(\frac{1}{2\sqrt{\delta }\mathrm{log}\delta }\left(\frac{P_{l1}^{g+1}(v+\sqrt{\delta },\delta )}{P_0^{g+1}(v+\sqrt{\delta },\delta )}\frac{P_{l1}^{g+1}(v\sqrt{\delta },\delta )}{P_0^{g+1}(v\sqrt{\delta },\delta )}\right)\right)|_{\delta =0}\hfill \end{array}$$
Using (2.21), (5.13), (5.82) and (5.58) we simplify the above relation to the form
$$\begin{array}{cc}& \left(\frac{1}{\mathrm{log}(\frac{\stackrel{~}{u}_1\stackrel{~}{u}_2}{2})^2(\stackrel{~}{u}_1\stackrel{~}{u}_2)}(\frac{R^{g+1}(\stackrel{~}{u}_1,\delta )}{P_0^{g+1}(\stackrel{~}{u}_1,\delta )}\frac{R^{g+1}(\stackrel{~}{u}_2,\delta )}{P_0^{g+1}(\stackrel{~}{u}_2,\delta )})\right)|_{\stackrel{~}{u}_1=\stackrel{~}{u}_2}=\frac{1}{\mu (v)_{Q^{}(v)}^{Q^+(v)}\sigma _0^g(\xi )}\times \hfill \\ & (2\mu ^2(v)_v\mathrm{\Psi }^g(v;\stackrel{}{u})+\underset{k=1}{\overset{g}{}}q_k^g(\stackrel{}{u})\underset{l=1}{\overset{k}{}}(2l1)\stackrel{~}{\mathrm{\Gamma }}_{kl}^g(\stackrel{}{u})(P_{l1}^g(v)P_0^g(v)\frac{_{Q^{}(v)}^{Q^+(v)}\sigma _{l1}^g(\xi )}{_{Q^{}(v)}^{Q^+(v)}\sigma _0^g(\xi )})\hfill \\ \hfill +& (2\mu (v)\mu ^{}(v)4\mu (v)\frac{P_0^g(v,\stackrel{}{u})}{_{Q^{}(v)}^{Q^+(v)}\sigma _0^g(\xi )})\mathrm{\Psi }^g(v;\stackrel{}{u}))\hfill \end{array}$$
Using repeatedly (5.13) we can write the above relation in the form
$$\begin{array}{cc}& \left(\frac{1}{\mathrm{log}(\frac{\stackrel{~}{u}_1\stackrel{~}{u}_2}{2})^2(\stackrel{~}{u}_1\stackrel{~}{u}_2)}(\frac{R^{g+1}(\stackrel{~}{u}_1,\delta )}{P_0^{g+1}(\stackrel{~}{u}_1,\delta )}\frac{R^{g+1}(\stackrel{~}{u}_2,\delta )}{P_0^{g+1}(\stackrel{~}{u}_2,\delta )})\right)|_{\stackrel{~}{u}_1=\stackrel{~}{u}_2}=\frac{1}{\mu (v)_{Q^{}(v)}^{Q^+(v)}\sigma _0^g(\xi )}\times \hfill \\ & \left(2\mu ^2(v)\mathrm{\Phi }^g(v;\stackrel{}{u})+R^g(v)\frac{P_0^g(v)}{_{Q^{}(v)}^{Q^+(v)}\sigma _0^g(\xi )}_{Q^{}(v)}^{Q^+(v)}\frac{2\mu ^2(\xi )\mathrm{\Phi }^g(\xi ;\stackrel{}{u})+R^g(\xi )}{\mu (\xi )}𝑑\xi \right),\hfill \end{array}$$
(5.85)
where $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ and $`R^g(r)`$ have been defined in (5.11) and (5.3) respectively.
From (5.82), (5.84) and (5.85) system (5.79) can be reduced to the form
$`\{\begin{array}{ccc}0=2\mu ^2(v)\mathrm{\Phi }^g(v;\stackrel{}{u})+Z^g(v)+{\displaystyle \frac{P_0^g(v)}{_{Q^{}(v)}^{Q^+(v)}\sigma _0^g(\xi )}}{\displaystyle _{Q^{}(v)}^{Q^+(v)}}{\displaystyle \frac{2\mu ^2(\xi )\mathrm{\Phi }^g(\xi ;\stackrel{}{u})+Z^g(\xi )}{\mu (\xi )}}𝑑\xi \hfill & & \\ & & \\ 0={\displaystyle _{Q^{}(v)}^{Q^+(v)}}{\displaystyle \frac{2\mu ^2(\xi )\mathrm{\Phi }^g(\xi ;\stackrel{}{u})+Z^g(\xi )}{\mu (\xi )}}𝑑\xi \hfill & & \\ & & \\ x=\left[t{\displaystyle \frac{P_1^g(r)}{P_0^g(r)}}+{\displaystyle \frac{R^g(r)}{P_0^g(r)}}\right]_{r=u_i},i=1,\mathrm{},2g+1\hfill & & \end{array}`$ (5.91)
where the polynomial $`Z^g(r)`$ has been defined in (5.2). From proposition 3.3 the last $`2g+1`$ equations in (5.91) are equivalent to the condition
$$Z^g(r)0,g>0.$$
Therefore system (5.91) is equivalent for $`g>0`$ to the system
$`\{\begin{array}{ccc}0=\mathrm{\Phi }^g(v;\stackrel{}{u})\hfill & & \\ & & \\ 0={\displaystyle _{Q^{}(v)}^{Q^+(v)}}\mu (\xi )\mathrm{\Phi }^g(\xi ;\stackrel{}{u})𝑑\xi \hfill & & \\ & & \\ x=\left[t{\displaystyle \frac{P_1^g(r)}{P_0^g(r)}}+{\displaystyle \frac{R^g(r)}{P_0^g(r)}}\right]_{r=u_i},i=1,\mathrm{},2g+1.\hfill & & \end{array}`$ (5.97)
Because of (5.14), when $`v(u_{2j},u_{2j1}),1jg+1,u_{2g+2}=a`$, we can split the integral of the second equation of (5.97) in the form
$$\begin{array}{cc}\hfill _{Q^{}(v)}^{Q^+(v)}\mu (\xi )\mathrm{\Phi }^g(\xi ;\stackrel{}{u})𝑑\xi =& _{Q^{}(v)}^{u_{2j1}}\mu (\xi )\mathrm{\Phi }^g(\xi ;\stackrel{}{u})𝑑\xi +_{u_{2j1}}^{Q^+(v)}\mu (\xi )\mathrm{\Phi }^g(\xi ;\stackrel{}{u})𝑑\xi \hfill \\ \hfill =& 2_{u_{2j1}}^{Q^+(v)}\mu (\xi )\mathrm{\Phi }^g(\xi ;\stackrel{}{u})𝑑\xi .\hfill \end{array}$$
Therefore (5.97) is equivalent to (5.26) for $`g>0`$.
For $`g=0`$ substituting (3.2) in (5.91) it is easy to check that we obtain a system equivalent to (5.26). $`\mathrm{}`$
Proof of Theorem 5.7.
The equations describing the point of gradient catastrophe of the $`g`$-phase solution can be obtained either considering the limit of the $`(g+1)`$phase solution when three Riemann invariants coalesce, or supposing that one of the $`2g+1`$ distinct Riemann invariants of the $`g`$-phase solution has a vertical inflection point for $`t>0`$. For proving Theorem 5.7 we follow the latter possibility.
On the solution of (3.6) $`_xu_l(x,t)=(_{u_l}(\lambda _lt+w_l))^1`$ therefore a point of gradient catastrophe of the $`g`$-phase solution is determined by the system
$`\{\begin{array}{ccc}_{u_l}(\lambda _lt+w_l)=0\hfill & & \\ & & \\ (_{u_l})^2(\lambda _lt+w_l)=0\hfill & & \\ & & \\ x=\left[t{\displaystyle \frac{P_1^g(r)}{P_0^g(r)}}+{\displaystyle \frac{R^g(r)}{P_0^g(r)}}\right]_{r=u_i},i=1,\mathrm{},2g+1\hfill & & \end{array}`$ (5.103)
where $`1l2g+1`$. We show that system (5.103) is equivalent to system (5.30). For the purpose we compute explicitly the derivative $`_{u_l}(\lambda _lt+w_l)`$. From a generalization of a result in we obtain:
$$\frac{}{u_l}\frac{P_k^g(u_l)}{P_0^g(u_l)}=\frac{1}{2}\frac{}{r}\frac{P_k^g(r)}{P_0^g(r)}|_{r=u_l},$$
(5.104)
so that
$$\begin{array}{cc}& \frac{}{u_l}(\lambda _l(\stackrel{}{u})t+w_l(\stackrel{}{u}))=6t\frac{}{r}\frac{P_1^g(r)}{P_0^g(r)}|_{r=u_l}+\frac{_{\stackrel{k=1}{kl}}^{2g+1}_{\stackrel{m=1}{mk,l}}^{2g+1}(u_lu_m)}{P_0^g(u_l)}_{u_l}q_g(\stackrel{}{u})\hfill \\ \hfill +& \frac{_{\stackrel{m=1}{ml}}^{2g+1}(u_lu_m)}{P_0^g(u_l)}(_{u_l})^2q_g(\stackrel{}{u})\frac{_{\stackrel{m=1}{ml}}^{2g+1}(u_lu_m)}{(P_0^g(u_l))^2}_{u_l}q_g(\stackrel{}{u})_{u_l}P_0^g(u_l)\hfill \\ \hfill +& \frac{1}{2}\underset{n=1}{\overset{g}{}}(2n1)\frac{}{r}\frac{P_{l1}^g(r)}{P_0^g(r)}|_{r=u_l}\underset{m=n}{\overset{g}{}}q_m(\stackrel{}{u})\stackrel{~}{\mathrm{\Gamma }}_{mn}+\underset{n=1}{\overset{g}{}}(2n1)\frac{P_{n1}^g(u_l)}{P_0^g(u_l)}\underset{m=n}{\overset{g}{}}_{u_l}(q_m(\stackrel{}{u})\stackrel{~}{\mathrm{\Gamma }}_{mn}).\hfill \end{array}$$
(5.105)
In the above relation we need to compute the derivative
$$_{u_l}P_0^g(u_l)=_rP_0^g(r)|_{r=u_l}+_{u_l}P_0^g(r)|_{r=u_l},$$
where $`P_0^g(r)=r^g+\alpha _1^0r^{g1}+\mathrm{}+\alpha _g^0`$. For computing the derivatives of the normalization constants $`\alpha _1^0,\alpha _2^0,\mathrm{},\alpha _g^0`$ in $`P_0^g(r)`$ we need the following proposition.
###### Proposition 5.15
Let $`\omega _1(r)`$ and $`\omega _2(r)`$ two normalized Abelian differentials on $`𝒮_g`$. Let be $`\xi ={\displaystyle \frac{1}{\sqrt{r}}}`$ the local coordinate at infinity and
$$\omega _1=\underset{k}{}a_k^1\xi ^kd\xi ,\omega _2=\underset{k}{}a_k^2\xi ^kd\xi .$$
Define the bilinear product
$$V_{\omega _1\omega _2}=\underset{k0}{}\frac{a_{k2}^1a_k^2}{k+1},$$
then
$$\frac{}{u_i}V_{\omega _1\omega _2}=\underset{[}{Res}r=u_i]\frac{\omega _1(r)\omega _2(r)}{dr},i=1,\mathrm{},2g+1,$$
(5.106)
where $`\underset{[}{Res}r=u_i]{\displaystyle \frac{\omega _1(r)\omega _2(r)}{dr}}`$ is the residue of the differential $`{\displaystyle \frac{\omega _1(r)\omega _2(r)}{dr}}`$ evaluated at $`r=u_i`$.
Applying the above proposition to $`\sigma _0`$ and $`\sigma _k`$, $`k=0,\mathrm{}g1`$ we obtain after non trivial simplifications
$$\begin{array}{cc}\hfill \frac{}{u_l}\left(\begin{array}{c}\alpha _1^0\\ \alpha _2^0\\ \mathrm{}\\ \alpha _{g1}^0\\ \alpha _g^0\end{array}\right)=& \frac{1}{2}\left(\begin{array}{c}1\\ u_l\\ \mathrm{}\\ u_l^{g1}\\ u_l^{g1}\end{array}\right)\frac{1}{2}\left(\begin{array}{ccccc}0& 0& 0& \mathrm{}& 0\\ 1& 0& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ u_l^{g3}& u_l^{g4}& \mathrm{}& 0& 0\\ u_l^{g2}& u_l^{g3}& \mathrm{}& 1& 0\end{array}\right)\left(\begin{array}{c}\alpha _1^0\\ \alpha _2^0\\ \mathrm{}\\ \alpha _{g1}^0\\ \alpha _g^0\end{array}\right)\hfill \\ \hfill +& \frac{1}{2}\frac{P_0^g(u_l)}{_{k=1,kl}^{2g+1}(u_lu_k)}\left(\begin{array}{ccccc}\stackrel{~}{\mathrm{\Gamma }}_0& 0& 0& \mathrm{}& 0\\ \stackrel{~}{\mathrm{\Gamma }}_1& \stackrel{~}{\mathrm{\Gamma }}_0& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \stackrel{~}{\mathrm{\Gamma }}_{g2}& \stackrel{~}{\mathrm{\Gamma }}_{g3}& \mathrm{}& \stackrel{~}{\mathrm{\Gamma }}_0& 0\\ \mathrm{\Gamma }_{g1}& \stackrel{~}{\mathrm{\Gamma }}_{g2}& \mathrm{}& \stackrel{~}{\mathrm{\Gamma }}_1& \stackrel{~}{\mathrm{\Gamma }}_0\end{array}\right)\left(\begin{array}{c}1\\ 3P_1^g(u_l)\\ \mathrm{}\\ (2g3)P_{g2}^g(u_l)\\ (2g1)P_{g1}^g(u_l)\end{array}\right)\hfill \end{array}$$
(5.107)
where the $`\stackrel{~}{\mathrm{\Gamma }}_k`$’s have been defined in (2.19). From the above formula we obtain
$$\begin{array}{cc}\hfill _{u_l}P_0^g(u_l)=& _rP_0^g(r)|_{r=u_l}+_{u_l}P_0^g(r)|_{r=u_l}\hfill \\ \hfill =& \frac{1}{2}_rP_0^g(r)|_{r=u_l}+\frac{1}{2}\frac{P_0^g(u_l)}{_{k=1,kl}^{2g+1}(u_lu_k)}\underset{n=1}{\overset{g}{}}(2n1)P_{n1}^g(u_l)\underset{m=n}{\overset{g}{}}u_l^m\stackrel{~}{\mathrm{\Gamma }}_{mn}.\hfill \end{array}$$
(5.108)
Using the relations (4.39), (4.43), (5.104) and (5.108) we simplify (5.105) to the form
$$\begin{array}{cc}\hfill \frac{}{u_l}(\lambda _l(\stackrel{}{u})t+& w_l(\stackrel{}{u}))=\frac{Z^g(u_l)}{2(P_0^g(u_l))^2}_rP^g_0(r)|_{r=u_l}+\frac{_{\stackrel{k=1}{kl}}^{2g+1}_{\stackrel{m=1}{mk,l}}^{2g+1}(u_lu_m)}{P_0^g(u_l)}_{u_l}q_g(\stackrel{}{u})\hfill \\ \hfill +& \frac{1}{2P_0^g(u_l)}_r\left(\underset{n=1}{\overset{g}{}}(2n1)P_{n1}^g(r)\underset{m=n}{\overset{g}{}}q_m\stackrel{~}{\mathrm{\Gamma }}_{mn}xP_0^g(r)12tP_1^g(r)\right)|_{r=u_l},\hfill \end{array}$$
(5.109)
where the polynomial $`Z^g(r)`$ has been defined in (5.2). Applying in the second term of (5.109) the relations (5.13) we can rewrite (5.109) in the compact form
$$\frac{}{u_l}(\lambda _l(\stackrel{}{u})t+w_l(\stackrel{}{u}))=\frac{}{r}\frac{Z^g(r)}{2P_0^g(r)}|_{r=u_l}+\frac{_{\stackrel{m=1}{ml}}^{2g+1}(u_lu_m)}{P_0^g(u_l)}\mathrm{\Phi }^g(u_l;\stackrel{}{u}),$$
(5.110)
where $`\mathrm{\Phi }^g(r;\stackrel{}{u})`$ has been defined in (5.11). From proposition 3.3 the last $`2g+1`$ equations in (5.103) are equivalent to the condition
$$Z^g(r)0,g>0.$$
Therefore we can simplify (5.110) to the form
$$\frac{}{u_l}(\lambda _l(\stackrel{}{u})t+w_l(\stackrel{}{u}))=\frac{_{\stackrel{m=1}{ml}}^{2g+1}(u_lu_m)}{P_0^g(u_l)}\mathrm{\Phi }^g(u_l;\stackrel{}{u}),$$
(5.111)
when $`u_1>u_2>\mathrm{}>u_{2g+1}`$ satisfy the $`g`$-phase Whitham equations. As regarding the second derivative $`{\displaystyle \frac{^2}{u_l^2}}(\lambda _l(\stackrel{}{u})t+w_l(\stackrel{}{u}))`$, from (5.111) we obtain
$$\frac{^2}{u_l^2}(\lambda _l(\stackrel{}{u})t+w_l(\stackrel{}{u}))=\frac{}{u_l}\left(\frac{_{\stackrel{m=1}{ml}}^{2g+1}(u_lu_m)}{P_0^g(u_l)}\right)\mathrm{\Phi }^g(u_l;\stackrel{}{u})+\frac{_{\stackrel{m=1}{ml}}^{2g+1}(u_lu_m)}{P_0^g(u_l)}_{u_l}\mathrm{\Phi }^g(u_l;\stackrel{}{u}).$$
(5.112)
Observing that
$$_{u_l}\mathrm{\Phi }^g(u_l;\stackrel{}{u})=_r\mathrm{\Phi }^g(r;\stackrel{}{u})|_{r=u_l}+_{u_l}\mathrm{\Phi }^g(r;\stackrel{}{u})|_{r=u_l},_r\mathrm{\Phi }^g(r;\stackrel{}{u})|_{r=u_l}=2_{u_l}\mathrm{\Phi }^g(r;\stackrel{}{u})|_{r=u_l},$$
we obtain the relation $`_{u_l}\mathrm{\Phi }^g(u_l;\stackrel{}{u})={\displaystyle \frac{3}{2}}_r\mathrm{\Phi }^g(r;\stackrel{}{u})|_{r=u_l}`$. Therefore
$$\frac{^2}{u_l^2}(\lambda _l(\stackrel{}{u})t+w_l(\stackrel{}{u}))=\frac{}{u_l}\left(\frac{_{m=1,ml}^{2g+1}(u_lu_m)}{P_0^g(u_l)}\right)\mathrm{\Phi }^g(u_l;\stackrel{}{u})+\frac{3}{2}\frac{_{\stackrel{m=1}{ml}}^{2g+1}(u_lu_m)}{P_0^g(u_l)}_r\mathrm{\Phi }^g(r;\stackrel{}{u})|_{r=u_l}.$$
(5.113)
From (5.111-5.113) and the fact that $`{\displaystyle \frac{_{m=1,ml}^{2g+1}(u_lu_m)}{P_0(u_l)}}0`$ for $`u_1>u_2>\mathrm{}>u_{2g+1}`$, it is clear that system (5.103) is equivalent to system (5.30). $`\mathrm{}`$ We remark that to avoid higher order degeneracies in system (5.30), we impose the condition
$$(_r)^2\mathrm{\Phi }^g(r;\stackrel{}{u})_{r=u_l(x_c,t_c)}0.$$
(5.114)
Indeed we prove that the above condition guarantees that the genus of solution of the Whitham equations increases at most by one in the neighborhood of the point of gradient catastrophe. It is sufficient to show that the transition to genus $`g+2`$ does not occur. For the purpose let us suppose $`l`$ even and let us consider the function
$$0(_r)^2\mathrm{\Phi }^g(r;\stackrel{}{u})_{r=u_l(x_c,t_c)}=\text{const}\times \mathrm{\Phi }^{g+2}(u_l;u_1,\mathrm{},u_{l1},u_l,u_l,u_l,u_l,u_l,u_{l+1},\mathrm{},u_{2g+1}).$$
(5.115)
By proposition 5.3 in order to have a genus $`g+2`$ solution in the neighborhood of the point of gradient catastrophe, the function
$$\mathrm{\Phi }^{g+2}(r;u_1,\mathrm{},u_{l1},u_l+ϵ_1,u_l+ϵ_2,u_l+ϵ_3,u_l+ϵ_4,u_l+ϵ_5,u_{l+1},\mathrm{},u_{2g+1})$$
has to change sign in each of the intervals $`(u_l+ϵ_2,u_l+ϵ_3)`$ and $`(u_l+ϵ_4,u_l+ϵ_5)`$ ($`l`$ even) for arbitrary small $`ϵ_1>ϵ_2>ϵ_3>ϵ_4>ϵ_5>0`$. Because of (5.115)
$$\mathrm{\Phi }^{g+2}(r;u_1,\mathrm{},u_{l1},u_l+ϵ_1,u_l+ϵ_2,u_l+ϵ_3,u_l+ϵ_4,u_l+ϵ_5,u_{l+1},\mathrm{},u_{2g+1})0$$
for sufficiently small $`ϵ_1>ϵ_2>ϵ_3>ϵ_4>ϵ_5>0`$ and for $`r(u_l+ϵ_2,u_l+ϵ_5)`$. Therefore the transition to genus $`g+2`$ does not occur. We can exclude transitions from genus $`g`$ to genus $`g+n`$, $`n>2`$, by perturbations arguments. Therefore it is legitimate to consider the point of gradient catastrophe that solve (5.30) and satisfies (5.114) as a point of the boundary between the domains $`D_g`$ and $`D_{g+1}`$. We remark that in the genus $`g=0`$ case, the condition (5.115) is not essential as illustrated by theorem (3.2)
Proof of Theorem (5.9)
A phase transition may occur between the zero-phase solution and the two-phase solution. As shown on Figure 5.5 we can have a double leading edge, a trailing-leading edge and a double trailing edge. There are also other types of boundary between the zero-phase solution and the two-phase solution which we call “point of gradient catastrophe $`\&`$ leading edge” and “point of gradient catastrophe $`\&`$ trailing edge”.
Multiple transitions may also occur between the $`g`$-phase solution and the $`(g+2)`$-phase solution. In order to determine the systems which describe such phase transitions we consider the Riemann surface $`𝒮_{g+2}`$ of genus $`g+2`$ defined by
$$\stackrel{~}{\mu }^2=(rv_1\sqrt{\delta }_1)(rv_1+\sqrt{\delta }_1)(rv_2\sqrt{\delta }_2)(rv_2+\sqrt{\delta }_2)\mu ^2,v\mathrm{I}\mathrm{R},$$
$$\mu ^2(r)=\underset{j=1}{\overset{2g+1}{}}(ru_j),u_1>u_2>\mathrm{}>u_{2g+1},$$
where $`v_ju_i`$, $`j=1,2`$, $`i=1,\mathrm{},2g+1`$. Then we study the hodograph transform for the distinct variables $`v_1\pm \sqrt{\delta }_1`$, $`v_2\pm \sqrt{\delta }_2`$ and $`u_1,\mathrm{},u_{2g+1}`$ in the independent limits $`\delta _10`$ and $`\delta _20`$. When $`v_1`$ and $`v_2`$ belong to the bands (5.15), we are considering the double leading edge. Repeating the calculations done for the single leading edge we can determine the equations which describe the phase transition for the double leading edge, namely
$`\{\begin{array}{ccc}_{v_1}\mathrm{\Phi }^g(v_1;\stackrel{}{u})=0\hfill & & \\ \mathrm{\Phi }^g(v_1;\stackrel{}{u})=0\hfill & & \\ _{v_2}\mathrm{\Phi }^g(v_2;\stackrel{}{u})=0\hfill & & \\ \mathrm{\Phi }^g(v_2;\stackrel{}{u})=0\hfill & & \\ x=\left[12t{\displaystyle \frac{P_1^g(r)}{P_0^g(r)}}+{\displaystyle \frac{R^g(r)}{P_0^g(r)}}\right]_{r=u_i}\hfill & & ,i=1,\mathrm{},2g+1,g>0\hfill \end{array}`$ (5.121)
From (5.20-5.26) analogous systems can be obtained for the trailing-leading edge and double trailing edge.
For studying the phase transitions between the $`g`$-phase solution and the $`(g+n)`$-phase solution, $`n1`$, having $`n_1`$ leading edges and $`n_2`$ trailing edges, $`n_1+n_2=n`$ we consider the Riemann surface
$$\stackrel{~}{\mu }^2=\underset{j=1}{\overset{n_1}{}}(rv_j\sqrt{\delta }_j)(rv_j+\sqrt{\delta }_j)\underset{k=1}{\overset{n_2}{}}(ry_k\sqrt{\rho }_k)(ry_k+\sqrt{\rho }_k)\mu ^2,$$
$$\mu ^2(r)=\underset{j=1}{\overset{2g+1}{}}(ru_j),u_1>u_2>\mathrm{}>u_{2g+1}.$$
Here $`v_j`$, $`j=1,\mathrm{},n_1`$, belongs to the bands (5.15) and $`y_k`$, $`k=1\mathrm{},n_2`$, belongs to the gaps (5.16), $`0<\delta _j1`$, $`j=1,\mathrm{},n_1`$, and $`0<\rho _k1`$, $`k=1\mathrm{},n_2`$.
We study the hodograph transform (3.6) for the variables $`v_j\pm \sqrt{\delta }_j`$, $`j=1\mathrm{},n_1`$, $`y_k\pm \sqrt{\rho }_k`$, $`k=1\mathrm{},n_2`$ and $`u_1,\mathrm{},u_{2g+1}`$ in the independent limits $`\delta _j0`$, $`j=1\mathrm{},n_1`$ and $`\rho _k0`$ $`k=1\mathrm{},n_2`$. Repeating the calculations done for proving Theorem 5.4 and Theorem 5.6 it is easy to show that system (5.45) describes the phase transition between the $`g`$-phase solution and the $`(g+n)`$-phase solution having $`n_1`$ leading edges, $`n_2`$ trailing edges and no points of gradient catastrophe, $`n_1+n_2=n>1`$. If we suppose that the $`g`$-phase solution has also $`n_3`$ points of gradient catastrophe, $`n_1+n_2+n_3=n`$, then combining Theorem 5.4, Theorem 5.6 and Theorem 5.7 we obtain proposition 5.12. $`\mathrm{}`$
## 6 Conclusion
In this work we have constructed, in implicit form, the $`g`$-phase solution of the Whitham equations for monotone increasing smooth initial data $`x=f(u)|_{t=0}`$. The goal is obtained solving the Tsarev system (3.4). We have shown that the solution of the Tsarev system which satisfies the natural boundary conditions (3.7-3.14) is unique. Then we have investigated the conditions for the solvability of the hodograph transform (3.6). For the purpose we have derived all the equations which describe a phase transition of the solution of the Whitham equations. Studying when phase transitions occur we have been able to prove the second main result of this work. Namely we have shown that when the initial data satisfies (5.1), the solution of the Whitham equations has genus $`gN`$ for all $`x`$ and $`t0`$. It is still an open problem to effectively determine, on the $`xt0`$ plane, the function $`0g(x,t)N`$ from the generic initial data $`x=f(u)|_{t=0}`$.
Acknowledgments. I am indebted with Professor Boris Dubrovin who posed me the problem of this work and gave me many hints to reach the solution. I am grateful to Professor Sergei Novikov for his suggestions during the preparation of the manuscript. This work was partially support by a CNR grant 203.01.70 and by a grant of S. Novikov.
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# 1 Introduction
## 1 Introduction
In the previous paper<sup>?</sup> we studied quantum geometrodynamics (QGD) of the Bianchi-IX model in the framework of extended phase space (EPS) approach elaborated by Batalin, Fradkin and Vilkovisky (BFV)<sup>?</sup><sup>-</sup><sup>?</sup>. It is generally accepted that in the BFV approach one singles out a BRST-invariant sector which is supposed to coincide with a gauge-invariant one. This scheme was realized in the works<sup>?</sup><sup>-</sup><sup>?</sup>; it has a manifest physical interpretation in the case of S-matrix theory. However, in quantum cosmology appropriate mathematical operations are just formal and the question arises, whether these operations are mathematically and physically justified. This question has been explored in our paper<sup>?</sup>. Our main result consisted in the demonstration that in a closed universe without asymptotical states BRST-invariance is not equivalent to gauge invariance. A mathematical indication to the nonequivalence of gauge and BRST invariance is given by the well-known parametrization noninvariance of the Wheeler – DeWitt (WDW) equation. As we have shown in the framework of the BFV approach, the latter is an ill-hidden gauge noninvariance. For this reason singling out a BRST-invariant sector is not motivated in QGD and is not used in our approach. This circumstance distinguished our version of QGD from other works on quantum geometrodynamics. A physical ground for our approach consists in the impossibility to remove an observer from a closed universe, and, as a consequence, the necessity to take into account his affecting physical processes. The lack of gauge symmetry, or, more precisely, breaking down this symmetry when considering observation means (OM) as a part of an integrated system in our version of QGD is described by eigenvalues of gravitational super-Hamiltonian; we refer to the latters as energy levels of gravitation vacuum condensate (GVC).
It is worth emphasizing that in a gauge-invariant approach to QGD it is assumed that the information about the reduction of a wave packet in the process of evolution of the Universe is contained in boundary conditions<sup>?</sup><sup>-</sup><sup>?</sup>. In contrast to the conventional approach, our version of QGD admits the Copenhagen interpretation and aims at describing the integrated system “a physical object (gravitational field) + observation means (gravitational vacuum condensate)”.
Approaches of quantum cosmology are traditionally tested and developed for the Bianchi-IX model and its particular case – the isotropic model<sup>?</sup><sup>-</sup><sup>?</sup>. As it is known, the Bianchi-IX model can be represented as an isotropic space where two transversal nonlinear gravitational waves are excited. So, in papers based on the conventional Wheeler – DeWitt QGD (or in those where the BRST-invariant sector in EPS is singled out) observables include 3-space volume and amplitudes of gravitational waves. The convential approach is developing applying new consepts such as supergavity and superstrings<sup>?</sup><sup>-</sup><sup>?</sup>. However, the interpretation of the Wheeler – De Witt theory is not changing: the information about the past of the Universe as well as its future is contained in boundary conditions. In our modification of QGD a new feature appears – the gravitational vacuum condensate with its degrees of freedom. Wave functions of the integrated system are represented by normalized wave packets; we suppose that their structure is determined while the Universe creating from “Nothing”. The influence of observation means on the evolution of the integrated system in the proposed version of QGD is described by transformation of the form of a wave packet.
The wave functions satisfy the dynamical Schrödinger equation (SE). For the Bianchi-IX model the normalized in the EPS general solution (GS) to the gauge-noninvariant SE reads
$$\mathrm{\Psi }(Q^a,Q^0,\theta ,\overline{\theta };t)=\mathrm{\Psi }_k(Q^a)\mathrm{exp}(iE_kt)(\overline{\theta }+i\theta )\delta (Q^0f(Q^a)k)𝑑E_k𝑑k,$$
(1)
where $`\mathrm{\Psi }_k(Q^a)`$ is a solution to the stationary equation
$$H_k^0\mathrm{\Psi }_k(Q^a)=E_k\mathrm{\Psi }_k(Q^a),$$
(2)
$`\theta ,\overline{\theta }`$ are the ghosts, $`Q^0`$ and $`Q^a`$ are the gauge and physical variables respectively, specification of which see below in Sec. 2.
The norm integral for the WF (1) over the full set of variables $`Q,\theta ,\overline{\theta }`$ results in
$$\mathrm{\Psi }_k^{}(Q^a,t)\mathrm{\Psi }_k(Q^a,t)M_k(Q^a)𝑑k\underset{a}{}dQ^a,$$
so the GS to the SE, under the condition the $`\mathrm{\Psi }_k(Q^a,t)`$ to be a sufficiently narrow packet over $`k`$, is normalizable with respect to the gauge variable, as well as to the ghosts and the physical variables.
The theory does not control the WF dependence on the parameter $`k`$ which in the classical dynamics determines an initial clock setting, and this additional degree of freedom has to be referred to an observer.
Investigation of the particular BRST-invariant solution has shown that it really satisfies the WDW equation but cannot be normalized. On the other hand, as we have already mentioned above, the well known parametrization noninvariance of the WDW equations is shown to be the ill-hidden gauge noninvariance. The latter means that the WDW theory does not achieve its object to give a gauge-invariant description of the Universe.
By this reason and because of the loss of probability interpretation we do not see any ground to be guided by the WDW theory, and explore the possibilities of the gauge-noninvariant QGD version.
The factored part of the GS (1) – $`\delta `$-function and ghosts – represents the OM described by the gauge-fixing term in the Lagrangian and by the appropriate “energy-momentum tensor” (quasi-EMT). The latter corresponds to the continuous medium that fills the whole space and can be called “gravitational vacuum condensate”; it is quantitatively fixed by an eigenvalue $`E_k`$ of the Hamiltonian $`H_k^0`$. The structure of the GS shows a GVC to be an important factor of the global evolution. The purpose of this paper is to investigate this effect on the base of the exact solution to the SE for the simplified model with the frozen degree of freedom $`Q^3`$, in a simple gauge.
In Sec. 2. all the necessary notations and the general equations for the Bianchi-IX model are given; in Sec. 3. the exact conditionally-classical solution for the simple case is considered. We shall show that GVC can play a decisive role in forming non-trivial cosmological scenarios. In Sec. 4. the exact solution to the SE is considered, we analyse the formation and properties of the wave packet representing the evolving universe. In Sec. 5. the transition to the semiclassical solution is considered and it is shown how the quantum cosmological effects become negligible in the semiclassical region; here various cosmological scenarios depending on the GVC and the homogeneous scalar field are considered. In Sec. 6. the concept of time in the gauge-noninvariant QGD is discussed.
## 2 The General Model Equations
The Bianchi-IX 4-interval is given by
$$ds^2=N^2(t)dt^2\eta _{ab}(t)e_i^ae_k^bdx^idx^k;$$
(3)
$$\eta _{ab}(t)=\mathrm{diag}(a^2(t),b^2(t),c^2(t)),$$
(4)
$$\begin{array}{c}e_i^1=(\mathrm{sin}x^3,\mathrm{cos}x^3\mathrm{sin}x^1,0),\hfill \\ e_i^2=\text{}(\mathrm{cos}x^3,\mathrm{sin}x^3\mathrm{sin}x^1,0),\hfill \\ e_i^3=\text{}(0,\mathrm{cos}x^1,1).\hfill \end{array}$$
$$a=\frac{1}{2}\mathrm{exp}\left[\frac{1}{2}\left(Q^1+Q^2+\sqrt{3}Q^3\right)\right];b=\frac{1}{2}\mathrm{exp}\left[\frac{1}{2}\left(Q^1+Q^2\sqrt{3}Q^3\right)\right];$$
$$c=\frac{1}{2}\mathrm{exp}\left(\frac{1}{2}Q^1Q^3\right).$$
The model is assumed to include an arbitrary number $`K`$ of real homogeneous scalar fields $`Q^4,\mathrm{},Q^{K+3}`$ with some potential $`U_s`$. The gauge coordinate $`Q^0`$ is defined by an arbitrary parametrization function
$$\zeta (Q^0,Q^1,\mathrm{},Q^{K+3})=\mathrm{ln}\left[\frac{1}{N}\mathrm{exp}\left(\frac{3}{2}Q^1\right)\right],$$
(5)
and gauges are supposed to be not depending on time,
$$Q^0=f(Q^1,\mathrm{})+k,k=\mathrm{const}.$$
(6)
The SE reads
$$i\frac{\mathrm{\Psi }(Q^0,Q^1,\mathrm{},\theta ,\overline{\theta };t)}{t}=H\mathrm{\Psi }(Q^0,Q^1,\mathrm{},\theta ,\overline{\theta };t),$$
(7)
where $`\mathrm{\Psi }(Q^\alpha ,\theta ,\overline{\theta };t)`$ is a universe WF $`(\alpha =0,a;a=1,\mathrm{},K+3),\theta ,\overline{\theta }`$, are the ghosts,
$$H=i\zeta ,_0\frac{}{\theta }\frac{}{\overline{\theta }}\frac{1}{2M}\frac{}{Q^\alpha }G^{\alpha \beta }\frac{}{Q^\beta }+\mathrm{e}^\zeta (UV),$$
(8)
$`\zeta ,_0=\zeta (Q^\alpha )/Q^0`$, probability measure
$$M=\mathrm{const}\zeta ,_0\mathrm{exp}\left(\frac{K+3}{2}\zeta \right),$$
(9)
$$U(Q)=\mathrm{e}^{2Q^1}U_g(Q^2,Q^3)+\mathrm{e}^{3Q^1}U_s(Q^4,\mathrm{}),$$
(10)
$`U_g(Q^2,Q^3)`$ $`=`$ $`{\displaystyle \frac{2}{3}}\{\mathrm{exp}\left[2(Q^2+\sqrt{3}Q^3)\right]+\mathrm{exp}\left[2(Q^2\sqrt{3}Q^3)\right]+\mathrm{exp}(4\chi )`$
$``$ $`2\mathrm{exp}[(Q^2+\sqrt{3}Q^3)]2\mathrm{exp}(Q^2+\sqrt{3}Q^3)2\mathrm{exp}(2Q^2)\}`$
$`V=`$ $``$ $`{\displaystyle \frac{3}{12}}{\displaystyle \frac{(\zeta ,_0)^a(\zeta ,_0)_a}{\zeta ,_0^2}}+{\displaystyle \frac{(\zeta ,_0)_a^a}{3\zeta ,_0}}+{\displaystyle \frac{K+1}{6\zeta ,_0}}\zeta _a(\zeta ,_0)^a`$ (11)
$`+`$ $`{\displaystyle \frac{1}{24}}(K^2+3K+14)\zeta _a\zeta ^a+{\displaystyle \frac{K+2}{6}}\zeta _a^a,`$
$`\zeta _a=\zeta /Q^a+f,_a\zeta /Q^0`$,
$$G^{\alpha \beta }=M\mathrm{e}^\zeta \left(\begin{array}{cc}f,_af^{,a}& f^{,a}\\ f^{,a}& \gamma ^{ab}\end{array}\right).$$
To make the analysis more visual, we begin with considering the conditionally-classical solution to the Einstein equations supplemented with the quasi-EMT of the ghosts and the GVC. In common case this set of equations, yielded by the “classical version” of the Hamiltonian (8), reads
$`\left(\mathrm{e}^\zeta \dot{Q}_a\right)^.`$ $`+`$ $`\mathrm{e}^\zeta U,_a\dot{\lambda }f,_a{\displaystyle \frac{1}{2}}\zeta ,_a\mathrm{e}^\zeta \dot{Q}^b\dot{Q}_b\zeta ,_a\mathrm{e}^\zeta U+{\displaystyle \frac{i\zeta ,_0,_a}{\zeta ,_0^2}}\dot{\overline{\theta }}\dot{\theta }=0,`$ (12)
$`{\displaystyle \frac{1}{2}}\zeta ,_0\mathrm{e}^\zeta \dot{Q}^a\dot{Q}_a`$ $`+`$ $`\zeta ,_0\mathrm{e}^\zeta U{\displaystyle \frac{i\zeta ,_0,_0}{\zeta ,_0^2}}\dot{\overline{\theta }}\dot{\theta }\dot{\lambda }=0,`$ (13)
$`\dot{Q}^0f,_a\dot{Q}^a`$ $`=`$ $`0,`$ (14)
$`\left(\zeta ,_0^1\dot{\theta }\right)^.`$ $`=`$ $`0,`$ (15)
$`\left(\zeta ,_0^1\dot{\overline{\theta }}\right)^.`$ $`=`$ $`0,`$ (16)
where $`\lambda =\pi +\dot{\overline{\theta }}\theta ,\pi `$ is the Lagrange multiple fixing the gauge (14), coupled to the integral of motion describing the GVC:
$$(\zeta ,_0)_k^1\dot{\lambda }=E_k.$$
## 3 The Conditionally-Classical Exact Solution
Taking the parametrization and the gauge
$$\zeta =Q^0=k$$
(17)
($`k=\mathrm{const}`$), one can obtain an exact particular solution to Eqs. (12) – (16) with $`Q^3=0`$. The existence of this solution gives the formal opportunity to consider the model without this degree of freedom.
Under the condition (17) the ghost variables vanish from Eqs. (12) – (14), and the latter form a closed set concerning the physical variables; the state equation of the GVC becomes extremely hard,
$$p=\epsilon =\frac{\dot{\lambda }}{2\pi ^2}\mathrm{exp}(k3Q^1),$$
(18)
$$\dot{\lambda }=E.$$
(19)
Note that conditionality of the classical approach shows here in the ghost presence in the integral of motion (19)
$$\dot{\lambda }=\dot{\pi }\dot{\overline{\theta }}\dot{\theta },$$
i.e. the forms made of the Grassmannian variables appear as parameters of the theory.
Let us turn to the case of a single massless linear scalar field $`\varphi =Q^4`$ and put
$$U_s(\varphi )=0,$$
in (10). Now we have the simple equation for $`\varphi `$
$$\ddot{\varphi }=0,$$
thus,
$$\dot{\varphi }=C_s=\mathrm{const}.$$
(20)
The scalar field behaves as a medium with positive energy density and with an extremely hard equation of state
$$p_{(scal)}=\epsilon _{(scal)}\mathrm{exp}(3Q^1)\dot{\varphi }^2=C_s^2\mathrm{exp}(3Q^1)$$
like that of the GVC (18).
As one can see, in the present model the Universe is filled with the two-component medium described by the parameters $`E`$ and $`C_s`$. Below we will show that relation between the two parameters essentially affects cosmological evolution at the quantum stage of the Universe existence as well as at the semiclassical one. Here is the difference between our consideration and the usual investigation of the Bianchi-IX model in general relativity.
The equations for $`Q^1,Q^2`$ take the form:
$$\ddot{Q}^1\frac{4}{3}\left[\mathrm{exp}(2Q^14Q^2)4\mathrm{exp}(2Q^1Q^2)\right]=0,$$
$$\ddot{Q^2}\frac{4}{3}[2\mathrm{exp}(2Q^14Q^2)2\mathrm{exp}(2Q^1Q^2)]=0.$$
Integration is simplified with the substitution
$$z_1=2Q^14Q^2,z_2=2Q^1Q^2;$$
(21)
after replacing
$$t\mathrm{e}^kt$$
the solution is available in the form
$$\mathrm{exp}\left(\frac{z_2}{2}\right)=\frac{\alpha }{\mathrm{cosh}[2\alpha (tt_0)]},\mathrm{exp}\left(\frac{z_1}{2}\right)=\frac{\beta }{\mathrm{cosh}[2\beta (tt_1)]},$$
(22)
where $`\alpha ,\beta ,t_0,t_1`$ are the integration constants. Without loss of generality, by shifting zero time one can put $`t_1=0`$. For the metric (3) – (4) one finds:
$$a^2=b^2=\frac{1}{4}\mathrm{exp}\left(z_2\frac{1}{2}z_1\right)=\frac{a^2\mathrm{cosh}(2\beta t)}{4\beta \mathrm{cosh}^2[2\alpha (tt_0)]};$$
$$c^2=\frac{1}{4}\mathrm{exp}\left(\frac{1}{2}z_1\right)=\frac{\beta }{4\mathrm{cosh}(2\beta t)}.$$
From the constraint equation (13) with $`Q^3=0`$, and (20) it follows:
$$\frac{1}{24}\left(\dot{z}_1^24\dot{z}_2^2\right)+\frac{2}{3}\left[\mathrm{exp}(z_1)4\mathrm{exp}(z_2)\right]=E_k\frac{1}{2}C_s^2,$$
where $`E_k=\mathrm{e}^kE`$; hence, in turn,
$$\alpha ^2\frac{1}{4}\beta ^2=\frac{3}{8}\left(E_k\frac{1}{2}C_s^2\right).$$
(23)
So, the dynamics of the model depends qualitatively on a relation between $`C_s`$ and $`E_k`$. Various cosmological effects of this dependence will be discussed in Sec. 5., after considering the quantum version of the exact solution.
## 4 The Exact Solution to the Schrödinger Equation
The task of constructing WF (1) is reduced to searching for a solution to stationary Eq. (2) for the physical part of the WF under the parametrization-and-gauge condition (17). This equation reads
$$\frac{1}{2}\frac{^2\mathrm{\Psi }_k}{Q_aQ^a}+U(Q^a)\mathrm{\Psi }_k(Q^a)=E_k\mathrm{\Psi }_k(Q^a),$$
(24)
$`a=(1,2,4)`$. Substitution of (21) enables to separate the variables in the equation, and it can be written in the following manner
$$\left(6\widehat{L}_1\frac{3}{2}\widehat{L}_2+\frac{1}{2}\widehat{L}_3kE\right)\mathrm{\Psi }_k(z_1,z_2,\varphi )=0,$$
where
$`\widehat{L}_1`$ $`=`$ $`{\displaystyle \frac{^2}{z_1^2}}+{\displaystyle \frac{1}{9}}\mathrm{exp}(z_1),`$
$`\widehat{L}_2`$ $`=`$ $`{\displaystyle \frac{^2}{z_2^2}}+{\displaystyle \frac{16}{9}}\mathrm{exp}(z_2),`$
$`\widehat{L}_3`$ $`=`$ $`{\displaystyle \frac{^2}{\varphi ^2}}.`$
The eigenfunctions of the operators $`\widehat{L}_1,\widehat{L}_2`$ appropriate to the positive eigenvalues $`\nu _1^2/4,\nu _2^2/4`$ are the modified Bessel functions with an imaginary index,
$$\psi _\nu (z)=\frac{1}{\sqrt{2\pi \mathrm{\Gamma }(i\nu )}}K_{i\nu }\left[A\mathrm{exp}\left(\frac{z}{2}\right)\right];(A_1;A_2)=(\frac{2}{3};\frac{8}{3}).$$
So far as
$$\mathrm{exp}\left(\frac{z_1}{2}\right)=4c^2;\mathrm{exp}\left(\frac{z_2}{2}\right)=4ac,$$
the quantum number $`\nu _1`$ determines probability distribution for the scale $`c`$, and so does the quantum number $`\nu _2`$ for the scale $`a=b`$ at a given $`c`$ value. Note, that for any $`\nu `$ there exists a semiclassical solution to the problem,
$`\psi _\nu (z)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}\mathrm{\Gamma }(i\nu )\mathrm{exp}\left({\displaystyle \frac{\pi \nu }{2}}\right)\left({\displaystyle \frac{\nu ^2}{4}}{\displaystyle \frac{A^2}{4}}\mathrm{exp}(z)\right)^{\frac{1}{4}}}}\mathrm{cos}[2(\sqrt{{\displaystyle \frac{\nu ^2}{4}}{\displaystyle \frac{A^2}{4}}\mathrm{exp}(z)}`$
$``$ $`{\displaystyle \frac{\nu }{2}}Artanh\sqrt{1{\displaystyle \frac{A^2}{\nu ^2}}\mathrm{exp}(z)})+{\displaystyle \frac{\pi }{4}}],`$
$`z<z_\nu `$;
$`\psi _\nu (z)`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}\mathrm{\Gamma }(i\nu )\mathrm{exp}\left({\displaystyle \frac{\pi \nu }{2}}\right)\left({\displaystyle \frac{A^2}{4}}\mathrm{exp}(z){\displaystyle \frac{\nu ^2}{4}}\right)^{\frac{1}{4}}}}\mathrm{exp}[2(\sqrt{{\displaystyle \frac{A^2}{4}}\mathrm{exp}(z){\displaystyle \frac{\nu ^2}{4}}}`$
$``$ $`{\displaystyle \frac{\nu }{2}}\mathrm{arctan}\sqrt{{\displaystyle \frac{A^2}{\nu ^2}}\mathrm{exp}(z)1})],`$
$`z>z_\nu ,z_\nu =\mathrm{ln}(\nu ^2/A^2)`$ being the classical turning-point.
The operator $`\widehat{L}_3`$ eigenfunctions are the plane waves
$$\psi _\varrho (\varphi )=\frac{1}{\sqrt{2\pi }}\mathrm{exp}(i\varrho \varphi )$$
that is in agreement with the classical solution (20).
The GS to Eq. (24) for a given value of the parameter $`E_k`$, describing a GVC state, is a superposition
$`\mathrm{\Psi }_{E_k}(z_1,z_2,\varphi )`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\varrho 𝑑\nu _1𝑑\nu _2c_1(\nu _1,\nu _2,\varrho )\psi _{\nu _1}(z_1)\psi _{\nu _2}(z_2)\psi _\varrho (\varphi )}`$ (25)
$`\times `$ $`\delta \left({\displaystyle \frac{3}{2}}\nu _1^2{\displaystyle \frac{3}{8}}\nu _2^2+{\displaystyle \frac{1}{2}}\varrho ^2E_k\right).`$
However, the stationary states, the wave functions (25) correspond to, are not physical, being unnormalizable because of continuity of the $`E_k`$ value spectrum. A time-dependent wave packet fits a physical state:
$$\mathrm{\Psi }_k(z_1,z_2,\varphi ,t)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑E_kc_2(E_k)\mathrm{\Psi }_{E_k}(z_1,z_2,\varphi )\mathrm{exp}\left[iE_k(tt_0)\right].$$
(26)
Note, that in the expressions (25), (26) the quantity $`E_k`$ appears as a controlling parameter providing, via the $`\delta `$-function, correlation of the quantum numbers $`\nu _1,\nu _2,\varrho `$, and, by that, a probability distribution of the space scales at the quantum stage of the Universe evolution.
## 5 The Semiclassical Regime and Cosmological Effects
One can obtain the classical evolution law by computing the operators $`\mathrm{exp}(z_1/2)`$, $`\mathrm{exp}(z_2/2)`$ mean values over the packet (26). The matrix elements will be required for that,
$`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z\mathrm{exp}\left({\displaystyle \frac{z}{2}}\right)\psi _\mu ^{}(z)\psi _\nu (z)`$ (27)
$`=`$ $`\left[2\pi \mathrm{\Gamma }(i\mu )\mathrm{\Gamma }(i\nu )\right]^1{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z\mathrm{exp}\left({\displaystyle \frac{z}{2}}\right)K_{i\mu }\left[A\mathrm{exp}\left({\displaystyle \frac{z}{2}}\right)\right]K_{i\nu }\left[A\mathrm{exp}\left({\displaystyle \frac{z}{2}}\right)\right]`$
$`=`$ $`\pi \left[4A\mathrm{\Gamma }(i\mu )\mathrm{\Gamma }(i\nu )\right]^1\left\{\mathrm{cosh}\left[{\displaystyle \frac{\pi }{2}}(\mu +\nu )\right]\mathrm{cosh}\left[{\displaystyle \frac{\pi }{2}}(\mu \nu )\right]\right\}^1.`$
To be able to describe really the classically evolving Universe, the packet (25) – (26) should be sufficiently narrow, i.e. $`c_1(\nu _1,\nu _2,\varrho )`$ and $`c_2(E_k)`$ should not deviate from zero values beyond a small vicinity of their arguments near $`(\overline{\nu }_1,\overline{\nu }_2,\overline{\varrho })`$ and $`\overline{E}_k`$. Therefore,
$$\mu +\nu 2\overline{\nu };\mu \nu \frac{A\omega }{2\overline{\nu }},$$
(28)
where $`\omega =A^1(\mu ^2\nu ^2)`$ is the difference between the two values of the parameter $`E_k`$ corresponding to the quantum numbers $`\mu `$ and $`\nu `$. Note that the matrix element (27) depends weakly on $`\nu `$ and is sharply decreasing when $`|\mu \nu |`$ increasing. And so, making use of the approximations (28) one obtaines, for the average exponents $`\mathrm{exp}(z/2)`$,
$`\overline{\mathrm{exp}\left({\displaystyle \frac{z}{2}}\right)}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{tanh}(\pi \overline{\nu }){\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑\omega {\displaystyle \frac{\mathrm{exp}\left[i\omega (tt_0)\right]}{\mathrm{cosh}\left({\displaystyle \frac{\pi A\omega }{4\overline{\nu }}}\right)}}`$ (29)
$`=`$ $`{\displaystyle \frac{\overline{\nu }\mathrm{tanh}(\pi \overline{\nu })}{A\mathrm{cosh}[(2A^1\overline{\nu }(tt_0)]}}.`$
In the classical limit $`\overline{\nu }`$ is large, hence $`\mathrm{tanh}(\pi \overline{\nu })1`$, and, comparing (29) with the classical expression (22), one concludes that
$$\alpha =\frac{\overline{\nu }_2}{A_2}=\frac{3}{8}\overline{\nu }_2;\beta =\frac{\overline{\nu }_1}{A_1}=\frac{3}{2}\overline{\nu }_1.$$
From the (25), the equation for the mean values follows $`(k=0)`$:
$$\frac{3}{2}\overline{\nu }_1^2\frac{3}{8}\overline{\nu }_2^2+\frac{1}{2}\overline{\varrho }^2=\overline{E},$$
(30)
because $`\overline{\nu _1^2}\overline{\nu }_1^2`$ an so on. Comparison of (23) and (30) gives
$$\overline{\varrho }=C_s.$$
These results reveal the following possible cosmological scenarios for the Bianchi-IX universe at the semiclassical stage of its evolution according to a relation between the parameters $`\overline{E}`$ and $`C_s`$ of the two condensates.
1. Empty space ($`C_s=0`$, $`\overline{E}_k=0`$); $`\alpha =\frac{\beta }{2}`$.
In the limit $`t=\pm \mathrm{}`$
$$a^2=b^2=\frac{\beta }{8};c^2=0,$$
i.e. the metric (3) asymptotically takes the form
$$ds^2=(\beta c^2)dt^2\frac{\beta }{8}(d\vartheta ^2+\mathrm{sin}^2\vartheta d\phi ^2),$$
$$\vartheta =x^1,\phi =x^2.$$
When reaching singularity in one of the dimensions, two others form a stationary space of constant curvature. Here one deals with a regime of dynamical compactification, a space with simple topology being compactified.
2. Space is filled with the medium having positive energy density: $`\overline{E}_k<\frac{1}{2}C_s^2;\alpha >\frac{1}{2}\beta `$.
For $`\alpha =\beta ,t_0=0`$ the model is isotropic.
For $`\alpha =\beta ,t_00`$ the model is unisotropic, but the singularity has an isotropic nature.
For $`\frac{1}{2}\beta <\alpha <\beta `$ in the pre-singular state $`a^2=b^2c^2`$, i.e. (2 + 1)-dimensional space-time arises, where the 2-space has constant curvature.
For $`\alpha >\beta `$ in the pre-singular state $`c^2a^2=b^2`$, however, the model is not reduced to the space of less dimensions.
In all the cases for $`\overline{E}_k<\frac{1}{2}C_s^2`$ space at singularity is contracted to a point.
3. Space is filled with the medium having negative energy density: $`\overline{E}_k>\frac{1}{2}C_s^2;\alpha <\frac{1}{2}\beta `$.
At $`t=\pm \mathrm{}`$ the third space dimension is compactified ($`c^20`$), and the remaining two-dimensional space of constant curvature is infinitely expanding. In the special case $`\alpha =\beta /4`$ the scale factor $`a=b`$ is increasing exponentially in proper time.
So, the GVC, affecting coupling between the constants $`\alpha `$ and $`\beta `$ through the controlling parameter $`\overline{E}_k`$, determines a cosmological scenario which may contain the following phenomena:
* cosmological expansion and contraction of space;
* cosmological singularity;
* compactification of space dimensions;
* asymptotically stationary space of less dimensions;
* inflation of the Universe.
One can see that even such a simplified model reveals a number of effects probable from the standpoint of the modern cosmological ideas. Introduction of the GVC to the theory enlarges the number of possible cosmological scenarios, a concrete value of the parameter $`E_k`$ being formed at the quantum stage of the Universe existence.
Evidently, every classical cosmological evolution scenario must be in correspondence with some configuration of the wave packet (25) – (26). But not all the solutions to the SE describe classical universes, i.e. form sufficiently narrow packets to satisfy the conditions (28); in addition, even those wave packets, for which a transition to the semiclassical regime is possible, may prove to be unstable. Therefore in our approach the known problem of initial conditions for classical evolution is formulated as the problem of choice of the Universe quantum state. A quantum state in the Bianchi-IX model is determined by a concrete kind of the function
$$\stackrel{~}{C}(\nu _1,\nu _2,\varrho ,E_k)=c_2(E_k)c_1(\nu _1,\nu _2,\varrho ),$$
describing a wave packet structure.
We do not know the way the choice of the quantum state is being made by. Perhaps, it is realized according to statistical laws in the process of the Universe creation from “Nothing”. In the next paper we intend to discuss the hypothesis according to which the act of the Universe creation is understood as, occurring out of time, quantum transition from the special singular state to one of the Universe physical states which the wave packets (25) – (26) correspond to.
## 6 Conception of Time
In the previous section the transition to the semiclassical regime was considered in the GVC time. As a matter of fact, there exist four different concepts of time in the ordinary quantum theory:
* the Heisenberg time – the time parameter in Heisenberg representation; it is introduced equally with space coordinates by the device preparing an initial state of an object (preparator);
* the Schrödinger time introduced by the registrating device(registrator);
* the time in the interaction representation defined by preparators and registrators in the asymptotical regions of space-time;
* the world time defined on the object itself in the semiclassical regime (semiclassical proper time).
In QGD the only clock carrier is a GVC, and the question about availability of the Heisenberg time is closely related to the question about the Universe creation. As to the semiclassical proper time (SPT), its existence is not predetermined by the theory for, as it was mentioned above, a semiclassical regime is not predetermined itself. But, having such a regime actually, let us see how to deal with the SPT in the present theory.
From the very beginning we should emphasize that the SPT has nothing to do with the GVC as the OM carrier. By transition to the semiclassical regime the integrity effects are weakening; as a result, the fact of the GVC existence shows only through the parameters responsible, in our case, for the wave packet satisfactory width. So, the SPT does not depend on the GVC time, and the question, whether we can co-ordinate our semiclassical clocks with the GVC ones, seems to be of considerable physical significance in view of obtaining information about the quantum stage of the Universe evolution.
A semiclassical wave packet may also be presented in the form
$`\mathrm{\Psi }_k(z_1,z_2,\varphi ,t)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑E_k𝑑\varrho 𝑑\nu _1𝑑\nu _2\stackrel{~}{C}(\nu _1,\nu _2,\varrho ,E_k)}`$
$`\times `$ $`\mathrm{exp}\left[iE_k(tt_0)+i\sigma _1(z_1)+i\sigma _2(z_2)\right]`$
$`\times `$ $`\psi _\varrho (\varphi )\delta \left({\displaystyle \frac{3}{2}}\nu _1^2{\displaystyle \frac{3}{8}}\nu _2^2+{\displaystyle \frac{1}{2}}\varrho ^2E_k\right).`$
Here the sum $`\sigma _1(z_1)+\sigma _2(z_2)`$ is the part of the classical action $`S(z_1,z_2,\varphi ,t)`$, determining its dependence on $`z_1`$ and $`z_2`$ (the scalar field is treated as essentially quantum). The functions $`\sigma (z)`$ satisfy the equations
$$\frac{\sigma }{z}=\sqrt{\frac{\overline{\nu }^2}{4}\frac{A^2}{4}\mathrm{exp}(z)}.$$
(31)
Knowing the dependence of the classical action on the variables $`z`$, one can reconstruct the evolution law (22) with the help of the standard procedure. But the two mentioned methods of going over to the classical limit are applicable owing to the explicit dependence of the GS on time. And this, in turn, is caused by the available, in the theory, indication of the concrete choice of an RS which the time $`t`$ is measured in.
In the classical limit a classical subsystem of the physical object itself can be considered as an RS. Such a subsystem cannot fill the whole space; it is admissible that it would occupy a limited region of space. So, we will refer to such an RS as to a local one. The appearance of the time $`\tau `$, introduced as a parameter along a classical path, is associated with that RS.
A derivative with respect to path length can be defined by the following way:
$$\frac{d}{dt}=u(\tau )S,$$
where $`u(\tau )`$ is an arbitrary function, $`S`$ is a tangent vector to the path;
$$=(2\sqrt{3}\frac{}{z_1},i\sqrt{3}\frac{}{z_2}).$$
On the other hand,
$$\frac{d}{d\tau }=\frac{dz_1}{d\tau }\frac{}{z_1}+\frac{dz_2}{d\tau }\frac{}{z_2},$$
whence
$$\frac{dz_1}{d\tau }=12u(\tau )\frac{d\sigma _1}{dz_1};\frac{dz_2}{d\tau }=3u(\tau )\frac{d\sigma _2}{dz_2}.$$
From (31) one obtaines
$$\frac{dz_1}{d\tau }=4u(\tau )\sqrt{\beta ^2\mathrm{exp}(z_1)};\frac{dz_2}{d\tau }=4u(\tau )\sqrt{\alpha ^2\mathrm{exp}(z_2)},$$
and, in the result of integration,
$$\mathrm{exp}\left(\frac{z_1}{2}\right)=\beta \mathrm{cosh}^1\left(2\beta \left[\stackrel{~}{u}(\tau )\tau _0^{}\right]\right);\mathrm{exp}\left(\frac{z_2}{2}\right)=\alpha \mathrm{cosh}^1\left(2\alpha \left[\stackrel{~}{u}(\tau )\tau _0\right]\right);$$
(32)
$$\stackrel{~}{u}(\tau )=u(\tau )𝑑\tau ;\tau _0,\tau _0^{}=\mathrm{const}.$$
The time $`\tau `$ of a local observer emerges irrespectively of the time $`t`$ existence, but both the times may correlate between them. To bring the expressions (22) and (32) in correspondence, it is sufficient to put $`u(\tau )=1`$. This act assumes the following operations:
* determining the rate of the Universe evolution by means of cosmological observations;
* fixing, on this data base, the GVC parameters and the state equation of it;
* reconstruction of the gauge from the state equation and so fixing the GVC time scale.
In the ordinary quantum mechanics the time variable involving in SE may fail to coincide with that appearing in Heisenberg operator equations, as well as with the world time in which semiclassical system dynamics can be described. In the quantum mechanics the hypothesis about equivalence of the mentioned time variables is used, though this is nowhere specified. In the considered example we have manifested that the time variables used for describing a quantum system evolution are different in general. In QGD times associated with different observers can be brought to agreement with each other by choice of a gauge condition.
References
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# 1 Introduction
## 1 Introduction
Following the recent understanding concerning the equivalence between non-commutative and commutative gauge theories , an immediate question to ask is how the theory behaves under the S-duality interchanging the strong and weak coupling regimes. Indeed, this question has been addressed in several recent works (See also related works ).
One expects an answer as simple as follows. The four-dimensional non-commutative theory describes the low-energy worldvolume dynamics of a D3-brane in the presence of a background of NS-NS two-form potential (Kalb-Ramond potential) $`B_2`$ but none others. Under the S-duality of the Type IIB string theory, the D3-brane is self-dual, while the NS-NS two-form potential $`B_2`$ is swapped with the R-R two-form potential $`C_2`$. Thus, after the S-duality, the dual theory appears to be the one describing the low-energy worldvolume dynamics of a D3-brane in the presence of a R-R two-form potential $`C_2`$, but none others. Since there is no NS-NS two-form potential present, noncommutative deformation via Seiberg-Witten map is not possible and the resulting theory ought to be a commutative theory. However, this is apparently not the answer one obtains . Starting with a noncommutative U(1) gauge theory with coupling constant $`g`$ and noncommutativity tensor $`\theta `$ and taking the standard duality transformation, one finds that the dual theory still remains a noncommutative U(1) gauge theory, but with coupling and noncommutativity parameters
$`g_\mathrm{D}={\displaystyle \frac{1}{g}}\mathrm{and}\theta _\mathrm{D}=g^2\stackrel{~}{\theta },`$ (1)
where $`\stackrel{~}{\theta }^{\mu \nu }=\frac{1}{2}ϵ_{}^{\mu \nu }{}_{\alpha \beta }{}^{}\theta ^{\alpha \beta }`$. Alternatively, as is done in , one may utilize the gauge invariance of $`(F+B_2)`$ to dial out the space-space noncommutativity and treat the theory as the standard gauge theory in a constant magnetic field. The S-duality would then turn this into a dual gauge theory, but now in a constant electric field, which is gauge equivalent to a theory with space-time noncommutativity. Actually, in the dual theory, it turns out to be impossible to take a field theory limit that the dual theory would be best described by a noncritical open string theory, whose tension is of the order of the noncommutativity scale.
The aim of this paper is to understand the S-duality via the following routes:
$`\begin{array}{ccc}\mathrm{D3}\mathrm{brane}& & \stackrel{~}{\mathrm{D3}}\mathrm{brane}\\ & & \\ |& & |\\ |& & |\\ & & \\ \mathrm{NCYM}& & \stackrel{~}{\mathrm{NCYM}}\end{array}`$ (8)
and, if possible, to reconcile these seemingly different results concerning the S-duality of the noncommutative gauge theory from various perspectives.
## 2 Naive S-Duality
In this section, we shall be studying S-duality of noncommutative U(1) gauge theory <sup>2</sup><sup>2</sup>2This section is based on the result of . Later, similar but independent derivation of the first half of this section has appeared in .. The latter is defined by the following action:
$`S`$ $`=`$ $`{\displaystyle \frac{1}{4g^2}}{\displaystyle d^4x\text{Tr}\left(\right)},`$ (9)
where the noncommutative field strength is defined by
$`=𝐅+𝐢\{𝐀_\mu ,𝐀_\nu \}_{},𝐅_{\mu \nu }=(_\mu 𝐀_\nu _\nu 𝐀_\mu )`$
and Tr refers to “matrix notation” for spacetime index contractions, for instance, $`\text{Tr}\left(\right)`$ means $`_{\mu }^{}{}_{}{}^{\nu }_{\nu }^{}{}_{}{}^{\mu }`$. Seiberg and Witten have found explicitly the map between $`_{\mu \nu }`$ and $`𝐅_{\mu \nu }`$ :
$`=𝐅\left[𝐅\theta 𝐅(𝐀\theta )𝐅\right]+\mathrm{}.`$ (10)
The map Eq.(10) allows to expand the action Eq.(9) in powers of dimensionless combination $`\theta 𝐅`$:
$`S={\displaystyle \frac{1}{4g^2}}{\displaystyle d^4x\left[\text{Tr}\left(\mathrm{𝐅𝐅}\right)+2\text{Tr}\left(\theta \mathrm{𝐅𝐅𝐅}\right)\frac{1}{2}\text{Tr}\left(\theta 𝐅\right)\text{Tr}\left(\mathrm{𝐅𝐅}\right)+\mathrm{}\right]}.`$ (11)
Adopting the standard rule of duality transformation, we can promote the field strength $`𝐅`$ (not $``$) into an unconstrained field by including a term $`\stackrel{~}{𝐆}𝐅`$ where $`\stackrel{~}{𝐆}_{\mu \nu }=\frac{1}{2}ϵ_{\mu \nu }^{\rho \sigma }𝐆_{\rho \sigma }`$ and $`𝐆_{\mu \nu }=_\mu 𝐁_\nu _\nu 𝐁_\mu `$. Varying with respect to $`𝐁`$ imposes the constraint $`d𝐅=0`$ which is solved by $`𝐅=d𝐀`$ and we get back the original theory. On the other hand we can now vary with respect to $`𝐅`$ which gives us
$`g^2\stackrel{~}{𝐆}=𝐅(\theta \mathrm{𝐅𝐅}+𝐅\theta 𝐅+\mathrm{𝐅𝐅}\theta )+{\displaystyle \frac{1}{4}}\text{Tr}\left(\mathrm{𝐅𝐅}\right)\theta +{\displaystyle \frac{1}{2}}\text{Tr}\left(\theta 𝐅\right)𝐅+\mathrm{}`$
which can be inverted, to the lowest order in $`\theta `$, as
$`𝐅=g^2\stackrel{~}{𝐆}+g^4\theta \stackrel{~}{𝐆}\stackrel{~}{𝐆}+g^4\stackrel{~}{𝐆}\theta \stackrel{~}{𝐆}+g^4\stackrel{~}{𝐆}\stackrel{~}{𝐆}\theta {\displaystyle \frac{g^4}{4}}\text{Tr}\left(\stackrel{~}{𝐆}\stackrel{~}{𝐆}\right)\theta {\displaystyle \frac{g^4}{2}}\text{Tr}\left(\theta \stackrel{~}{𝐆}\right)\stackrel{~}{𝐆}+\mathrm{}.`$
After the duality transformation, the action becomes
$`S={\displaystyle \frac{g^2}{4}}{\displaystyle d^4x\left[\text{Tr}\left(\stackrel{~}{𝐆}\stackrel{~}{𝐆}\right)+2g^2\text{Tr}\left(\theta \stackrel{~}{𝐆}\stackrel{~}{𝐆}\stackrel{~}{𝐆}\right)\frac{g^2}{2}\text{Tr}\left(\theta \stackrel{~}{𝐆}\right)\text{Tr}\left(\stackrel{~}{𝐆}\stackrel{~}{𝐆}\right)+\mathrm{}\right]}.`$
This can be rewritten as
$`S={\displaystyle \frac{1}{4g_\mathrm{D}^2}}{\displaystyle d^4x\left[\text{Tr}\left(\mathrm{𝐆𝐆}\right)+2\text{Tr}\left(\theta _D\mathrm{𝐆𝐆𝐆}\right)\frac{1}{2}\text{Tr}\left(\theta _D𝐆\right)\text{Tr}\left(\mathrm{𝐆𝐆}\right)+\mathrm{}\right]},`$ (12)
where $`g_\mathrm{D}`$ and $`\theta _\mathrm{D}`$ are given as in Eq.(1). It is evident that the above action can be reorganized into a self-dual form
$`S={\displaystyle \frac{1}{4g_\mathrm{D}^2}}{\displaystyle d^4x\text{Tr}\left(𝒢\stackrel{~}{}𝒢\right)}`$ (13)
when expanded in powers of $`\theta _\mathrm{D}𝐆`$ and expressed noncommutativity via $`\theta _\mathrm{D}`$, where
$`𝒢_{\mu \nu }=𝐆_{\mu \nu }+𝐢\{𝐁_\mu ,𝐁_\nu \}_\stackrel{~}{},𝐆_{\mu \nu }=_\mu 𝐁_\nu _\nu 𝐁_\mu .`$
Thus, we find that the dual action Eq.(13) is again noncommutative U(1) gauge theory, but with dual coupling parameters Eq.(1). Note here that if we started with a theory with space-space non-commutativity, the dual theory will have space-time non-commutativity.
The foregoing analysis may be extended to the situation where R-R zero-form $`C`$ and two-form $`C_2`$ background is turned on. In this case, the full action takes the form:
$`S_{\mathrm{total}}=S+{\displaystyle d^4x\left[\frac{1}{2}C\text{Tr}\left(𝐅\stackrel{~}{𝐅}\right)+\text{Tr}\left(𝐅\stackrel{~}{C}_2\right)+\mathrm{}\right]}.`$
Here, one may wonder what form of the coupling one ought to use. One could for instance imagine that the correct coupling should be given by a standard form but in terms of the noncommutative field strength, for example, $`C\text{Tr}\left(\stackrel{~}{}\right)`$, etc. We have checked, again to the lowest order in noncommutativity parameter, that modifying the coupling in this fashion does not change the final result we will draw.
It is obvious that the effect of the R-R two-form potential is simply to shift $`𝐆𝐆+C_2`$ everywhere and, for the R-R zero-form, a straightforward calculation shows that the action is self-dual (i.e. of the same form as the original one) under the duality transformation provided that we define the dual parameters and backgrounds as:
$`g_D^2`$ $`=`$ $`{\displaystyle \frac{g^2}{\left(g^4+C^2\right)}}`$
$`C_D`$ $`=`$ $`{\displaystyle \frac{C}{\left(g^4+C^2\right)}}`$ (14)
$`\theta _D`$ $`=`$ $`C_D\theta {\displaystyle \frac{1}{g_D^2}}\stackrel{~}{\theta }.`$
These are the results anticipated from Type IIB S-duality:
$`\mathrm{S}_{\mathrm{IIB}}:(C+{\displaystyle \frac{i}{g^2}})`$ $``$ $`\left(C+{\displaystyle \frac{i}{g^2}}\right)^1`$
$`\left(\theta +i\stackrel{~}{\theta }\right)`$ $``$ $`\left(C+{\displaystyle \frac{i}{g^2}}\right)^1\left(\theta +i\stackrel{~}{\theta }\right).`$ (15)
Note that, had we started with an original theory having purely ‘electric’ or ‘magnetic’ noncommutativity, viz. $`\theta `$ is a tensor of rank-1, in the presence of the R-R scalar background, the dual theory would have noncommutativity tensor $`\theta _\mathrm{D}`$ of full rank.
## 3 Closer Look - Noncritical Open String
We have seen, in the previous section, that the naive S-dual of a noncommutative gauge theory is again a noncommutative gauge theory, but with the noncommutativity parameter $`\theta _\mathrm{D}=g^2\stackrel{~}{\theta }`$. Thus, if the original theory were defined with ‘magnetic’ noncommutativity, the dual theory would have ‘electric’ noncommutativity and vice versa. We also note that, in performing the S-duality map, as we have expanded the original and the dual actions as in Eq.(11) and Eq.(12), respectively, the result would be valid only for $`\theta 𝐅1`$ and $`\theta _\mathrm{D}𝐆1`$. The latter requires that $`g1`$ and, in this case, the physics in the original and the dual theories wouldn’t be different much from the standard commutative gauge theories <sup>3</sup><sup>3</sup>3Recall that, for pure U(1) gauge theories, the S-duality map is exact for all values of $`g`$..
What can one say for the $`g1`$ case, viz. the original theory is strongly coupled? In order to understand this limit, for definiteness, we will take the noncommutativity purely ‘magnetic’: $`\theta =\theta _{23}`$. The spectrum in this theory includes, in addition to the U(1) gauge boson, the magnetic monopole and the dyon (See and references therein). They may be viewed as noncommutative deformations of the photon, the magnetic monopoles and the dyons in the standard U(1) gauge theory. In noncommutative gauge theory, the U(1) gauge boson can be visualized as an induced electric dipole.
Alternatively, one may analyze the spectrum as fundamental (F-) or D-strings, respectively, ending on the D3-brane on which background magnetic field is turned on. The latter should be describable in terms of the Dirac-Born-Infeld Lagrangian ($`T_\mathrm{F}=1/2\pi \alpha ^{}`$):
$`L_{\mathrm{DBI}}={\displaystyle \frac{T_\mathrm{F}^2}{\lambda _{\mathrm{st}}}}\sqrt{g_{00}g_{11}}\sqrt{g_{22}g_{33}+𝐅_{23}^2/T_\mathrm{F}^2}`$
where, in the Seiberg-Witten decoupling limit, the bulk coupling parameter are related to the gauge theory parameters in Eq.(9) as
$`\lambda _{\mathrm{st}}={\displaystyle \frac{g^2}{T_\mathrm{F}\theta }},g_{00}=g_{11}=1,g_{22}=g_{33}=\left({\displaystyle \frac{1}{T_\mathrm{F}\theta }}\right)^2,𝐅_{23}=\left({\displaystyle \frac{1}{\theta }}\right)ϵ_{23}.`$
The magnetic monopoles are not part of the physical spectrum as, being represented by semi-infinite D-string ending on the D3-brane, they are counterpart of the Dirac magnetic monopole with an infinite self-energy. Among the physical excitations, however, are magnetic dipoles (as well as their dyonic counterparts) composed of monopole-antimonopole pair. See figure 1. Consider an open D-string (lying entirely on the D3-brane) of length $`\mathrm{\Delta }𝐱`$. It represents a magnetic dipole carrying a dipole moment $`𝐦=\mathrm{\Delta }𝐱`$ (measured in string unit) and total mass
$`M_{\mathrm{dipole}}=\left({\displaystyle \frac{T_\mathrm{F}}{\lambda _{\mathrm{st}}}}\right)|\mathrm{\Delta }𝐱|𝐦𝐇.`$ (16)
where the negative sign in the second term refers to relative opposite orientation between the dipole and the background magnetic field. The last term represents the interaction energy of the magnetic dipole with the background magnetic field:
$`𝐇={\displaystyle \frac{L_{\mathrm{DBI}}}{𝐅_{23}}}=𝐇_\mathrm{c}\left[1+\left({\displaystyle \frac{1}{T_\mathrm{F}\theta }}\right)^2\right]^{1/2},`$ (17)
where $`𝐇_\mathrm{c}`$ denotes the critical magnetic field strength
$`𝐇_\mathrm{c}={\displaystyle \frac{T_\mathrm{F}}{\lambda _{\mathrm{st}}}}.`$
In the field theory limit $`T_\mathrm{F}\mathrm{}`$, the magnetic dipole remain as low-energy excitaitons – they are noncritical open D-strings with an effective tension
$`T_{\mathrm{eff}}={\displaystyle \frac{M_{\mathrm{dipole}}}{|\mathrm{\Delta }𝐱|}}{\displaystyle \frac{1}{2g^2\theta }}.`$ (18)
In Eq.(16), the negative sign in the second term refers to relative opposite orientation between the dipole and the background magnetic field. It implies that the noncritical open D-string representing the magnetic dipole ought to be chiral: open D-string with opposite orientation, which represents magnetic anti-dipole, is separated by an infinite mass gap from the magnetic dipoles. Thus, from the open D-string point of view, the field theory limit $`T_\mathrm{F}\mathrm{}`$ amounts to taking non-relativistic limit and allows to expand Eq.(17) in power-series of $`(1/T_\mathrm{F}\theta )^2`$. This also account for physical origin of the numerical factor 1/2 in Eq.(18) <sup>4</sup><sup>4</sup>4A related observation was made recently by Klebanov and Maldacena in the context of (1+1)-dimensional noncommutative open string theory ..
Taking the strong coupling limit, $`g1`$, unlike the magnetic monopoles, the magnetic dipoles are nearly tensionless, weakly interacting degrees of freedom, while the U(1) gauge bosons are tensionless, strongly interacting degrees of freedom. Performing S-duality Eq.(1) to the dual gauge theory, the two are interchanged with each other: the dual electric dipoles are nearly tensionless, weakly interacting degrees of freedom and the dual magnetic dipoles are tensionless but strongly interacting degrees of freedom. They are made out of open F- and D-strings ending on the dual D3-brane worldvolume, on which a background dual electric field $`𝐆_{01}`$ is turned on. The background dual electric field ought to be near critical, as is anticipated from S-duality and evidenced by the fact that, for fixed $`\theta `$, $`\theta _\mathrm{D}=g^2\theta `$ becomes infinitely large. Indeed, combining the S-duality transformation of gauge theory parameters Eq.(1) and of bulk coupling parameters
$`\lambda _{\mathrm{D},\mathrm{st}}={\displaystyle \frac{1}{\lambda }}_{\mathrm{st}}\mathrm{and}g_{\mathrm{D},\mu \nu }={\displaystyle \frac{1}{\lambda }}_{\mathrm{st}}g_{\mu \nu },`$
we find that $`𝐅_{23}`$ is mapped to displacement field $`𝐃=(L_{\mathrm{DBI}}/𝐆_{01})`$ and the dual electric field $`𝐄_\mathrm{D}:=𝐆_{01}`$ is given by
$`𝐄_\mathrm{D}=𝐄_\mathrm{c}\left[1+\left({\displaystyle \frac{1}{g_\mathrm{D}^2T_\mathrm{F}\theta _\mathrm{D}}}\right)^2\right]^{1/2}\mathrm{where}𝐄_\mathrm{c}=g_\mathrm{D}^4T_\mathrm{F}^2\theta _\mathrm{D}.`$
Likewise, magnetic dipoles are mapped into electric dipoles made out of open F-strings, whose effective tension is given by the S-dual of Eq.(18):
$`\stackrel{~}{T}_{\mathrm{eff}}={\displaystyle \frac{1}{2\theta _\mathrm{D}}},`$ (19)
where again the factor of 1/2 ought to signify chiral nature of the open F-string. An immediate question is: are these nearly tensionless, open F-strings identifiable within the dual gauge theory?
A set of gauge invariant operators in the dual gauge theory is given by the following open Wilson lines :
$`\stackrel{~}{W}_𝐤[C]={\displaystyle d^2𝐱\mathrm{exp}_\stackrel{~}{}\left[i_C\dot{y}(t)𝐁(x+y(t))\right]\stackrel{~}{}e^{i𝐤𝐱}}.`$ (20)
Here, $`t=[0,1]`$ denotes the affine parameter along the open Wilson line, $`x^\mu `$ refers to the spacetime position of the $`\tau =0`$ point, $`𝐱`$ to the projection of $`x^\mu `$ onto the two-dimensional noncommutative spacetime and $`𝐤`$ to the Fourier-momentum along the noncommutative spacetime. All the multiplications are defined in terms of the ‘generalized Moyal product’:
$`A(x)\stackrel{~}{}B(y):=\mathrm{exp}\left({\displaystyle \frac{i}{2}}\theta _\mathrm{D}^{\mu \nu }_\mu ^x_\nu ^y\right)A(x)B(y).`$
It is then straightforward to check that the Wilson line Eq.(20) is gauge invariant provided the following relation holds between the momentum $`𝐤_\mu `$ and the endpoint separation distance:
$`𝐤_\mu =\left({\displaystyle \frac{1}{\theta _\mathrm{D}}}\right)_{\mu \nu }\mathrm{\Delta }y^\nu \mathrm{where}\mathrm{\Delta }yy(1)y(0).`$ (21)
What happens here is that, once the relation Eq.(21), the extra phase factor $`e^{i𝐤𝐱}`$ effectively parallel transports the gauge transformation parameter at $`t=1`$ back to that at $`t=0`$. This then ensures that the noncommutative Wilson line, despite being an open string, is gauge invariant. Being so, much as the closed Wilson loops form a complete set of gauge invariant observables in Yang-Mills theory, we can take the noncommutative Wilson lines Eq.(20) as a complete set of gauge invariant observables of the dual noncommutative U(1) gauge theory <sup>5</sup><sup>5</sup>5For noncommutative U(N) gauge theory, via Morita equivalence, we expect that the noncommutative Wilson lines Eq.(20) still form a complete set of gauge invariant observables..
For small total energy or momentum, $`|𝐤|1/\sqrt{\theta _\mathrm{D}}`$, separation between the Wilson line endpoints is shorter than the non-commutativity scale, $`|\mathrm{\Delta }y|\sqrt{\theta _\mathrm{D}}`$. Hence, the Wilson line reduces effectively to a (Fourier-transform) of the standard closed Wilson loop. For energies smaller than the non-commutativity scale $`1/\sqrt{\theta _\mathrm{D}}`$, we would expect the dual theory behaves as in the standard gauge theory. This agrees with the conclusions of . On the other hand, if $`|𝐤|1/\sqrt{\theta _\mathrm{D}}`$, then the separation of the Wilson line endpoints would be larger than the noncommutativity scale, $`|\mathrm{\Delta }y|\sqrt{\theta _\mathrm{D}}`$. These excitations are string-like.
Taking the zeroth component of Eq.(21), one finds
$`E_{\mathrm{YM}}:=\stackrel{~}{T}_{\mathrm{eff}}\mathrm{\Delta }y^1,`$ (22)
where $`\stackrel{~}{T}_{\mathrm{eff}}`$ coincides precisely with the effective tension Eq.(19). Thus, we are prompted to identify the Wilson line with a noncritical open F-string, whose effective tension is given by $`\stackrel{~}{T}_{\mathrm{eff}}`$. Recall that, under the S-duality Eq.(1), $`T_{\mathrm{eff}}=1/2g^2\theta `$ is mapped to $`\stackrel{~}{T}_{\mathrm{eff}}=1/2\theta _\mathrm{D}`$. Hence, the open Wilson lines in the dual gauge theory are the right candidates for the S-dual of the magnetic dipoles in the original gauge theory.
Characteristic size of the open Wilson lines is $`𝒪(\sqrt{\theta _\mathrm{D}})`$ and become macroscopically large in the strong coupling limit $`g1`$. They represent a complete set of excitations in the weakly coupled, dual gauge theory for $`|𝐤|1/\sqrt{\theta _\mathrm{D}}0`$. It clearly suggests that the noncommutative gauge theory captures more of the description than we had the right to expect from the naive duality argument in Section 2. It also indicates that a suitable formulation of the dual gauge theory is in terms of macroscopic open strings.
In extracting tension of the open Wilson line from Eq.(22), we were able to match it to Eq.(19) modulo a numerical factor of 2. Recall that the numerical factor of 1/2 in Eqs.(18, 19) has originated from chiral or, equivalently, non-relativistic nature of the open D- and F-strings. What then would cause the open Wilson lines of the noncommutative Yang-Mills theory chiral and eventually account for cancellation or disappearance of the factor of 2 in Eq. (22)?
We believe an answer to this question comes from the fact that ‘electric’ noncommutative Yang-Mills theory is parity-violating. For fixed $`\theta _\mathrm{D}`$, this is easily seen from non-invariance of the electric noncommutativity $`[x^0,x^1]=\theta _\mathrm{D}`$ under $`𝐱𝐱`$ <sup>6</sup><sup>6</sup>6Note, however, that ‘magnetic’ noncommutative Yang-Mills theory is parity-conserving: magnetic noncommutativity $`[x^2,x^3]=\theta `$ is invariant under $`𝐱𝐱`$.. This then implies that, along the $`(x^0x^1)`$-directions along which $`𝐤_\mu `$ and $`\mathrm{\Delta }y^\mu `$ in Eq.(21) point, the open Wilson lines ought to be chiral, stretching the two endpoints such that $`\mathrm{\Delta }y^1`$ positive. This chirality also implies that, from Eq.(22), only positive energy Wilson lines are physical excitations but not negative energy ones. The net result is essentially the same as that of the infinite mass gap opening up in the non-relativistic limit.
There exists one more piece of evidence that the Wilson lines are indeed identifiable with a sort of macroscopic string. Utilizing $`E\mathrm{\Delta }E\mathrm{}/\mathrm{\Delta }y^0`$, we observe that the Wilson line exhibits a version of the spacetime uncertainty relation:
$`\mathrm{\Delta }y^0\mathrm{\Delta }y^1{\displaystyle \frac{1}{2}}\mathrm{}\theta _\mathrm{D}.`$
As emphasized by Yoneya , the spacetime uncertainty relation is a distinguishing feature of any string theory with worldsheet conformal invariance.
## 4 Yet Another Look – Strong Noncommutativity Limit
We would like to discuss yet another piece of physics associated with the S-duality, Eq.(1). In the previous section, we have seen that, due to the dual electric field background, the open F-string is oriented predominantly along the $`x^0x^1`$ directions. See Eq.(21). There, we have also argued that the open string is macroscopically stretched. According to Eq.(20), the open string is made out of the dual gauge field $`𝐁_\mu `$ as a sort of coherent state configuration. Thus, it ought to be possible to visualize the open string out of the dual gauge theory in the semiclassical limit. In this section, under suitable condition, we show that the dual gauge theory Eq.(13) describes worldsheet dynamics of $`N_\mathrm{F}`$ coincident macroscopic F-strings propagating in four dimensional spacetime.
Let us begin with the following elementary observation. Strongly coupled noncommutative U(1) gauge theory with a finite noncommutativity $`\theta `$, as is seen from Eq.(1), is dual to weakly coupled noncommutative U(1) gauge theory with an infinite noncommutativity $`\theta _\mathrm{D}`$. In dimensionless measure, this implies that
$`\theta _\mathrm{D}\{𝐁_\mu ,𝐁_\nu \}\theta _\mathrm{D}\left(𝐁\right)\mathrm{\hspace{0.17em}\hspace{0.17em}1}`$ (23)
and hence corresponds to high field-strength, high-energy limit <sup>7</sup><sup>7</sup>7This is the limit considered originally by . Because of the infinitely large noncommutativity, dynamics of the dual ‘electric’ U(1) gauge theory is considerably simplified. At leading order in the noncommutativity Eq.(23), the dual gauge theory action Eq.(13) is reduced as:
$`S{\displaystyle \frac{1}{4g_\mathrm{D}^2}}V_{}{\displaystyle 𝑑x^0𝑑x^1<\left(𝒢_{\mu \nu }\right)_\stackrel{~}{}^2>},`$
where the Lorentz indices are contracted with gauge theory metric. We have also introduced notations
$`V_{}{\displaystyle 𝑑x^2𝑑x^3}\mathrm{and}<𝒪>{\displaystyle \frac{1}{V_{}}}{\displaystyle 𝑑x^2𝑑x^3𝒪}.`$ (24)
In the limit of infinitely many coincident noncommutative D3-branes, it is known that nonabelian generalization of the dual gauge theory Eq.(13) may be interpreted as a theory describing low-energy dynamics of (F1-D3) bound states . It is known that, in this case, the Yang-Mills gauge coupling is not arbitrary but is determined by the F-string and the D3-brane charges $`N_\mathrm{F},N_{\mathrm{D3}}`$ :
$`g_\mathrm{D}^2={\displaystyle \frac{V_{}}{\theta _\mathrm{D}}}\left({\displaystyle \frac{N_{\mathrm{D3}}}{N_\mathrm{F}}}\right).`$
Thus, using this parameter relation and introducing new covariant operator variables
$`X^\mu =\theta _\mathrm{D}^{\mu \nu }\left(i_\nu +𝐁_\nu \right),`$
we can re-express the dual ‘electric’ noncommutative Yang-Mills theory action as:
$`S={\displaystyle \frac{N_\mathrm{F}T_{\mathrm{eff}}}{2}}{\displaystyle 𝑑x^0𝑑x^1\left(\frac{1}{2}\left(\frac{1}{\theta _\mathrm{D}}\{X^\mu ,X^\nu \}_\stackrel{~}{}\right)^2+\mathrm{}\right)}`$ (25)
where again the Lorentz indices are contracted with respect to the gauge theory metric and the ellipses denote $`𝒪(1)`$ sub-leading terms. We have thus found that the resulting action Eq.(25) has precisely the same form as string worldsheet action for $`N_\mathrm{F}`$ multiple noncritical, F-strings of tension $`T_{\mathrm{eff}}`$, except that the action is expressed in the so-called Schild gauge . The F-strings ought to be interpreted as open strings, albeit infinitely stretched, as they propagate in a spacetime governed by the gauge theory metric. Note that the worldsheet direction is along $`(x^0x^1)`$-directions and the strings are delocalized along $`(x^2x^3)`$-directions.
Actually, what we have gotten is not precisely the Schild-gauge action but a deformation quantization of it. Namely, plaquette element $`\mathrm{\Sigma }^{\mu \nu }`$ of the string worldsheet is deformed into
$`\mathrm{\Sigma }^{\mu \nu }ϵ^{ab}_aX^\mu _bX^\nu =\{X^\mu ,X^\nu \}_{\mathrm{PB}}\mathrm{\Sigma }_\stackrel{~}{}^{\mu \nu }{\displaystyle \frac{1}{\theta _\mathrm{D}}}\{X^\mu ,X^\nu \}_\stackrel{~}{}.`$ (26)
We trust the deformation is correctly normalized, as $`\mathrm{\Sigma }_\stackrel{~}{}^{\mu \nu }=\mathrm{\Sigma }^{\mu \nu }+𝒪(\theta _D)`$ in small $`\theta _\mathrm{D}`$ limit.
What conclusion can one draw out of the above result? Had we considered $`N_3`$ coincident D3-branes with $`N_\mathrm{F}`$ units of center-of-mass U(1) electric flux turned on, we would have obtained the standard Nambu-Goto or Schild action of $`N_\mathrm{F}`$ F-strings. Recalling that noncommutative U(1) gauge theory is equivalent to $`\mathrm{U}(\mathrm{})`$ gauge theory at high-energy regime, we may interpret that the dual ‘electric’ noncommutative gauge theory indeed describes worldsheet dynamics of $`N_\mathrm{F}`$ coincident F-strings provided the latter carry high energy-momentum and become open strings (see Eq.(21)), and are delocalized along $`(x^2x^3)`$-directions. The result seems consistent with what one finds from supergravity dual .
Acknowledgement
We thank M.R. Douglas, J. Klusoň, U. Lindström and G. Moore for useful discussions. SJR thanks warm hospitality of Marsaryk University, New High-Energy Theory Center at Rutgers University, and Theory Division at CERN during completion of the work.
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# Random graphs with arbitrary degree distributions and their applications
## I Introduction
A random graph is a collection of points, or vertices, with lines, or edges, connecting pairs of them at random (Fig. 1a). The study of random graphs has a long history. Starting with the influential work of Paul Erdős and Alfréd Rényi in the 1950s and 1960s , random graph theory has developed into one of the mainstays of modern discrete mathematics, and has produced a prodigious number of results, many of them highly ingenious, describing statistical properties of graphs, such as distributions of component sizes, existence and size of a giant component, and typical vertex–vertex distances.
In almost all of these studies the assumption has been made that the presence or absence of an edge between two vertices is independent of the presence or absence of any other edge, so that each edge may be considered to be present with independent probability $`p`$. If there are $`N`$ vertices in a graph, and each is connected to an average of $`z`$ edges, then it is trivial to show that $`p=z/(N1)`$, which for large $`N`$ is usually approximated by $`z/N`$. The number of edges connected to any particular vertex is called the degree $`k`$ of that vertex, and has a probability distribution $`p_k`$ given by
$$p_k=\left(\genfrac{}{}{0pt}{}{N}{k}\right)p^k(1p)^{Nk}\frac{z^k\mathrm{e}^z}{k!},$$
(1)
where the second equality becomes exact in the limit of large $`N`$. This distribution we recognize as the Poisson distribution: the ordinary random graph has a Poisson distribution of vertex degrees, a point which turns out to be crucial, as we now explain.
Random graphs are not merely a mathematical toy; they have been employed extensively as models of real-world networks of various types, particularly in epidemiology. The passage of a disease through a community depends strongly on the pattern of contacts between those infected with the disease and those susceptible to it. This pattern can be depicted as a network, with individuals represented by vertices and contacts capable of transmitting the disease by edges. The large class of epidemiological models known as susceptible/infectious/recovered (or SIR) models makes frequent use of the so-called fully mixed approximation, which is the assumption that contacts are random and uncorrelated, i.e., that they form a random graph.
Random graphs however turn out to have severe shortcomings as models of such real-world phenomena. Although it is difficult to determine experimentally the structure of the network of contacts by which a disease is spread , studies have been performed of other social networks such as networks of friendships within a variety of communities , networks of telephone calls , airline timetables , and the power grid , as well as networks in physical or biological systems, including neural networks , the structure and conformation space of polymers , metabolic pathways , and food webs . It is found that the distribution of vertex degrees in many of these networks is measurably different from a Poisson distribution—often wildly different—and this strongly suggests, as has been emphasized elsewhere , that there are features of such networks which we would miss if we were to approximate them by an ordinary (Poisson) random graph.
Another very widely studied network is the internet, whose structure has attracted an exceptional amount of scrutiny, academic and otherwise, following its meteoric rise to public visibility starting in 1993. Pages on the world-wide web may be thought of as the vertices of a graph and the hyperlinks between them as edges. Empirical studies have shown that this graph has a distribution of vertex degree which is heavily right-skewed and possesses a fat (power-law) tail with an exponent between $`2`$ and $`3`$. (The underlying physical structure of the internet also has a degree distribution of this type .) This distribution is very far from Poisson, and therefore we would expect that a simple random graph would give a very poor approximation of the structural properties of the web. However, the web differs from a random graph in another way also: it is directed. Links on the web lead from one page to another in only one direction (see Fig. 1b). As discussed by Broder et al. this has a significant practical effect on the typical accessibility of one page from another, and this effect also will not be captured by a simple (undirected) random graph model.
A further class of networks that has attracted scrutiny is the class of collaboration networks. Examples of such networks include the boards of directors of companies , co-ownership networks of companies , and collaborations of scientists and movie actors . As well as having strongly non-Poisson degree distributions , these networks have a bipartite structure; there are two distinct kinds of vertices on the graph with links running only between vertices of unlike kinds —see Fig. 2. In the case of movie actors, for example, the two types of vertices are movies and actors, and the network can be represented as a graph with edges running between each movie and the actors that appear in it. Researchers have also considered the projection of this graph onto the unipartite space of actors only, also called a one-mode network . In such a projection two actors are considered connected if they have appeared in a movie together. The construction of the one-mode network however involves discarding some of the information contained in the original bipartite network, and for this reason it is more desirable to model collaboration networks using the full bipartite structure.
Given the high current level of interest in the structure of many of the graphs described here, and given their substantial differences from the ordinary random graphs that have been studied in the past, it would clearly be useful if we could generalize the mathematics of random graphs to non-Poisson degree distributions, and to directed and bipartite graphs. In this paper we do just that, demonstrating in detail how the statistical properties of each of these graph types can be calculated exactly in the limit of large graph size. We also give examples of the application of our theory to the modeling of a number of real-world networks, including the world-wide web and collaboration graphs.
## II Random graphs with arbitrary degree distributions
In this section we develop a formalism for calculating a variety of quantities, both local and global, on large unipartite undirected graphs with arbitrary probability distribution of the degrees of their vertices. In all respects other than their degree distribution, these graphs are assumed to be entirely random. This means that the degrees of all vertices are independent identically-distributed random integers drawn from a specified distribution. For a given choice of these degrees, also called the “degree sequence,” the graph is chosen uniformly at random from the set of all graphs with that degree sequence. All properties calculated in this paper are averaged over the ensemble of graphs generated in this way. In the limit of large graph size an equivalent procedure is to study only one particular degree sequence, averaging uniformly over all graphs with that sequence, where the sequence is chosen to approximate as closely as possible the desired probability distribution. The latter procedure can be thought of as a “microcanonical ensemble” for random graphs, where the former is a “canonical ensemble.”
Some results are already known for random graphs with arbitrary degree distributions: in two beautiful recent papers , Molloy and Reed have derived formulas for the position of the phase transition at which a giant component first appears, and the size of the giant component. (These results are calculated within the microcanonical ensemble, but apply equally to the canonical one in the large system size limit.) The formalism we present in this paper yields an alternative derivation of these results and also provides a framework for obtaining other quantities of interest, some of which we calculate. In Sections III and IV we extend our formalism to the case of directed graphs (such as the world-wide web) and bipartite graphs (such as collaboration graphs).
### A Generating functions
Our approach is based on generating functions, the most fundamental of which, for our purposes, is the generating function $`G_0(x)`$ for the probability distribution of vertex degrees $`k`$. Suppose that we have a unipartite undirected graph—an acquaintance network, for example—of $`N`$ vertices, with $`N`$ large. We define
$$G_0(x)=\underset{k=0}{\overset{\mathrm{}}{}}p_kx^k,$$
(2)
where $`p_k`$ is the probability that a randomly chosen vertex on the graph has degree $`k`$. The distribution $`p_k`$ is assumed correctly normalized, so that
$$G_0(1)=1.$$
(3)
The same will be true of all generating functions considered here, with a few important exceptions, which we will note at the appropriate point. Because the probability distribution is normalized and positive definite, $`G_0(x)`$ is also absolutely convergent for all $`|x|1`$, and hence has no singularities in this region. All the calculations of this paper will be confined to the region $`|x|1`$.
The function $`G_0(x)`$, and indeed any probability generating function, has a number of properties that will prove useful in subsequent developments.
Derivatives The probability $`p_k`$ is given by the $`k^{\mathrm{th}}`$ derivative of $`G_0`$ according to
$$p_k=\frac{1}{k!}\frac{\mathrm{d}^kG_0}{\mathrm{d}x^k}|_{x=0}.$$
(4)
Thus the one function $`G_0(x)`$ encapsulates all the information contained in the discrete probability distribution $`p_k`$. We say that the function $`G_0(x)`$ “generates” the probability distribution $`p_k`$.
Moments The average over the probability distribution generated by a generating function—for instance, the average degree $`z`$ of a vertex in the case of $`G_0(x)`$—is given by
$$z=k=\underset{k}{}kp_k=G_0^{}(1).$$
(5)
Thus if we can calculate a generating function we can also calculate the mean of the probability distribution which it generates. Higher moments of the distribution can be calculated from higher derivatives also. In general, we have
$$k^n=\underset{k}{}k^np_k=\left[\left(x\frac{\mathrm{d}}{\mathrm{d}x}\right)^nG_0(x)\right]_{x=1}.$$
(6)
Powers If the distribution of a property $`k`$ of an object is generated by a given generating function, then the distribution of the total of $`k`$ summed over $`m`$ independent realizations of the object is generated by the $`m^{\mathrm{th}}`$ power of that generating function. For example, if we choose $`m`$ vertices at random from a large graph, then the distribution of the sum of the degrees of those vertices is generated by $`[G_0(x)]^m`$. To see why this is so, consider the simple case of just two vertices. The square $`[G_0(x)]^2`$ of the generating function for a single vertex can be expanded as
$`[G_0(x)]^2`$ $`=`$ $`\left[{\displaystyle \underset{k}{}}p_kx^k\right]^2={\displaystyle \underset{jk}{}}p_jp_kx^{j+k}`$ (7)
$`=`$ $`p_0p_0x^0+(p_0p_1+p_1p_0)x^1`$ (10)
$`+(p_0p_2+p_1p_1+p_2p_0)x^2`$
$`+(p_0p_3+p_1p_2+p_2p_1+p_3p_0)x^3+\mathrm{}`$
It is clear that the coefficient of the power of $`x^n`$ in this expression is precisely the sum of all products $`p_jp_k`$ such that $`j+k=n`$, and hence correctly gives the probability that the sum of the degrees of the two vertices will be $`n`$. It is straightforward to convince oneself that this property extends also to all higher powers of the generating function.
All of these properties will be used in the derivations given in this paper.
Another quantity that will be important to us is the distribution of the degree of the vertices that we arrive at by following a randomly chosen edge. Such an edge arrives at a vertex with probability proportional to the degree of that vertex, and the vertex therefore has a probability distribution of degree proportional to $`kp_k`$. The correctly normalized distribution is generated by
$$\frac{\underset{k}{}kp_kx^k}{_kkp_k}=x\frac{G_0^{}(x)}{G_0^{}(1)}.$$
(12)
If we start at a randomly chosen vertex and follow each of the edges at that vertex to reach the $`k`$ nearest neighbors, then the vertices arrived at each have the distribution of remaining outgoing edges generated by this function, less one power of $`x`$, to allow for the edge that we arrived along. Thus the distribution of outgoing edges is generated by the function
$$G_1(x)=\frac{G_0^{}(x)}{G_0^{}(1)}=\frac{1}{z}G_0^{}(x),$$
(13)
where $`z`$ is the average vertex degree, as before. The probability that any of these outgoing edges connects to the original vertex that we started at, or to any of its other immediate neighbors, goes as $`N^1`$ and hence can be neglected in the limit of large $`N`$. Thus, making use of the “powers” property of the generating function described above, the generating function for the probability distribution of the number of second neighbors of the original vertex can be written as
$$\underset{k}{}p_k[G_1(x)]^k=G_0(G_1(x)).$$
(14)
Similarly, the distribution of third-nearest neighbors is generated by $`G_0(G_1(G_1(x)))`$, and so on. The average number $`z_2`$ of second neighbors is
$$z_2=\left[\frac{\mathrm{d}}{\mathrm{d}x}G_0(G_1(x))\right]_{x=1}=G_0^{}(1)G_1^{}(1)=G_0^{\prime \prime }(1),$$
(15)
where we have made use of the fact that $`G_1(1)=1`$. (One might be tempted to conjecture that since the average number of first neighbors is $`G_0^{}(1)`$, Eq. (5), and the average number of second neighbors is $`G_0^{\prime \prime }(1)`$, Eq. (15), then the average number of $`m`$th neighbors should be given by the $`m`$th derivative of $`G_0`$ evaluated at $`x=1`$. As we show in Section II F, however, this conjecture is wrong.)
### B Examples
To make things more concrete, we immediately introduce some examples of specific graphs to illustrate how these calculations are carried out.
#### a Poisson-distributed graphs
The simplest example of a graph of this type is one for which the distribution of degree is binomial, or Poisson in the large $`N`$ limit. This distribution yields the standard random graph studied by many mathematicians and discussed in Section I. In this graph the probability $`p=z/N`$ of the existence of an edge between any two vertices is the same for all vertices, and $`G_0(x)`$ is given by
$`G_0(x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{N}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{N}{k}}\right)p^k(1p)^{Nk}x^k`$ (16)
$`=`$ $`(1p+px)^N=\mathrm{e}^{z(x1)},`$ (17)
where the last equality applies in the limit $`N\mathrm{}`$. It is then trivial to show that the average degree of a vertex is indeed $`G_0^{}(1)=z`$ and that the probability distribution of degree is given by $`p_k=z^k\mathrm{e}^z/k!`$, which is the ordinary Poisson distribution. Notice also that for this special case we have $`G_1(x)=G_0(x)`$, so that the distribution of outgoing edges at a vertex is the same, regardless of whether we arrived there by choosing a vertex at random, or by following a randomly chosen edge. This property, which is peculiar to the Poisson-distributed random graph, is the reason why the theory of random graphs of this type is especially simple.
#### b Exponentially distributed graphs
Perhaps the next simplest type of graph is one with an exponential distribution of vertex degrees
$$p_k=(1\mathrm{e}^{1/\kappa })\mathrm{e}^{k/\kappa },$$
(18)
where $`\kappa `$ is a constant. The generating function for this distribution is
$$G_0(x)=(1\mathrm{e}^{1/\kappa })\underset{k=0}{\overset{\mathrm{}}{}}\mathrm{e}^{k/\kappa }x^k=\frac{1\mathrm{e}^{1/\kappa }}{1x\mathrm{e}^{1/\kappa }},$$
(19)
and
$$G_1(x)=\left[\frac{1\mathrm{e}^{1/\kappa }}{1x\mathrm{e}^{1/\kappa }}\right]^2.$$
(20)
An example of a graph with an exponential degree distribution is given in Section V A.
#### c Power-law distributed graphs
The recent interest in the properties of the world-wide web and of social networks leads us to investigate the properties of graphs with a power-law distribution of vertex degrees. Such graphs have been discussed previously by Barabási et al. and by Aiello et al.. In this paper, we will look at graphs with degree distribution given by
$$p_k=Ck^\tau \mathrm{e}^{k/\kappa }\text{for }k1\text{.}$$
(21)
where $`C`$, $`\tau `$, and $`\kappa `$ are constants. The reason for including the exponential cutoff is two-fold: first many real-world graphs appear to show this cutoff ; second it makes the distribution normalizable for all $`\tau `$, and not just $`\tau 2`$.
The constant $`C`$ is fixed by the requirement of normalization, which gives $`C=[Li_\tau (\mathrm{e}^{1/\kappa })]^1`$ and hence
$$p_k=\frac{k^\tau \mathrm{e}^{k/\kappa }}{Li_\tau (\mathrm{e}^{1/\kappa })}\text{for }k1\text{,}$$
(22)
where $`Li_n(x)`$ is the $`n`$th polylogarithm of $`x`$, a function familiar to those who have worked with Feynman integrals.
Substituting (22) into Eq. (2), we find that the generating function for graphs with this degree distribution is
$$G_0(x)=\frac{\underset{\tau }{Li}(x\mathrm{e}^{1/\kappa })}{Li_\tau (\mathrm{e}^{1/\kappa })}.$$
(23)
In the limit $`\kappa \mathrm{}`$—the case considered in Refs. and —this simplifies to
$$G_0(x)=\frac{\underset{\tau }{Li}(x)}{\zeta (\tau )},$$
(24)
where $`\zeta (x)`$ is the Riemann $`\zeta `$-function.
The function $`G_1(x)`$ is given by
$$G_1(x)=\frac{\underset{\tau 1}{Li}(x\mathrm{e}^{1/\kappa })}{xLi_{\tau 1}(\mathrm{e}^{1/\kappa })}.$$
(25)
Thus, for example, the average number of neighbors of a randomly-chosen vertex is
$$z=G_0^{}(1)=\frac{\underset{\tau 1}{Li}(\mathrm{e}^{1/\kappa })}{Li_\tau (\mathrm{e}^{1/\kappa })},$$
(26)
and the average number of second neighbors is
$$z_2=G_0^{\prime \prime }(1)=\frac{\underset{\tau 2}{Li}(\mathrm{e}^{1/\kappa })\underset{\tau 1}{Li}(\mathrm{e}^{1/\kappa })}{Li_\tau (\mathrm{e}^{1/\kappa })}.$$
(27)
#### d Graphs with arbitrary specified degree distribution
In some cases we wish to model specific real-world graphs which have known degree distributions—known because we can measure them directly. A number of the graphs described in the introduction fall into this category. For these graphs, we know the exact numbers $`n_k`$ of vertices having degree $`k`$, and hence we can write down the exact generating function for that probability distribution in the form of a finite polynomial:
$$G_0(x)=\frac{\underset{k}{}n_kx^k}{_kn_k},$$
(28)
where the sum in the denominator ensures that the generating function is properly normalized. As a example, suppose that in a community of 1000 people, each person knows between zero and five of the others, the exact numbers of people in each category being, from zero to five: $`\{86,150,363,238,109,54\}`$. This distribution will then be generated by the polynomial
$$G_0(x)=\frac{86+150x+363x^2+238x^3+109x^4+54x^5}{1000}.$$
(29)
### C Component sizes
We are now in a position to calculate some properties of interest for our graphs. First let us consider the distribution of the sizes of connected components in the graph. Let $`H_1(x)`$ be the generating function for the distribution of the sizes of components which are reached by choosing a random edge and following it to one of its ends. We explicitly exclude from $`H_1(x)`$ the giant component, if there is one; the giant component is dealt with separately below. Thus, except when we are precisely at the phase transition where the giant component appears, typical component sizes are finite, and the chances of a component containing a closed loop of edges goes as $`N^1`$, which is negligible in the limit of large $`N`$. This means that the distribution of components generated by $`H_1(x)`$ can be represented graphically as in Fig. 3; each component is tree-like in structure, consisting of the single site we reach by following our initial edge, plus any number (including zero) of other tree-like clusters, with the same size distribution, joined to it by single edges. If we denote by $`q_k`$ the probability that the initial site has $`k`$ edges coming out of it other than the edge we came in along, then, making use of the “powers” property of Section II A, $`H_1(x)`$ must satisfy a self-consistency condition of the form
$$H_1(x)=xq_0+xq_1H_1(x)+xq_2[H_1(x)]^2+\mathrm{}$$
(30)
However, $`q_k`$ is nothing other than the coefficient of $`x^k`$ in the generating function $`G_1(x)`$, Eq. (13), and hence Eq. (30) can also be written
$$H_1(x)=xG_1(H_1(x)).$$
(31)
If we start at a randomly chosen vertex, then we have one such component at the end of each edge leaving that vertex, and hence the generating function for the size of the whole component is
$$H_0(x)=xG_0(H_1(x)).$$
(32)
In principle, therefore, given the functions $`G_0(x)`$ and $`G_1(x)`$, we can solve Eq. (31) for $`H_1(x)`$ and substitute into Eq. (32) to get $`H_0(x)`$. Then we can find the probability that a randomly chosen vertex belongs to a component of size $`s`$ by taking the $`s`$th derivative of $`H_0`$. In practice, unfortunately, this is usually impossible; Eq. (31) is a complicated and frequently transcendental equation, which rarely has a known solution. On the other hand, we note that the coefficient of $`x^s`$ in the Taylor expansion of $`H_1(x)`$ (and therefore also the $`s`$th derivative) are given exactly by only $`s+1`$ iterations of Eq. (32), starting with $`H_1=1`$, so that the distribution generated by $`H_0(x)`$ can be calculated exactly to finite order in finite time. With current symbolic manipulation programs, it is quite possible to evaluate the first one hundred or so derivatives in this way. Failing this, an approximate solution can be found by numerical iteration and the distribution of cluster sizes calculated from Eq. (4) by numerical differentiation. Since direct evaluation of numerical derivatives is prone to machine-precision problems, we recommend evaluating the derivatives by numerical integration of the Cauchy formula, giving the probability distribution $`P_s`$ of cluster sizes thus:
$$P_s=\frac{1}{s!}\frac{\mathrm{d}^sH_0}{\mathrm{d}z^s}|_{z=0}=\frac{1}{2\pi \mathrm{i}}\frac{H_0(z)}{z^{s+1}}dz.$$
(33)
The best numerical precision is obtained by using the largest possible contour, subject to the condition that it enclose no poles of the generating function. The largest contour for which this condition is satisfied in general is the unit circle $`|z|=1`$ (see Section II A), and we recommend using this contour for Eq. (33). It is possible to find the first thousand derivatives of a function without difficulty using this method .
### D The mean component size, the phase transition, and the giant component
Although it is not usually possible to find a closed-form expression for the complete distribution of cluster sizes on a graph, we can find closed-form expressions for the average properties of clusters from Eqs. (31) and (32). For example, the average size of the component to which a randomly chosen vertex belongs, for the case where there is no giant component in the graph, is given in the normal fashion by
$$s=H_0^{}(1)=1+G_0^{}(1)H_1^{}(1).$$
(34)
From Eq. (31) we have
$$H_1^{}(1)=1+G_1^{}(1)H_1^{}(1),$$
(35)
and hence
$$s=1+\frac{G_0^{}(1)}{1G_1^{}(1)}=1+\frac{z_1^2}{z_1z_2},$$
(36)
where $`z_1=z`$ is the average number of neighbors of a vertex and $`z_2`$ is the average number of second neighbors. We see that this expression diverges when
$$G_1^{}(1)=1.$$
(37)
This point marks the phase transition at which a giant component first appears. Substituting Eqs. (2) and (13) into Eq. (37), we can also write the condition for the phase transition as
$$\underset{k}{}k(k2)p_k=0.$$
(38)
Indeed, since this sum increases monotonically as edges are added to the graph, it follows that the giant component exists if and only if this sum is positive. This result has been derived by different means by Molloy and Reed . An equivalent and intuitively reasonable statement, which can also be derived from Eq. (36), is that the giant component exists if and only if $`z_2>z_1`$.
Our generating function formalism still works when there is a giant component in the graph, but, by definition, $`H_0(x)`$ then generates the probability distribution of the sizes of components excluding the giant component. This means that $`H_0(1)`$ is no longer unity, as it is for the other generating functions considered so far, but instead takes the value $`1S`$, where $`S`$ is the fraction of the graph occupied by the giant component. We can use this to calculate the size of the giant component from Eqs. (31) and (32) thus:
$$S=1G_0(u),$$
(39)
where $`uH_1(1)`$ is the smallest non-negative real solution of
$$u=G_1(u).$$
(40)
This result has been derived in a different but equivalent form by Molloy and Reed , using different methods.
The correct general expression for the average component size, excluding the (formally infinite) giant component, if there is one, is
$`s`$ $`=`$ $`{\displaystyle \frac{H_0^{}(1)}{H_0(1)}}`$ (41)
$`=`$ $`{\displaystyle \frac{1}{H_0(1)}}\left[G_0(H_1(1))+{\displaystyle \frac{G_0^{}(H_1(1))G_1(H_1(1))}{1G_1^{}(H_1(1))}}\right]`$ (42)
$`=`$ $`1+{\displaystyle \frac{zu^2}{[1S][1G_1^{}(u)]}},`$ (43)
which is equivalent to (36) when there is no giant component ($`S=0`$, $`u=1`$).
For example, in the ordinary random graph with Poisson degree distribution, we have $`G_0(x)=G_1(x)=\mathrm{e}^{z(x1)}`$ (Eq. (17)), and hence we find simply that $`1S=u`$ is a solution of $`u=G_0(u)`$, or equivalently that
$$S=1\mathrm{e}^{zS}.$$
(44)
The average component size is given by
$$s=\frac{1}{1z+zS}.$$
(45)
These are all well-known results .
For graphs with purely power-law distributions (Eq. (22) with $`\kappa \mathrm{}`$), $`S`$ is given by (39) with $`u`$ the smallest non-negative real solution of
$$u=\frac{\underset{\tau 1}{Li}(u)}{u\zeta (\tau 1)}.$$
(46)
For all $`\tau 2`$ this gives $`u=0`$, and hence $`S=1`$, implying that a randomly chosen vertex belongs to the giant component with probability tending to 1 as $`\kappa \mathrm{}`$. For graphs with $`\tau >2`$, the probability of belonging to the giant component is strictly less than 1, even for infinite $`\kappa `$. In other words, the giant component essentially fills the entire graph for $`\tau 2`$, but not for $`\tau >2`$. These results have been derived by different means by Aiello et al. .
### E Asymptotic form of the cluster size distribution
A variety of results are known about the asymptotic properties of the coefficients of generating functions, some of which can usefully be applied to the distribution of cluster sizes $`P_s`$ generated by $`H_0(x)`$. Close to the phase transition, we expect the tail of the distribution $`P_s`$ to behave as
$$P_ss^\alpha \mathrm{e}^{s/s^{}},$$
(47)
where the constants $`\alpha `$ and $`s^{}`$ can be calculated from the properties of $`H_0(x)`$ as follows.
The cutoff parameter $`s^{}`$ is simply related to the radius of convergence $`|x^{}|`$ of the generating function , according to
$$s^{}=\frac{1}{\mathrm{log}|x^{}|}.$$
(48)
The radius of convergence $`|x^{}|`$ is equal to the magnitude of the position $`x^{}`$ of the singularity in $`H_0(x)`$ nearest to the origin. From Eq. (32) we see that such a singularity may arise either through a singularity in $`G_0(x)`$ or through one in $`H_1(x)`$. However, since the first singularity in $`G_0(x)`$ is known to be outside the unit circle (Section II A), and the first singularity in $`H_1(x)`$ tends to $`x=1`$ as we go to the phase transition (see below), it follows that, sufficiently close to the phase transition, the singularity in $`H_0(x)`$ closest to the origin is also a singularity in $`H_1(x)`$. With this result $`x^{}`$ is easily calculated.
Although we do not in general have a closed-form expression for $`H_1(x)`$, it is easy to derive one for its functional inverse. Putting $`w=H_1(x)`$ and $`x=H_1^1(w)`$ in Eq. (31) and rearranging, we find
$$x=H_1^1(w)=\frac{w}{G_1(w)}.$$
(49)
The singularity of interest corresponds to the point $`w^{}`$ at which the derivative of $`H_1^1(w)`$ is zero, which is a solution of
$$G_1(w^{})w^{}G_1^{}(w^{})=0.$$
(50)
Then $`x^{}`$ (and hence $`s^{}`$) is given by Eq. (49). Note that there is no guarantee that (50) has a finite solution, and that if it does not, then $`P_s`$ will not in general follow the form of Eq. (47).
When we are precisely at the phase transition of our system, we have $`G_1(1)=G_1^{}(1)=1`$, and hence the solution of Eq. (50) gives $`w^{}=x^{}=1`$—a result which we used above—and $`s^{}\mathrm{}`$. We can use the fact that $`x^{}=1`$ at the transition to calculate the value of the exponent $`\alpha `$ as follows. Expanding $`H_1^1(w)`$ about $`w^{}=1`$ by putting $`w=1+ϵ`$ in Eq. (49), we find that
$$H_1^1(1+ϵ)=1\frac{1}{2}G_1^{\prime \prime }(1)ϵ^2+\mathrm{O}(ϵ^3),$$
(51)
where we have made use of $`G_1(1)=G_1^{}(1)=1`$ at the phase transition. So long as $`G_1^{\prime \prime }(1)0`$, which in general it is not, this implies that $`H_1(x)`$ and hence also $`H_0(x)`$ are of the form
$$H_0(x)(1x)^\beta \text{as }x1\text{,}$$
(52)
with $`\beta =\frac{1}{2}`$. This exponent is related to the exponent $`\alpha `$ as follows. Equation (47) implies that $`H_0(x)`$ can be written in the form
$$H_0(x)=\underset{s=0}{\overset{a1}{}}P_sx^s+C\underset{s=a}{\overset{\mathrm{}}{}}s^\alpha \mathrm{e}^{s/s^{}}x^s+ϵ(a),$$
(53)
where $`C`$ is a constant and the last (error) term $`ϵ(a)`$ is assumed much smaller than the second term. The first term in this expression is a finite polynomial and therefore has no singularities on the finite plane; the singularity resides in the second term. Using this equation, the exponent $`\beta `$ can be written:
$`\beta `$ $`=`$ $`\underset{x1}{lim}\left[1+(x1){\displaystyle \frac{H_0^{\prime \prime }(x)}{H_0^{}(x)}}\right]`$ (54)
$`=`$ $`\underset{a\mathrm{}}{lim}\underset{x1}{lim}\left[{\displaystyle \frac{1}{x}}+{\displaystyle \frac{x1}{x}}{\displaystyle \frac{\underset{s=a}{\overset{\mathrm{}}{}}s^{2\alpha }x^{s1}}{_{s=a}^{\mathrm{}}s^{1\alpha }x^{s1}}}\right]`$ (55)
$`=`$ $`\underset{a\mathrm{}}{lim}\underset{x1}{lim}\left[{\displaystyle \frac{1}{x}}+{\displaystyle \frac{1x}{x\mathrm{log}x}}{\displaystyle \frac{\mathrm{\Gamma }(3\alpha ,a\mathrm{log}x)}{\mathrm{\Gamma }(2\alpha ,a\mathrm{log}x)}}\right],`$ (56)
where we have replaced the sums with integrals as $`a`$ becomes large, and $`\mathrm{\Gamma }(\nu ,\mu )`$ is the incomplete $`\mathrm{\Gamma }`$-function. Taking the limits in the order specified and rearranging for $`\alpha `$, we then get
$$\alpha =\beta +1=\frac{3}{2},$$
(57)
regardless of degree distribution, except in the special case where $`G_1^{\prime \prime }(1)`$ vanishes (see Eq. (51)). The result $`\alpha =\frac{3}{2}`$ was known previously for the ordinary Poisson random graph , but not for other degree distributions.
### F Numbers of neighbors and average path length
We turn now to the calculation of the number of neighbors who are $`m`$ steps away from a randomly chosen vertex. As shown in Section II A, the probability distributions for first- and second-nearest neighbors are generated by the functions $`G_0(x)`$ and $`G_0(G_1(x))`$. By extension, the distribution of $`m`$th neighbors is generated by $`G_0(G_1(\mathrm{}G_1(x)\mathrm{}))`$, with $`m1`$ iterations of the function $`G_1`$ acting on itself. If we define $`G^{(m)}(x)`$ to be this generating function for $`m`$th neighbors, then we have
$$G^{(m)}(x)=\{\begin{array}{cc}G_0(x)\hfill & \text{for }m=1\hfill \\ G^{(m1)}(G_1(x))\hfill & \text{for }m2\text{.}\hfill \end{array}$$
(58)
Then the average number $`z_m`$ of $`m`$th-nearest neighbors is
$$z_m=\frac{\mathrm{d}G^{(m)}}{\mathrm{d}x}|_{x=1}=G_1^{}(1)G_{}^{(m1)}{}_{}{}^{}(1)=G_1^{}(1)z_{m1}.$$
(59)
Along with the initial condition $`z_1=z=G_0^{}(1)`$, this then tells us that
$$z_m=[G_1^{}(1)]^{m1}G_0^{}(1)=\left[\frac{z_2}{z_1}\right]^{m1}z_1.$$
(60)
From this result we can make an estimate of the typical length $`\mathrm{}`$ of the shortest path between two randomly chosen vertices on the graph. This typical path length is reached approximately when the total number of neighbors of a vertex out to that distance is equal to the number of vertices on the graph, i.e., when
$$1+\underset{m=1}{\overset{\mathrm{}}{}}z_m=N.$$
(61)
Using Eq. (60) this gives us
$$\mathrm{}=\frac{\mathrm{log}[(N1)(z_2z_1)+z_1^2]\mathrm{log}z_1^2}{\mathrm{log}(z_2/z_1)}.$$
(62)
In the common case where $`Nz_1`$ and $`z_2z_1`$, this reduces to
$$\mathrm{}=\frac{\mathrm{log}(N/z_1)}{\mathrm{log}(z_2/z_1)}+1.$$
(63)
This result is only approximate for two reasons. First, the conditions used to derive it are only an approximation; the exact answer depends on the detailed structure of the graph. Second, it assumes that all vertices are reachable from a randomly chosen starting vertex. In general however this will not be true. For graphs with no giant component it is certainly not true and Eq. (63) is meaningless. Even when there is a giant component however, it is usually not the case that it fills the entire graph. A better approximation to $`\mathrm{}`$ may therefore be given by replacing $`N`$ in Eq. (63) by $`NS`$, where $`S`$ is the fraction of the graph occupied by the giant component, as in Section II D.
Such shortcomings notwithstanding, there are a number of remarkable features of Eq. (63):
1. It shows that the average vertex–vertex distance for all random graphs, regardless of degree distribution, should scale logarithmically with size $`N`$, according to $`\mathrm{}=A+B\mathrm{log}N`$, where $`A`$ and $`B`$ are constants. This result is of course well-known for a number of special cases.
2. It shows that the average distance, which is a global property, can be calculated from a knowledge only of the average numbers of first- and second-nearest neighbors, which are local properties. It would be possible therefore to measure these numbers empirically by purely local measurements on a graph such as an acquaintance network and from them to determine the expected average distance between vertices. For some networks at least, this gives a surprisingly good estimate of the true average distance .
3. It shows that only the average numbers of first- and second-nearest neighbors are important to the calculation of average distances, and thus that two random graphs with completely different distributions of vertex degrees, but the same values of $`z_1`$ and $`z_2`$, will have the same average distances.
For the case of the purely theoretical example graphs we discussed earlier, we cannot make an empirical measurement of $`z_1`$ and $`z_2`$, but we can still employ Eq. (63) to calculate $`\mathrm{}`$. In the case of the ordinary (Poisson) random graph, for instance, we find from Eq. (17) that $`z_1=z`$, $`z_2=z^2`$, and so $`\mathrm{}=\mathrm{log}N/\mathrm{log}z`$, which is the standard result for graphs of this type . For the graph with degree distributed according to the truncated power law, Eq. (22), $`z_1`$ and $`z_2`$ are given by Eqs. (26) and (27), and the average vertex–vertex distance is
$$\mathrm{}=\frac{\mathrm{log}N+\mathrm{log}\left[\underset{\tau }{Li}(\mathrm{e}^{1/\kappa })/\underset{\tau 1}{Li}(\mathrm{e}^{1/\kappa })\right]}{\mathrm{log}\left[Li_{\tau 2}(\mathrm{e}^{1/\kappa })/Li_{\tau 1}(\mathrm{e}^{1/\kappa })1\right]}+1.$$
(64)
In the limit $`\kappa \mathrm{}`$, this becomes
$$\mathrm{}=\frac{\mathrm{log}N+\mathrm{log}\left[\zeta (\tau )/\zeta (\tau 1)\right]}{\mathrm{log}\left[\zeta (\tau 2)/\zeta (\tau 1)1\right]}+1.$$
(65)
Note that this expression does not have a finite positive real value for any $`\tau <3`$, indicating that one must specify a finite cutoff $`\kappa `$ for the degree distribution to get a well-defined average vertex–vertex distance on such graphs.
### G Simulation results
As a check on the results of this section, we have performed extensive computer simulations of random graphs with various distributions of vertex degree. Such graphs are relatively straightforward to generate. First, we generate a set of $`N`$ random numbers $`\left\{k_i\right\}`$ to represent the degrees of the $`N`$ vertices in the graph. These may be thought of as the “stubs” of edges, emerging from their respective vertices. Then we choose pairs of these stubs at random and place edges on the graph joining them up. It is simple to see that this will generate all graphs with the given set of vertex degrees with equal probability. The only small catch is that the sum $`_ik_i`$ of the degrees must be even, since each edge added to the graph must have two ends. This is not difficult to contrive however. If the set $`\left\{k_i\right\}`$ is such that the sum is odd, we simply throw it away and generate a new set.
As a practical matter, integers representing vertex degrees with any desired probability distribution can be generated using the transformation method if applicable, or failing that, a rejection or hybrid method . For example, degrees obeying the power-law-plus-cutoff form of Eq. (22) can be generated using a two-step hybrid transformation/rejection method as follows. First, we generate random integers $`k1`$ with distribution proportional to $`\mathrm{e}^{k/\kappa }`$ using the transformation
$$k=\kappa \mathrm{log}(1r),$$
(66)
where $`r`$ is a random real number uniformly distributed in the range $`0r<1`$. Second, we accept this number with probability $`k^\tau `$, where by “accept” we mean that if the number is not accepted we discard it and generate another one according to Eq. (66), repeating the process until one is accepted.
In Fig. 4 we show results for the size of the giant component in simulations of undirected unipartite graphs with vertex degrees distributed according to Eq. (22) for a variety of different values of $`\tau `$ and $`\kappa `$. On the same plot we also show the expected value of the same quantity derived by numerical solution of Eqs. (39) and (40). As the figure shows, the agreement between simulation and theory is excellent.
## III Directed graphs
We turn now to directed graphs with arbitrary degree distributions. An example of a directed graph is the world-wide web, since every hyperlink between two pages on the web goes in only one direction. The web has a degree distribution that follows a power-law, as discussed in Section I.
Directed graphs introduce a subtlety that is not present in undirected ones, and which becomes important when we apply our generating function formalism. In a directed graph it is not possible to talk about a “component”—i.e., a group of connected vertices—because even if vertex A can be reached by following (directed) edges from vertex B, that does not necessarily mean that vertex B can be reached from vertex A. There are two correct generalizations of the idea of the component to a directed graph: the set of vertices which are reachable from a given vertex, and the set from which a given vertex can be reached. We will refer to these as “out-components” and “in-components” respectively. An in-component can also be thought of as those vertices reachable by following edges backwards (but not forwards) from a specified vertex. It is possible to study directed graphs by allowing both forward and backward traversal of edges (see Ref. , for example). In this case, however, the graph effectively becomes undirected and should be treated with the formalism of Section II.
With these considerations in mind, we now develop the generating function formalism appropriate to random directed graphs with arbitrary degree distributions.
### A Generating functions
In a directed graph, each vertex has separate in-degree and out-degree for links running into and out of that vertex. Let us define $`p_{jk}`$ to be the probability that a randomly chosen vertex has in-degree $`j`$ and out-degree $`k`$. It is important to realize that in general this joint distribution of $`j`$ and $`k`$ is not equal to the product $`p_jp_k`$ of the separate distributions of in- and out-degree. In the world-wide web, for example, it seems likely (although this question has not been investigated to our knowledge) that sites with a large number of outgoing links also have a large number of incoming ones, i.e., that $`j`$ and $`k`$ are correlated, so that $`p_{jk}p_jp_k`$. We appeal to those working on studies of the structure of the web to measure the joint distribution of in- and out-degrees of sites; empirical data on this distribution would make theoretical work much easier!
We now define a generating function for the joint probability distribution of in- and out-degrees, which is necessarily a function of two independent variables, $`x`$ and $`y`$, thus:
$$𝒢(x,y)=\underset{jk}{}p_{jk}x^jy^k.$$
(67)
Since every edge on a directed graph must leave some vertex and enter another, the net average number of edges entering a vertex is zero, and hence $`p_{jk}`$ must satisfy the constraint
$$\underset{jk}{}(jk)p_{jk}=0.$$
(68)
This implies that $`𝒢(x,y)`$ must satisfy
$$\frac{𝒢}{x}|_{x,y=1}=\frac{𝒢}{y}|_{x,y=1}=z,$$
(69)
where $`z`$ is the average degree (both in and out) of vertices in the graph.
Using the function $`𝒢(x,y)`$, we can, as before, define generating functions $`G_0`$ and $`G_1`$ for the number of out-going edges leaving a randomly chosen vertex, and the number leaving the vertex reached by following a randomly chosen edge. We can also define generating functions $`F_0`$ and $`F_1`$ for the number arriving at such a vertex. These functions are given by
$`F_0(x)`$ $`=`$ $`𝒢(x,1),F_1(x)={\displaystyle \frac{1}{z}}{\displaystyle \frac{𝒢}{y}}|_{y=1},`$ (70)
$`G_0(y)`$ $`=`$ $`𝒢(1,y),G_1(y)={\displaystyle \frac{1}{z}}{\displaystyle \frac{𝒢}{x}}|_{x=1}.`$ (71)
Once we have these functions, many results follow as before. The average numbers of first and second neighbors reachable from a randomly chosen vertex are given by Eq. (69) and
$$z_2=G_0^{}(1)G_1^{}(1)=\frac{^2𝒢}{xy}|_{x,y=1}.$$
(72)
These are also the numbers of first and second neighbors from which a random vertex can be reached, since Eqs. (69) and (72) are manifestly symmetric in $`x`$ and $`y`$. We can also make an estimate of the average path length on the graph from
$$\mathrm{}=\frac{\mathrm{log}(N/z_1)}{\mathrm{log}(z_2/z_1)}+1,$$
(73)
as before. However, this equation should be used with caution. As discussed in Section II F, the derivation of this formula assumes that we are in a regime in which the bulk of the graph is reachable from most vertices. On a directed graph however, this may be far from true, as appears to be the case with the world-wide web .
The probability distribution of the numbers of vertices reachable from a randomly chosen vertex in a directed graph—i.e., of the sizes of the out-components—is generated by the function $`H_0(y)=yG_0(H_1(y))`$, where $`H_1(y)`$ is a solution of $`H_1(y)=yG_1(H_1(y))`$, just as before. (A similar and obvious pair of equations governs the sizes of the in-components.) The results for the asymptotic behavior of the component size distribution from Section II E generalize straightforwardly to directed graphs. The average out-component size for the case where there is no giant component is given by Eq. (36), and thus the point at which a giant component first appears is given once more by $`G_1^{}(1)=1`$. Substituting Eq. (67) into this expression gives the explicit condition
$$\underset{jk}{}(2jkjk)p_{jk}=0$$
(74)
for the first appearance of the giant component. This expression is the equivalent for the directed graph of Eq. (38). It is also possible, and equally valid, to define the position at which the giant component appears by $`F_1^{}(1)=1`$, which provides an alternative derivation for Eq. (74).
Just as with the individual in- and out-components for vertices, the size of the giant component on a directed graph can also be defined in different ways. The giant component can be represented using the “bow-tie” diagram of Broder et al. , which we depict (in simplified form) in Fig. 5. The diagram has three parts. The strongly connected portion of the giant component, represented by the central circle, is that portion in which every vertex can be reached from every other. The two sides of the bow-tie represent (1) those vertices from which the strongly connected component can be reached but which it is not possible to reach from the strongly connected component and (2) those vertices which can be reached from the strongly connected component but from which it is not possible to reach the strongly connected component. The solution of Eqs. (39) and (40) with $`G_0(x)`$ and $`G_1(x)`$ defined according to Eq. (71) gives the number of vertices, as a fraction of $`N`$, in the giant strongly connected component plus those vertices from which the giant strongly connected component can be reached. Using $`F_0(x)`$ and $`F_1(x)`$ (Eq. (70)) in place of $`G_0(x)`$ and $`G_1(x)`$ gives a different solution, which represents the fraction of the graph in the giant strongly connected component plus those vertices which can be reached from it.
### B Simulation results
We have performed simulations of directed graphs as a check on the results above. Generation of random directed graphs with known joint degree distribution $`p_{jk}`$ is somewhat more complicated than generation of undirected graphs discussed in Section II G. The method we use is as follows. First, it is important to ensure that the averages of the distributions of in- and out-degree of the graph are the same, or equivalently that $`p_{jk}`$ satisfies Eq. (68). If this is not the case, at least to good approximation, then generation of the graph will be impossible. Next, we generate a set of $`N`$ in/out-degree pairs $`(j_i,k_i)`$, one for each vertex $`i`$, according to the joint distribution $`p_{jk}`$, and calculate the sums $`_ij_i`$ and $`_ik_i`$. These sums are required to be equal if there are to be no dangling edges in the graph, but in most cases we find that they are not. To rectify this we use a simple procedure. We choose a vertex $`i`$ at random, discard the numbers $`(j_i,k_i)`$ for that vertex and generate new ones from the distribution $`p_{jk}`$. We repeat this procedure until the two sums are found to be equal. Finally, we choose random in/out pairs of edges and join them together to make a directed graph. The resulting graph has the desired number of vertices and the desired joint distribution of in and out degree.
We have simulated directed graphs in which the distribution $`p_{jk}`$ is given by a simple product of independent distributions of in- and out-degree. (As pointed out in Section III A, this is not generally the case for real-world directed graphs, where in- and out-degree may be correlated.) In Fig. 6 we show results from simulations of graphs with identically distributed (but independent) in- and out-degrees drawn from the exponential distribution, Eq. (18). For this distribution, solution of the critical-point equation $`G_1^{}(1)=1`$ shows that the giant component first appears at $`\kappa _c=[\mathrm{log}2]^1=1.4427`$. The three curves in the figure show the distribution of numbers of vertices accessible from each vertex in the graph for $`\kappa =0.5`$, $`0.8`$, and $`\kappa _c`$. The critical distribution follows a power-law form (see Section II C), while the others show an exponential cutoff. We also show the exact distribution derived from the coefficients in the expansion of $`H_1(x)`$ about zero. Once again, theory and simulation are in good agreement. A fit to the distribution for the case $`\kappa =\kappa _c`$ gives a value of $`\alpha =1.50\pm 0.02`$, in good agreement with Eq. (57).
## IV Bipartite graphs
The collaboration graphs of scientists, company directors, and movie actors discussed in Section I are all examples of bipartite graphs. In this section we study the theory of bipartite graphs with arbitrary degree distributions. To be concrete, we will speak in the language of “actors” and “movies,” but clearly all the developments here are applicable to academic collaborations, boards of directors, or any other bipartite graph structure.
### A Generating functions and basic results
Consider then a bipartite graph of $`M`$ movies and $`N`$ actors, in which each actor has appeared in an average of $`\mu `$ movies and each movie has a cast of average size $`\nu `$ actors. Note that only three of these parameters are independent, since the fourth is given by the equality
$$\frac{\mu }{M}=\frac{\nu }{N}.$$
(75)
Let $`p_j`$ be the probability distribution of the degree of actors (i.e., of the number of movies in which they have appeared) and $`q_k`$ be the distribution of degree (i.e., cast size) of movies. We define two generating functions which generate these probability distributions thus:
$$f_0(x)=\underset{j}{}p_jx^j,g_0(x)=\underset{k}{}q_kx^k.$$
(76)
(It may be helpful to think of $`f`$ as standing for “film,” in order to keep these two straight.) As before, we necessarily have
$$f_0(1)=g_0(1)=1,f_0^{}(1)=\mu ,g_0^{}(1)=\nu .$$
(77)
If we now choose a random edge on our bipartite graph and follow it both ways to reach the movie and actor which it connects, then the distribution of the number of other edges leaving those two vertices is generated by the equivalent of (13):
$$f_1(x)=\frac{1}{\mu }f_0^{}(x),g_1(x)=\frac{1}{\nu }g_0^{}(x).$$
(78)
Now we can write the generating function for the distribution of the number of co-stars (i.e., actors in shared movies) of a randomly chosen actor as
$$G_0(x)=f_0(g_1(x)).$$
(79)
If we choose a random edge, then the distribution of number of co-stars of the actor to which it leads is generated by
$$G_1(x)=f_1(g_1(x)).$$
(80)
These two functions play the same role in the one-mode network of actors as the functions of the same name did for the unipartite random graphs of Section II. Once we have calculated them, all the results from Section II follow exactly as before.
The numbers of first and second neighbors of a randomly chosen actor are
$`z_1`$ $`=`$ $`G_0^{}(1)=f_0^{}(1)g_1^{}(1),`$ (81)
$`z_2`$ $`=`$ $`G_0^{}(1)G_1^{}(1)=f_0^{}(1)f_1^{}(1)[g_1^{}(1)]^2.`$ (82)
Explicit expressions for these quantities can be obtained by substituting from Eqs. (76) and (78). The average vertex–vertex distance on the one-mode graph is given as before by Eq. (63). Thus, it is possible to estimate average distances on such graphs by measuring only the numbers of first and second neighbors.
The distribution of the sizes of the connected components in the one-mode network is generated by Eq. (32), where $`H_1(x)`$ is a solution of Eq. (31). The asymptotic results of Section II E generalize simply to the bipartite case, and the average size of a connected component in the absence of a giant component is
$$s=1+\frac{G_0^{}(1)}{1G_1^{}(1)},$$
(83)
as before. This diverges when $`G_1^{}(1)=1`$, marking the first appearance of the giant component. Equivalently, the giant component first appears when
$$f_0^{\prime \prime }(1)g_0^{\prime \prime }(1)=f_0^{}(1)g_0^{}(1).$$
(84)
Substituting from Eq. (76), we then derive the explicit condition for the first appearance of the giant component:
$$\underset{jk}{}jk(jkjk)p_jq_k=0.$$
(85)
The size $`S`$ of the giant component, as a fraction of the total number $`N`$ of actors, is given as before by the solution of Eqs. (39) and (40).
Of course, all of these results work equally well if “actors” and “movies” are interchanged. One can calculate the average distance between movies in terms of common actors shared, the size and distribution of connected components of movies, and so forth, using the formulas given above, with only the exchange of $`f_0`$ and $`f_1`$ for $`g_0`$ and $`g_1`$. The formula (84) is, not surprisingly, invariant under this interchange, so that the position of the onset of the giant component is the same regardless of whether one is looking at actors or movies.
### B Clustering
Watts and Strogatz have introduced the concept of clustering in social networks, also sometimes called network transitivity. Clustering refers to the increased propensity of pairs of people to be acquainted with one another if they have another acquaintance in common. Watts and Strogatz defined a clustering coefficient which measures the degree of clustering on a graph. For our purposes, the definition of this coefficient is
$$C=\frac{3\times \text{ number of triangles on the graph}}{\text{number of connected triples of vertices}}=\frac{3N_{\mathrm{}}}{N_3}.$$
(86)
Here “triangles” are trios of vertices each of which is connected to both of the others, and “connected triples” are trios in which at least one is connected to both the others. The factor of 3 in the numerator accounts for the fact that each triangle contributes to three connected triples of vertices, one for each of its three vertices. With this factor of 3, the value of $`C`$ lies strictly in the range from zero to one. In the directed and undirected unipartite random graphs of Sections II and III, $`C`$ is trivially zero in the limit $`N\mathrm{}`$. In the one-mode projections of bipartite graphs, however, both the actors and the movies can be expected to have non-zero clustering. We here treat the case for actors. The case for movies is easily derived by swapping $`f`$s and $`g`$s.
An actor who has $`zz_1`$ co-stars in total contributes $`\frac{1}{2}z(z1)`$ connected triples to $`N_3`$, so that
$$N_3=\frac{1}{2}N\underset{z}{}z(z1)r_z,$$
(87)
where $`r_z`$ is the probability of having $`z`$ co-stars. As shown above (Eq. (79)), the distribution $`r_z`$ is generated by $`G_0(x)`$ and so
$$N_3=\frac{1}{2}NG_0^{\prime \prime }(1).$$
(88)
A movie which stars $`k`$ actors contributes $`\frac{1}{6}k(k1)(k2)`$ triangles to the total triangle count in the one-mode graph. Thus the total number of triangles on the graph is the sum of $`\frac{1}{6}k(k1)(k2)`$ over all movies, which is given by
$$N_{\mathrm{}}=\frac{1}{6}M\underset{k}{}k(k1)(k2)q_k=\frac{1}{6}Mg_0^{\prime \prime \prime }(1).$$
(89)
Substituting into Eq. (86), we then get
$$C=\frac{M}{N}\frac{g_0^{\prime \prime \prime }(1)}{G_0^{\prime \prime }(1)}.$$
(90)
Making use of Eqs. (75), (76), and (79), this can also be written as
$$\frac{1}{C}1=\frac{(\mu _2\mu _1)(\nu _2\nu _1)^2}{\mu _1\nu _1(2\nu _13\nu _2+\nu _3)},$$
(91)
where $`\mu _n=_kk^np_k`$ is the $`n`$th moment of the distribution of numbers of movies in which actors have appeared, and $`\nu _n`$ is the same for cast size (number of actors in a movie).
### C Example
To give an example, consider a random bipartite graph with Poisson-distributed numbers of both movies per actor and actors per movie. In this case, following the derivation of Eq. (17), we find that
$$f_0(x)=\mathrm{e}^{\mu (x1)},g_0(x)=\mathrm{e}^{\nu (x1)},$$
(92)
and $`f_1(x)=f_0(x)`$ and $`g_1(x)=g_0(x)`$. Thus
$$G_0(x)=G_1(x)=\mathrm{e}^{\mu (\mathrm{e}^{\nu (x1)}1)}.$$
(93)
This implies that $`z_1=\mu \nu `$ and $`z_2=(\mu \nu )^2`$, so that
$$\mathrm{}=\frac{\mathrm{log}N}{\mathrm{log}\mu \nu }=\frac{\mathrm{log}N}{\mathrm{log}z},$$
(94)
just as in an ordinary Poisson-distributed random graph. From Eq. (83), the average size $`s`$ of a connected component of actors, below the phase transition, is
$$s=\frac{1}{1\mu \nu },$$
(95)
which diverges, yielding a giant component, at $`\mu \nu =z=1`$, also as in the ordinary random graph. From Eqs. (39) and (40), the size $`S`$ of the giant component as a fraction of $`N`$ is a solution of
$$S=1\mathrm{e}^{\mu (\mathrm{e}^{\nu S}1)}.$$
(96)
And from Eq. (90), the clustering coefficient for the one-mode network of actors is
$$C=\frac{M\nu ^3}{N\nu ^2(\mu ^2+\mu )}=\frac{1}{\mu +1},$$
(97)
where we have made use of Eq. (75).
Another quantity of interest is the distribution of numbers of co-stars, i.e., of the numbers of people with whom each actor has appeared in a movie. As discussed above, this distribution is generated by the function $`G_0(x)`$ defined in Eq. (79). For the case of the Poisson degree distribution, we can perform the derivatives, Eq. (4), and setting $`x=0`$ we find that the probability $`r_z`$ of having appeared with a total of exactly $`z`$ co-stars is
$$r_z=\frac{\nu ^z}{z!}\mathrm{e}^{\mu (\mathrm{e}^\nu 1)}\underset{k=1}{\overset{z}{}}\left\{\genfrac{}{}{0pt}{}{z}{k}\right\}\left[\mu \mathrm{e}^\nu \right]^k,$$
(98)
where the coefficients $`\left\{\genfrac{}{}{0pt}{}{z}{k}\right\}`$ are the Stirling numbers of the second kind
$$\left\{\genfrac{}{}{0pt}{}{z}{k}\right\}=\underset{r=1}{\overset{k}{}}\frac{(1)^{kr}}{r!(kr)!}r^z.$$
(99)
### D Simulation results
Random bipartite graphs can be generated using an algorithm similar to the one described in Section III B for directed graphs. After making sure that the required degree distributions for both actor and movie vertices have means consistent with the required total numbers of actors and movies according to Eq. (75), we generate vertex degrees for each actor and movie at random and calculate their sum. If these sums are unequal, we discard the degree of one actor and one movie, chosen at random, and replace them with new degrees drawn from the relevant distributions. We repeat this process until the total actor and movie degrees are equal. Then we join vertices up in pairs.
In Fig. 7 we show the results of such a simulation for a bipartite random graph with Poisson degree distribution. (In fact, for the particular case of the Poisson distribution, the graph can be generated simply by joining up actors and movies at random, without regard for individual vertex degrees.) The figure shows the distribution of the number of co-stars of each actor, along with the analytic solution, Eqs. (98) and (99). Once more, numerical and analytic results are in good agreement.
## V Applications to real-world networks
In this section we construct random graph models of two types of real-world networks, namely collaboration graphs and the world-wide web, using the results of Sections III and IV to incorporate realistic degree distributions into the models. As we will show, the results are in reasonably good agreement with empirical data, although there are some interesting discrepancies also, perhaps indicating the presence of social phenomena that are not incorporated in the random graph.
### A Collaboration networks
In this section we construct random bipartite graph models of the known collaboration networks of company directors , movie actors , and scientists . As we will see, the random graph works well as a model of these networks, giving good order-of-magnitude estimates of all quantities investigated, and in some cases giving results of startling accuracy.
Our first example is the collaboration network of the members of the boards of directors of the Fortune 1000 companies (the one thousand US companies with the highest revenues). The data come from the 1999 Fortune 1000 and in fact include only 914 of the 1000, since data on the boards of the remaining 86 were not available. The data form a bipartite graph in which one type of vertex represents the boards of directors, and the other type the members of those boards, with edges connecting boards to their members. In Fig. 8 we show the frequency distribution of the numbers of boards on which each member sits, and the numbers of members of each board. As we see, the former distribution is close to exponential, with the majority of directors sitting on only one board, while the latter is strongly peaked around 10, indicating that most boards have about 10 members.
Using these distributions, we can define generating functions $`f_0(x)`$ and $`g_0(x)`$ as in Eq. (28), and hence find the generating functions $`G_0(x)`$ and $`G_1(x)`$ for the distributions of numbers of co-workers of the directors. We have used these generating functions and Eqs. (81) and (90) to calculate the expected clustering coefficient $`C`$ and the average number of co-workers $`z`$ in the one-mode projection of board directors on a random bipartite graph with the same vertex degree distributions as the original dataset. In Table I we show the results of these calculations, along with the same quantities for the real Fortune 1000. As the table shows the two are in remarkable—almost perfect—agreement.
It is not just the average value of $`z`$ that we can calculate from our generating functions, but the entire distribution: since the generating functions are finite polynomials in this case, we can simply perform the derivatives to get the probability distribution $`r_z`$. In Fig. 9, we show the results of this calculation for the Fortune 1000 graph. The points in the figure show the actual distribution of $`z`$ for the real-world data, while the solid line shows the theoretical results. Again the agreement is excellent. The dashed line in the figure shows the distribution for an ordinary Poisson random graph with the same mean. Clearly this is a significantly inferior fit.
In fact, within the business world, attention has focussed not on the collaboration patterns of company directors, but on the “interlocks” between boards, i.e., on the one-mode network in which vertices represent boards of directors and two boards are connected if they have one or more directors in common . This is also simple to study with our model. In Fig. 10 we show the distribution of the numbers of interlocks that each board has, along with the theoretical prediction from our model. As we see, the agreement between empirical data and theory is significantly worse in this case than for the distribution of co-directors. In particular, it appears that our theory significantly underestimates the number of boards which are interlocked with very small or very large numbers of other boards, while over estimating those with intermediate numbers of interlocks. One possible explanation of this is that “big-shots work with other big-shots.” That is, the people who sit on many boards tend to sit on those boards with other people who sit on many boards. And conversely the people who sit on only one board (which is the majority of all directors), tend to do so with others who sit on only one board. This would tend to stretch the distribution of numbers of interlocks, just as seen in figure, producing a disproportionately high number of boards with very many or very few interlocks to others. To test this hypothesis, we have calculated, as a function of the number of boards on which a director sits, the average number of boards on which each of their codirectors sit. The results are shown in the inset of Fig. 10. If these two quantities were uncorrelated, the plot would be flat. Instead, however, it slopes clearly upwards, indicating indeed that on the average the big-shots work with other big-shots. (This idea is not new. It has been discussed previously by a number of others—see Refs. and , for example.)
The example of the boards of directors is a particularly instructive one. What it illustrates is that the cases in which our random graph models agree well with real-world phenomena are not necessarily the most interesting. Certainly it is satisfying, as in Fig. 9, to have the theory agree well with the data. But probably Fig. (10) is more instructive: we have learned something about the structure of the network of the boards of directors by observing the way in which the pattern of board interlocks differs from the predictions of the purely random network. Thus it is perhaps best to regard our random graph as a null model—a baseline from which our expectations about network structure should be measured. It is deviation from the random graph behavior, not agreement with it, that allows us to draw conclusions about real-world networks.
We now look at three other graphs for which our theory also works well, although again there are some noticeable deviations from the random graph predictions, indicating the presence of social or other phenomena at work in the network.
We consider the graph of movie actors and the movies in which they appear and graphs of scientists and the papers they write in physics and biomedical research . In Table I we show results for the clustering coefficients and average coordination numbers of the one-mode projections of these graphs onto the actors or scientists. As the table shows, our theory gives results for these figures which are of the right general order of magnitude, but typically deviate from the empirically measured figures by a factor of two or so. In the insets of Fig. 9 we show the distributions of numbers of collaborators in the movie actor and physics graphs, and again the match between theory and real data is good, but not as good as with the Fortune 1000.
The figures for clustering and mean numbers of collaborators are particularly revealing. The former is uniformly about twice as high in real life as our model predicts for the actor and scientist networks. This shows that there is a significant tendency to clustering in these networks, in addition to the trivial clustering one expects on account of the bipartite structure. This may indicate, for example, that scientists tend to introduce pairs of their collaborators to one another, thereby encouraging clusters of collaboration. The figures for average numbers of collaborators show less deviation from theory than the clustering coefficients, but nonetheless there is a clear tendency for the numbers of collaborators to be smaller in the real-world data than in the models. This probably indicates that scientists and actors collaborate repeatedly with the same people, thereby reducing their total number of collaborators below the number that would naively be expected if we consider only the numbers of papers that they write or movies they appear in. It would certainly be possible to take effects such as these into account in a more sophisticated model of collaboration practices.
### B The world-wide web
In this section we consider the application of our theory of random directed graphs to the modeling of the world-wide web. As we pointed out in Section III A, it is not at present possible to make a very accurate random-graph model of the web, because to do so we need to know the joint distribution $`p_{jk}`$ of in- and out-degrees of vertices, which has not to our knowledge been measured. However, we can make a simple model of the web by assuming in- and out-degree to be independently distributed according to their known distributions. Equivalently, we assume that the joint probability distribution factors according to $`p_{jk}=p_jq_k`$.
Broder et al. give results showing that the in- and out-degree distributions of the web are approximately power-law in form with exponents $`\tau _{\mathrm{in}}=2.1`$ and $`\tau _{\mathrm{out}}=2.7`$, although there is some deviation from the perfect power law for small degree. In Fig. 11, we show histograms of their data with bins chosen to be of uniform width on the logarithmic scales used. (This avoids certain systematic errors known to afflict linearly histogrammed data plotted on log scales.) We find both distributions to be well-fitted by the form
$$p_k=C(k+k_0)^\tau ,$$
(100)
where the constant $`C`$ is fixed by the requirement of normalization, taking the value $`1/\zeta (\tau ,k_0)`$, were $`\zeta (x,y)`$ is the generalized $`\zeta `$-function . The constants $`k_0`$ and $`\tau `$ are found by least-squares fits, giving values of $`0.58`$ and $`3.94`$ for $`k_0`$, and $`2.17`$ and $`2.69`$ for $`\tau `$, for the in- and out-degree distributions respectively, in reasonable agreement with the fits performed by Broder et al. With these choices, the data and Eq. (100) match closely (see Fig. 11 again).
Neither the raw data nor our fits to them satisfy the constraint (68), that the total number of links leaving pages should equal the total number arriving at them. This is because the data set is not a complete picture of the web. Only about 200 million of the web’s one billion or so pages were included in the study. Within this subset, our estimate of the distribution of out-degree is presumably quite accurate, but many of the outgoing links will not connect to other pages within the subset studied. At the same time, no incoming links which originate outside the subset of pages studied are included, because the data are derived from “crawls” in which web pages are found by following links from one to another. In such a crawl one only finds links by finding the pages that they originate from. Thus our data for the incoming links is quite incomplete, and we would expect the total number of incoming links in the dataset to fall short of the number of outgoing ones. This indeed is what we see. The totals for incoming and outgoing links are approximately $`2.3\times 10^8`$ and $`1.1\times 10^9`$.
The incompleteness of the data for incoming links limits the information we can at present extract from a random graph model of the web. There are however some calculations which only depend on the out-degree distribution.
Given Eq. (100), the generating functions for the out-degree distribution take the form
$$G_0(x)=G_1(x)=\frac{\mathrm{\Phi }(x,\tau ,k_0)}{\zeta (\tau ,k_0)},$$
(101)
where $`\mathrm{\Phi }(x,y,z)`$ is the Lerch $`\mathrm{\Phi }`$-function . The corresponding generating functions $`F_0`$ and $`F_1`$ we cannot calculate accurately because of the incompleteness of the data. The equality $`G_0=G_1`$ (and also $`F_0=F_1`$) is a general property of all directed graphs for which $`p_{jk}=p_jq_k`$ as above. It arises because in such graphs in- and out-degree are uncorrelated, and therefore the distribution of the out-degree of a vertex does not depend on whether you arrived at it by choosing a vertex at random, or by following a randomly chosen edge.
One property of the web which we can estimate from the generating functions for out-degree alone is the fraction $`S_{\mathrm{in}}`$ of the graph taken up by the giant strongly connected component plus those sites from which the giant strongly connected component can be reached. This is given by
$$S_{\mathrm{in}}=1G_0(1S_{\mathrm{in}}).$$
(102)
In other words, $`1S_{\mathrm{in}}`$ is a fixed point of $`G_0(x)`$. Using the measured values of $`k_0`$ and $`\tau `$, we find by numerical iteration that that $`S_{\mathrm{in}}=0.527`$, or about 53%. The direct measurements of the web made by Broder et al. show that in fact about 49% of the web falls in $`S_{\mathrm{in}}`$, in reasonable agreement with our calculation. Possibly this implies that the structure of the web is close to that of a directed random graph with a power-law degree distribution, though it is possible also that it is merely coincidence. Other comparisons between random graph models and the web will have to wait until we have more accurate data on the joint distribution $`p_{jk}`$ of in- and out-degree.
## VI Conclusions
In this paper we have studied in detail the theory of random graphs with arbitrary distributions of vertex degree, including directed and bipartite graphs. We have shown how, using the mathematics of generating functions, one can calculate exactly many of the statistical properties of such graphs in the limit of large numbers of vertices. Among other things, we have given explicit formulas for the position of the phase transition at which a giant component forms, the size of the giant component, the average and distribution of the sizes of the other components, the average numbers of vertices a certain distance from a given vertex, the clustering coefficient, and the typical vertex–vertex distance on a graph. We have given examples of the application of our theory to the modeling of collaboration graphs, which are inherently bipartite, and the World-Wide web, which is directed. We have shown that the random graph theory gives good order-of-magnitude estimates of the properties of known collaboration graphs of business-people, scientists and movie actors, although there are measurable differences between theory and data which point to the presence of interesting sociological effects in these networks. For the web we are limited in what calculations we can perform because of the lack of appropriate data to determine the generating functions. However, the calculations we can perform agree well with empirical results, offering some hope that the theory will prove useful once more complete data become available.
## Acknowledgements
The authors would like to thank Lada Adamic, Andrei Broder, Jon Kleinberg, and Cris Moore for useful comments and suggestions, and Jerry Davis, Paul Ginsparg, Oleg Khovayko, David Lipman, Grigoriy Starchenko, and Janet Wiener for supplying data used in this study. This work was funded in part by the National Science Foundation, the Army Research Office, the Electric Power Research Institute, and Intel Corporation.
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# Charge conservation and time-varying speed of light
## Abstract
It has been recently claimed that cosmologies with time dependent speed of light might solve some of the problems of the standard cosmological scenario, as well as inflationary scenarios. In this letter we show that most of these models, when analyzed in a consistent way, lead to large violations of charge conservation. Thus, they are severly constrained by experiment, including those where $`c`$ is a power of the scale factor and those whose source term is the trace of the energy-momentum tensor. In addition, early Universe scenarios with a sudden change of $`c`$ related to baryogenesis are discarded.
Since one of the key hypothesis of special relativity is the frame independence of the velocity of light $`c`$, it is implied in this statement the time and space independence of this velocity. As well established that it may seem, this constancy principle has been recently contested to provide an alternative account of the horizon, flatness and cosmological constant problems present in the standard big bang scenario. Instead of the superluminal expansion of the Universe present in inflationary scenarios, a period in which light traveled much faster than today would explain the homogeneity we see today in the Universe. Some cosmological models have also been analyzed afterwards to test the dynamical viability of this scenario.
These ideas are highly provocative, not only from the observational viewpoint but also from the conceptual one. Indeed one of the key aspects of Einstein equivalence principle is the time-independence of the so called “fundamental constants” of physics . The replacement of these parameters by one or more dynamical fields can lead to time- as well as space-dependent local fundamental constants. Unification schemes such as superstring theories and Kaluza-Klein theories have cosmological solutions in which the low-energy fundamental constants are functions of time. Usually low-energy phenomena are used to analize the variation rate of the fundamental constants -. If the cosmological dynamics of a field is such that its large-scale value is invariant under local Lorentz transformations, or if the local coupling with matter is strong enough so that it depends on the local environment (*e.g.* electromagnetism with the absorber condition), then the local field equations will be Lorentz invariant. If on the other hand the cosmological evolution is non-trivial and the field couples softly with the local matter, it will act as an external bath, breaking local Lorentz invariance. A variable speed of light theory may belong to the latter set of theories. Any VSL theory poses an additional problem, namely that $`c`$ is a dimensional constant, and talking about a varying dimensional constant is not an invariant statement: we can change our units and obtain a different time dependence of such a parameter. Of course, once we fix our units, every claim about a dimensional parameter is an invariant claim, since we are implicitly referring to a dimensionless ratio: that between the parameter and the unit .
Any scientific theory has to be stated in clear and precise terms. Beckenstein’s theory of a variable fine structure constant was based on Lorentz invariance, explicitly protecting charge conservation. In the case of VSL theories, local Lorentz invariance is relaxed, and the inhomogeneous Maxwell equations are assumed to be
$$\frac{1}{c}_\mu \left(cF^{\mu \nu }\right)=4\pi j^\nu ,$$
(1)
where $`j^\mu =(\rho ,𝐣/c)`$ is the electric charge current. In reference it was suggested that charge is conserved, implying a variation of the fine structure constant $`\alpha =\frac{e^2}{\mathrm{}c}`$. The constancy of $`e`$ can be derived, for instance, from Dirac’s equation, written in Hamiltonian form:
$$i\mathrm{}_t\psi =(i\mathrm{}𝜶\mathbf{}+\mathrm{h}.\mathrm{c}.)\psi +mc^2\psi $$
which implies that $`Q=e_V\psi ^{}\psi d^3x`$ is conserved. The above form is, however, not unique since powers of $`c`$ can be introduced in the equation in several ways, followed by appropriate symmetrization.
However, from equation (1) it can be easily seen that $`j^\mu `$ is no longer a conserved current, but satisfies the equation:
$$4\pi _\mu (cj^\mu )=_\mu _\nu (cF^{\mu \nu })$$
(2)
The right hand side is not null, because the partial derivatives do not commute:
$$[_0,_i]=\frac{c_{,i}}{c^2}_t$$
It is usually assumed that $`c`$ is a function of a scalar field and that it depends only on time in the comoving cosmological frame. In the local frame of the solar system, moving with a velocity $`v`$ with respect to the cosmological frame, a small space dependence will arise, with gradients $`O(v/c)`$ with respect to the time derivative, which can be neglected for the present purposes. So, the right hand side of equation 2 is effectively zero. The left hand side, however, is not a four divergence, because $`_{x^0}=1/c(t)_t`$. The fully expanded expression is:
$$_t\rho +\frac{\dot{c}}{c}\rho +𝐣=0$$
(3)
If equation (3) is integrated over a volume $`V`$ containing the charges, we obtain
$$\frac{\dot{Q}}{Q}=\frac{\dot{c}}{c}$$
(4)
where $`Q=_V\rho d^3x`$ is the total electric charge. This is our main result.
Equation (4) provides very stringent tests of the variation of $`c`$, since there have been many experiments to test the conservation of charge . Depending on the details of the theory, several cases arise.
If we assume, as it is usually done in this context that the electron charge $`e`$ is constant, charge conservation can only be broken by processes that change charge discontinuously, such as the dissapearance of electrons or the transformation of neutrons into protons. A generic model for these processes has been proposed in references , under the assumption of total energy conservation. The first three entries of Table 1 show some sample upper limits obtained from these processes with different techniques and hypothesis. These upper limits on $`\dot{c}/c`$ are much smaller than those obtained by a direct measurement, for instance in reference , namely $`\dot{c}/c<10^{13}\mathrm{yr}^1`$.
On the other hand, if $`e`$ varies continuously in such a way that $`ce`$ is conserved, then $`\alpha c^3`$ and strict limits can be obtained from geophysical or astronomical data, such as the Oklo phenomenon or the line spectra of distant quasars . Furthermore, evidence for the time variation of the fine structure constant has been claimed by Webb et al. and confirmed while we were correcting this paper . Thus, if the systematic errors are well estimated, the requirement of conservation of charge suggests that, the mechanism for $`\alpha `$ variation should not be driven by the change in the speed of light.
These limits discard a great number of cosmological models with varying velocity of light. For instance, the family introduced in parameterizes light velocity in the form: $`c=c_0\left(\frac{a}{a_0}\right)^n`$ with $`a`$ the cosmological scale factor. Even though, was not proposing that this behaviour exists at all times, suggests that this solutions could be extended to radiation and matter-dominated universes. It was shown in reference that $`n<1/2`$ is necessary to solve the flatness and horizon problems, and $`n<3/2`$ solves the cosmological constant problem.
On the other hand, any complete and consistent VSL theory will predict dynamically the value of c via a wave-like equation for $`\psi =\mathrm{ln}\frac{c}{c_0}`$, whose source term is proportional to the trace of the energy-momentum tensor $`T`$. Thus, in the neighborhood of a quasi-static system, such as a star or a virialized galaxy cluster, a generic expresion for a varying $`c`$ will be :
$$c=c_c(t)\left(1\frac{\lambda GM}{c_0^2r}\right)$$
(5)
Here $`\lambda `$ is a constant that depends on the specific VSL theory. Furthermore, in the limit $`r\mathrm{}`$ , the expression for $`c`$ reduces to the cosmological one ($`cc_c(t)`$).
For the purposes of this paper, our interest is focused on violation of charge conservation. Expanding equation(2) we obtain:
$$4\pi _\mu j^\mu =4\pi \frac{_\mu c}{c}j^\mu _tc\frac{\stackrel{}{}c.\stackrel{}{E}}{c^3}+\frac{\stackrel{}{}.c}{c}_0\stackrel{}{E}$$
(6)
In a static situation $`\stackrel{}{j}=0`$ and $`_0E=0`$, thus the effects of $`\stackrel{}{}c`$ are of second order and can be neglected. Hence, there are two contributions to the violation of charge conservation, one accounts for time-variations over cosmological time-scales and the other for the motion of the Earth with respect to massive bodies such as clusters or galaxies:
$$\frac{\dot{c}}{c}=\left(\frac{\dot{c}}{c}\right)_{cosmological}+\left(\frac{\dot{c}}{c}\right)_{local}=n\frac{\dot{a}}{a}+\frac{\lambda GM}{c_0^2r^3}\stackrel{}{r}.\stackrel{}{v}$$
(7)
where $`\stackrel{}{v}`$ is the velocity of the Earth with respect to the lump. ($`\stackrel{}{v}1000\frac{km}{seg}`$ for the Virgo supercluster).
From the first term of equation (7) we find the limits of table 1 on $`n`$, which contradict the above requirements. (We use $`H_0=65\mathrm{km}/\mathrm{s}/\mathrm{Mpc}=6.65\times 10^{11}\mathrm{yr}^1`$). Similar bounds can be obtained for other similar models, such as those studied in reference (See also ).
For the Virgo supercluster, $`\frac{GM}{c_0^2r}2\times 10^7`$. The last three results in table 1 for $`\lambda `$ are obtained from the second term of equation (7) . These results rule out essentially any VSL whose source of variation is the trace of the energy-momentum tensor $`T`$.
Finally, models similar to the original Albrecht-Magueijo one , involving a sudden change of $`c`$ between two different constant values in the very early Universe, are not affected by the above limits. Orito and Yoshimura observed that if charge conservation is broken in the very early Universe, a large charge excess should have been formed through a mechanism similar to that of baryogenesis: violation of $`Q`$, $`C`$ and $`CP`$ conservation while the system is out of thermodynamic equilibrium . In the above mentioned models, the net charge excess will be produced by way-out-of-equilibrium production and decay of heavy mesons .
Let $`X`$ be an unstable meson that produces a mean baryon number $`ϵ_B`$ and a mean charge excess $`ϵ_Q`$ per decay, and assume that matter is created during the charge transition period. Then, the equations for the evolution of the number densities of $`X,B,Q`$ will be :
$`(a^3n_X)^{}+\lambda _X(a^3n_X)`$ $`=`$ $`{\displaystyle \frac{3Kc\dot{c}}{4\pi Gm_X}}a`$
$`\dot{(a^3n_B)}^{}`$ $`=`$ $`ϵ_B\lambda _Xn_Xa^3`$
$`\dot{(a^3n_Q)}^{}`$ $`=`$ $`ϵ_Q\lambda _Xn_Xa^3`$
In these equations, $`K`$ is the curvature parameter, $`m_X`$ is the mass of the X meson, and all densities are effective densities (particle - antiparticle). Let us assume that the change in $`c`$ occurs in a short interval of time $`t_c<<\tau <<\frac{1}{\lambda _x}`$. Charge conservation will be broken only during this interval, but the $`X`$ meson decay will always produce a baryon excess. With these hypothesis, the above equations have the following solutions:
$`a^3n_X`$ $`{\displaystyle \frac{3K(c_0^2c_P^2)}{8\pi Gm_X}}a(0)e^{\lambda t}`$ $`=a^3(0)n_X^0e^{\lambda t}`$
$`a^3n_B`$ $`ϵ_Ba^3(0)n_X^0\left(1e^{\lambda t}\right)`$ $`ϵ_Ba^3(0)n_X^0`$
$`a^3n_Q`$ $`ϵ_Q\lambda _X\tau a^3(0)n_X^0`$ $`{\displaystyle \frac{ϵ_Q}{ϵ_B}}\lambda _X\tau a^3n_B`$
After the transition, $`n_Q`$ will be fixed but $`n_B`$ will be diluted from the above estimate by thermal processes . So we finally get a lower bound on the charge excess:
$$\left|\frac{n_Q}{n_B}\right|\frac{ϵ_Q}{ϵ_B}\lambda _X\tau $$
(8)
As we have mentioned before, we expect on general grounds that $`\tau >t_{Pl}`$, while $`ϵ_Qϵ_B`$, since these fractions depend both on the $`C`$ and $`CP`$ breaking terms in the lagrangean. Thus, equation (8) predicts a firm lower limit for the charge excess. Orito and Yoshimura have given limits on any charge excess in the Universe, shown in table 1. These limits are many orders of magnitude below the prediction of equation (8). The last column of the table shows rough estimates of $`\tau `$ taken from the observational limits, assuming $`1/\lambda _Xt_{GUT}`$.
Although these results do not rule out all varying velocity of light theories, they put very stringent bounds on them through the conservation of charge requirement. Moreover, similar bounds will hold for any theory with varying speed of light velocity in the early universe. These bounds, which may be lowered through improvements in the experimental techniques , will lead into deeper understanding of these interesting theories.
The authors want to thank Diego F. Torres for many interesting discussions and advice. The authors acknowledge partial economic support through the project 011/G035, Universidad Nacional de La Plata.
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# A realistic interpretation of the density matrix II: The non-relativistic case
## 1 Introduction
In a recent paper we proposed a realistic interpretation of the Schrödinger equation for density matrices, in which the difference between the position arguments of the density matrix is considered as an objective extra space dimension. In the case of a free particle, where the potential $`V\left(x\right)`$ vanishes, we found solutions which are perfectly localized both in position space and in momentum space; these solutions behave exactly as non-relativistic point-like particles moving at constant speed, with the correct values for all observable quantities. In the general case, where $`V\left(x\right)0`$, we pointed out that the natural frequencies of the stationary states in the density matrix representation correspond to the difference between two energy levels in the original quantum system; since the “jumps” between energy levels are observable, while the individual energy levels are not, we trivially deduced that the observable natural frequencies are the same in both representations.
In the first paper the non-relativistic case was treated mainly as an introduction to the relativistic case; in this second paper we will study it in more detail, examining the correspondence between our new representation and the standard representation of non-relativistic quantum mechanics, based on the Schrödinger equation for pure states. In Section 2 we will first summarize the main results of the previous paper, and then we will give a qualitative definition of a quantum trajectory (describing an individual physical system), as opposed to a statistical mixture (describing an ensemble of physical systems); as a consequence, we will reject two basic principles of standard quantum mechanics, namely the superposition principle and the wave function collapse. In Section 3 we will show the results obtained by applying our approach to simple potentials: by examining the linear harmonic oscillator we will shed a new light on the supposed equivalence between energy eigenvectors and stationary states; besides, the delta barrier potential will provide a semi-classical explanation of the tunnel effect, in which an essential feature is the introduction of internal degrees of freedom even for spinless particles. Finally in Section 4 we will present our conclusions.
## 2 The density matrix representation and the superposition principle
The standard quantum description of a non-relativistic particle moving in one space dimension is based on a wave function $`\psi \left(x\right)`$, whose time evolution is given by the Schrödinger equation:
$$i\mathrm{}\frac{\psi }{t}=\text{H}\psi =\frac{\mathrm{}^2}{2m}\frac{^2\psi }{x^2}+V\left(x\right)\psi $$
(1)
We will refer to equation (1) as to the Schrödinger equation for pure states. The time evolution generated by (1) belongs to the group of unitary transformations; a new representation of this group, the density matrix representation, is obtained by means of the fundamental relation
$$\phi (x,y)=\psi \left(x\right)\psi ^{}\left(y\right)$$
(2)
If we apply an arbitrary unitary transformation to the pure state $`\psi \left(x\right)`$, the relation (2) enables us to obtain the same transformation in the density matrix representation. Specifically, if we consider the hamiltonian operator H, the momentum operator P and the position operator Q as generators respectively of time translations, space translations and momentum translations, we easily obtain the following expressions for the same generators in the new representation:
$`\text{H}\phi `$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}\left({\displaystyle \frac{^2\phi }{x^2}}{\displaystyle \frac{^2\phi }{y^2}}\right)+\left[V\left(x\right)V\left(y\right)\right]\phi `$
$`\text{Q}\phi `$ $`=`$ $`\left(xy\right)\phi `$ (3)
$`\text{P}\phi `$ $`=`$ $`i\mathrm{}\left({\displaystyle \frac{\phi }{x}}+{\displaystyle \frac{\phi }{y}}\right)`$
The time evolution in the density matrix representation is then simply
$$i\mathrm{}\frac{\phi }{t}=\frac{\mathrm{}^2}{2m}\left(\frac{^2\phi }{x^2}\frac{^2\phi }{y^2}\right)+\left[V\left(x\right)V\left(y\right)\right]\phi $$
(4)
In the usual interpretation, equation (4) describes the time evolution of mixed states, and therefore its solutions have only a statistical meaning: they are simply useful tools for computing probability amplitudes and expectation values. On the contrary, in our approach the solutions of equation (4) are considered as real objective fields describing individual physical systems. To avoid misunderstandings, in the rest of this paper we will use the expression “density matrix” when we mean the usual statistical interpretation and we will use the new expression “quantum matrix” when we mean a real objective field.
A quantum matrix $`\phi (x,y)`$ depends upon two position coordinates $`x`$ and $`y`$; if we want to consider $`\phi (x,y)`$ as a real objective field then we must introduce an objective extra space dimension. Therefore, we will define a new pair of position coordinates
$$x_S=\frac{1}{2}\left(x+y\right)x_D=xy$$
(5)
and we will interpret $`x_S`$ as the “physical” position coordinate, while $`x_D`$ will be an “auxiliary” position coordinate, having observable effects only around the point $`x_D=0`$; both $`x_S`$ and $`x_D`$, however, will be considered as objective position coordinates (note that in the definition of $`x_D`$ we have inverted the sign convention with respect to our first paper ). In the rest of this paper we will write loosely $`\phi (x,y)`$ or $`\phi (x_S,x_D)`$, meaning the same field expressed in two different coordinate systems.
The transition from position space to momentum space is obtained by means of the two-dimensional Fourier transform:
$$\mathrm{\Phi }(k_x,k_y)=\frac{1}{2\pi }\phi (x,y)e^{ik_xx+ik_yy}dxdy$$
(6)
Again we define a new pair of momentum coordinates:
$$k_S=\frac{1}{2}\left(k_x+k_y\right)k_D=k_xk_y$$
(7)
and again $`k_S`$ will be the “physical” momentum coordinate, while $`k_D`$ will be an “auxiliary” momentum coordinate, having observable effects only around the point $`k_D=0`$.
Since we have two position coordinates and two momentum coordinates, we naturally can define two position operators and two momentum operators. The operators already defined in (3) may be rewritten as:
$$\text{Q}_D\phi =x_D\phi \text{P}_D\phi =i\mathrm{}\frac{\phi }{x_S}$$
(8)
In addition, we introduce the new operators $`\text{Q}_S`$ and $`\text{P}_S`$ by means of the following definition:
$$\text{Q}_S\phi =x_S\phi \text{P}_S\phi =i\mathrm{}\frac{\phi }{x_D}$$
(9)
In momentum space the operators $`\text{P}_D`$ and $`\text{P}_S`$ are simply given by:
$$\text{P}_D\mathrm{\Phi }=\mathrm{}k_D\mathrm{\Phi }\text{P}_S\mathrm{\Phi }=\mathrm{}k_S\mathrm{\Phi }$$
(10)
It is very important here to note that $`\text{Q}_S`$ and $`\text{P}_S`$ are commuting operators, i.e. $`[\text{Q}_S,\text{P}_S]=0`$; this property is fundamental for our approach, because it means that they have common eigenvectors. These eigenvectors have the form
$$\phi =\delta \left(x_Sx_0\right)e^{ik_0x_D}$$
(11)
or, in momentum space,
$$\mathrm{\Phi }=\delta \left(k_Sk_0\right)e^{ik_Dx_0}$$
(12)
Since $`x_S`$ and $`k_S`$ are the “physical” position and momentum coordinates, it is clear that the field defined in (11) and (12) is perfectly localized in position space at $`x_S=x_0`$ and in momentum space at $`p_S=p_0=\mathrm{}k_0`$. Therefore in the density matrix representation the Heisenberg uncertainty principle is not true and we conclude, as in our previous paper, that the Heisenberg principle has no ontological meaning; rather, it is just a shortcoming of the standard quantum formalism based on pure states. A similar point of view is expressed by Olavo .
While the two operators $`\text{Q}_D`$ and $`\text{P}_D`$ are interesting just because they generate translations respectively in momentum space and in position space, the two operators $`\text{Q}_S`$ and $`\text{P}_S`$ are more strictly related to the measurement of the corresponding physical quantities, as we will now see. To define the observable quantities in the density matrix representation, we start from the expressions for the mean values in the original representation and then exploit the fundamental relation (2); we obtain:
$`Q`$ $`=`$ $`{\displaystyle \left(\text{Q}_S\phi \right)}|_{x_D=0}\text{d}x_S`$ (13)
$`P`$ $`=`$ $`{\displaystyle \left(\text{P}_S\phi \right)}|_{x_D=0}\text{d}x_S`$ (14)
$`E`$ $`=`$ $`{\displaystyle [\frac{1}{2m}\text{P}_S^2+V\left(\text{Q}_S\right)]\phi }|_{x_D=0}\text{d}x_S`$ (15)
The above definitions confirm our interpretation of $`x_D`$ as an “auxiliary” position coordinate, having observable effects only around $`x_D=0`$. The same quantities may be written in momentum space, obtaining analogous expressions; for instance, the momentum $`P`$ may be written as
$$P=\left(\text{P}_S\mathrm{\Phi }\right)|_{k_D=0}\text{d}k_S$$
(16)
Since we do not assign a statistical meaning to the quantum matrix $`\phi `$, we will not interpret the expressions (13)-(15) as mean values. On the contrary, we will consider (13) as the center of mass of the system, while (14) and (15) will be respectively the total momentum and the total energy.
If we Fourier transform the quantum matrix $`\phi `$ with respect to $`x_D`$ alone, we obtain the well known Moyal-Wigner transformation:
$$F(x_S,p_S)=\frac{1}{2\pi \mathrm{}}\phi (x_S,x_D)e^{i\frac{p_S}{\mathrm{}}x_D}dx_D$$
(17)
first introduced by Wigner in 1932. The Wigner function $`F(x_S,p_S)`$ has some valuable properties: the common eigenvectors of the operators $`\text{Q}_S`$ and $`\text{P}_S`$ are represented simply by products of delta functions, i.e. $`F(x_S,p_S)=\delta \left(x_Sx_0\right)\delta \left(p_Sp_0\right)`$; besides, those operators which depend only on $`\text{Q}_S`$ and $`\text{P}_S`$ become c-numbers in the Wigner representation; finally, if $`G`$ is an observable quantity depending on position and momentum through the function $`g(x,p)=g_1\left(x\right)+g_2\left(p\right)`$, then in the Wigner representation we can write the expression
$$G=g(x_S,p_S)F(x_S,p_S)dx_Sdp_S$$
(18)
i.e. $`F(x_S,p_S)`$ weighs the function $`g(x_S,p_S)`$ over the whole phase space $`(x_S,p_S)`$; the application of (18) to the quantities defined in (13)-(15) is straightforward. Thus the Wigner function is a useful tool for establishing relations between the quantum description and the classical description based on phase space; indeed, it provides the basis for Moyal’s deformation quantization , which is an autonomous formulation of quantum mechanics, alternative to the more familiar Hilbert space and path integral quantizations.
However, in the usual statistical interpretation of the density matrix, expression (18) yields the mean value of the quantity $`G`$, and therefore the Wigner function seems to play the role of a joint probability distribution in phase space. But even if the function $`\phi (x,y)`$ satisfies all the properties required for being a well defined density matrix, it is in any case possible for the Wigner function to be negative in some phase space region. This is the well known “negativity problem”, which prevents a complete analogy between the classical description in phase space and the quantum description in the Wigner representation. On the contrary, if we consider the quantum matrix $`\phi (x,y)`$ as a real objective field, describing the state of an individual physical system, we do not assign to the Wigner function any statistical meaning, and therefore it may safely have negative values: in our interpretation, the “negativity problem” is no problem at all.
Now we want to study the difference between the time evolution of a quantum system in the density matrix representation and the time evolution of the corresponding classical system in phase space. We will use the Wigner representation, but the results may be easily carried to position space and momentum space, since we will express them in operator form. The quantum evolution, obtained from (4), is simply given by:
$$i\mathrm{}\frac{F_Q}{t}=H_QF_Q=\left[\frac{1}{m}\text{P}_S\text{P}_D+V(\text{Q}_S+\frac{1}{2}\text{Q}_D)V(\text{Q}_S\frac{1}{2}\text{Q}_D)\right]F_Q$$
(19)
To obtain the corresponding classical evolution we start from the continuity equation for the classical joint probability density in phase space:
$$\frac{F_C}{t}+\frac{F_C}{x}\frac{p}{m}\frac{F_C}{p}V^{}=0$$
(20)
where $`V^{}`$ is the space derivative of the potential $`V`$. Translating (20) to operator form we then obtain
$$i\mathrm{}\frac{F_C}{t}=H_CF_C=\left[\frac{1}{m}\text{P}_S\text{P}_D+V^{}\left(\text{Q}_S\right)\text{Q}_D\right]F_C$$
(21)
Expressions (19) and (21) are quite similar; indeed, if we expand $`V\left(x\right)`$ in powers of $`x`$, we can write:
$$H_QH_C=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{\left(2n+1\right)!\mathrm{\hspace{0.17em}2}^{2n}}V^{(2n+1)}\left(\text{Q}_S\right)\text{Q}_D^{(2n+1)}$$
(22)
The difference between the quantum evolution and the classical evolution involves the space derivatives $`V^{(2n+1)}`$ with $`n1`$. Therefore the quantum evolution and the classical evolution are equivalent if $`V=0`$ (free particle), $`V=kx`$ (particle moving in a uniform force field) and $`V=kx^2`$ (linear harmonic oscillator). This is a rather surprising result since we know that the quantum energy spectrum of the linear harmonic oscillator is discrete, and this is a highly non-classical feature; we will examine this point in detail in the next section.
Now we turn to the problem of finding localized solutions in the density matrix representation; the existence of such solutions is fundamental for our approach. It is well known that to each density matrix satisfying the Schrödinger equation (4) we can associate a quasi-classical trajectory, namely the time evolution of the mean values $`x\left(t\right)`$ and $`p\left(t\right)`$ of position and momentum; this is of course a consequence of the Ehrenfest theorem:
$$\dot{x}=\frac{1}{m}p\dot{p}=V^{}\left(x\right)$$
(23)
From the trivial relation
$$V^{}\left(x\right)=V^{}\left(x\right)+\underset{n=2}{\overset{\mathrm{}}{}}\frac{1}{n!}V^{\left(n+1\right)}\left(x\right)\left(xx\right)^n$$
(24)
it then follows that the equations (23) define a classical trajectory only if the derivatives $`V^{\left(n\right)}`$ vanish for $`n>2`$ (again the free particle and the linear harmonic oscillator), while in the general case the trajectory will be quasi-classical if the moments $`\left(xx\right)^n`$ are small, i.e. if the wave-packet is well localized in position space around its mean value $`x\left(t\right)`$.
In standard quantum mechanics, due to the Heisenberg uncertainty principle, wave-packets normally spread out with time; only in some special cases this does not happen, the most famous example being the coherent states of the linear harmonic oscillator first introduced by Schrödinger in 1926 . However, in the density matrix representation the Heisenberg uncertainty principle is not true, and therefore it should be always possible to find solutions which do not spread out with time and remain well localized around their center of mass $`x\left(t\right)`$, in the sense that their “matter density” $`\phi (x_S,x_D,t)|_{x_D=0}`$ is significantly different from zero only in a small region around $`x_S=x\left(t\right)`$. In our approach, the existence of such solutions is of the utmost importance and allows us to state the following “localization postulate”: the only physical solutions of the Schrödinger equation (4) are those which remain well localized around their center of mass $`x\left(t\right)`$ in the limit $`t\pm \mathrm{}`$; by “physical” we mean “representing individual physical systems”, as opposed to statistical mixtures or ensembles. In the rest of this paper, we will use the definition “quantum trajectories” to label such solutions: thus a quantum trajectory is the time evolution of a quantum matrix, as defined at the beginning of the present section.
Clearly, our localization postulate is expressed as a qualitative statement, since the concept of “well localized” wave packet is not sharply defined. However, even in its rudimentary form, this postulate implies an important difference between our approach and the standard interpretation of quantum mechanics: in standard QM, immediately before the measure of an observable $`O`$ (for instance the position), the state of the system is usually a linear superposition of eigenstates of the associated hermitian operator; then, when the measure is performed and the result $`O_n`$ is obtained, the state suddenly collapses to the eigenvector $`|O_n`$. This wave function collapse has always been an obscure feature of the standard QM interpretation: many alternative explanations have been proposed (many worlds splitting , environment induced decoherence , spontaneous localization , …) but until now none of these alternative explanations seems to be universally accepted by the academic community.
On the contrary, in our approach the physical wave-packets are always well localized around their center of mass $`x\left(t\right)`$ and therefore no wave function collapse is needed. When a measure of position is performed, and the result is a certain position $`x_0`$, this means that the wave-packet immediately before the measure was already well localized around the point $`x_0`$: our localization postulate explicitly prevents the spreading of the physical wave function $`\phi (x,y,t)`$ over a wide space region as $`t\mathrm{}`$.
Besides eliminating the need of the wave function collapse, our localization postulate implicitly negates the validity of the superposition principle: in our approach, a linear combination of physical solutions is not a physical solution any more, even if it satisfies the linear motion equation (4). This is a trivial consequence of the fact that in general an arbitrary superposition of localized wave-packets is not a localized wave-packet.
However, a linear superposition of quantum trajectories may have a statistical interpretation. Let’s consider the set $`\left\{\phi _\lambda \right\}`$ of all quantum trajectories, where $`\lambda `$ is some appropriate index (for instance, $`\lambda `$ may be the center of mass and momentum at time $`t=0`$, together with some non-classical internal state). Then the integral
$$\phi (x,y,t)=f\left(\lambda \right)\phi _\lambda (x,y,t)\text{d}\lambda $$
(25)
where $`f\left(\lambda \right)`$ is a non-negative probability distribution, represents a statistical ensemble of particles; in our approach, the solutions defined by (25) play the same role as the density matrices in standard quantum mechanics. Besides, our approach rejects those solutions of the Schrödinger equation (4) which are neither quantum trajectories nor positive superpositions of quantum trajectories: since they do not represent neither individual particles nor statistical ensembles of particles, we conclude that they are just mathematical objects with no physical meaning.
If we now apply the Moyal-Wigner transformation to the density matrix (25), we obtain the Wigner function
$$F(x_S,p_S,t)=f\left(\lambda \right)F_\lambda (x_S,p_S,t)\text{d}\lambda $$
(26)
where $`F_\lambda (x_S,p_S,t)`$ are the individual Wigner functions associated to the quantum trajectories $`\phi _\lambda (x,y,t)`$. Note that both the individual functions $`F_\lambda `$ and the average function $`F`$ may safely have negative values, since neither represents a probability density. The only probability density in (26) is the function $`f\left(\lambda \right)`$, which is non-negative by definition.
Let’s now examine the problem of finding the quantum trajectories for a specified potential $`V\left(x\right)`$. We will suppose that $`V\left(x\right)`$ vanishes as $`x\pm \mathrm{}`$: in classical terms, this means that the particle is subject to a force only in a limited space region, while outside this region it moves freely at constant speed. The first step consists in writing the Schrödinger equation (4) in the form
$$\frac{^2}{x_Sx_D}\phi =\frac{m}{\mathrm{}^2}\left[V\left(x_S+\frac{1}{2}x_D\right)V\left(x_S\frac{1}{2}x_D\right)i\mathrm{}\frac{}{t}\right]\phi $$
(27)
From (27) it is easily seen that, given two arbitrary functions $`f(x_S,t)`$ and $`g(x_D,t)`$, we may formally impose the boundary conditions
$$\phi (x_S,x_D,t)|_{x_D=0}=f(x_S,t)\text{and}\frac{\phi }{x_D}(x_S,x_D,t)|_{x_S=0}=g(x_D,t)$$
(28)
and then we may extend (28) to a complete solution $`\phi (x_S,x_D,t)`$ by means of the Taylor expansion
$$\phi (x_S,x_D,t)=f(x_S,t)+\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}\frac{^n\phi }{x_D^n}|_{x_D=0}x_D^n$$
(29)
where the derivatives are obtained by differentiating $`n`$ times (27) with respect to $`x_D`$ and then imposing $`x_D=0`$.
At first sight, this result is rather paradoxical: it implies that we may choose an arbitrary time evolution for the center of mass of our system and then obtain a solution of the Schrödinger equation which satisfies this time evolution! However, there are at least two reasons why the above procedure does not always produce acceptable results: the first reason follows from the fact that an arbitrary solution of the Schrödinger equation (4) must be a linear combination of the form
$$\phi (x,y,t)=\text{d}k_1\text{d}k_2C(k_1,k_2)f_{k1}\left(x\right)f_{k2}^{}\left(y\right)e^{\frac{i}{\mathrm{}}\left(E_{k1}E_{k2}\right)t}$$
(30)
where $`f_k\left(x\right)`$ are the energy eigenvectors of the original Schrödinger equation (1) with associated eigenvalues $`E_k`$; of course, if the energy spectrum is discrete then the integrals in (30) must be replaced by summations. Therefore the function $`f(x_S,t)`$ in (28) cannot be chosen arbitrarily, since it must be equal to (30), in the case $`x=y=x_S`$, for some choice of the coefficients $`C(k_1,k_2)`$. For instance, we know that for the linear harmonic oscillator the solutions must be periodic in time, i.e. $`\phi (x,y,t)=\phi (x,y,t+T_0)`$ where $`T_0`$ is the revolution time of the oscillator; if we choose a function $`f(x_S,t)`$ which is not periodic, then it will be impossible to obtain from it a complete solution $`\phi (x_S,x_D,t)`$ by means of (28), i.e. the Taylor series (29) will diverge.
The second reason is that the solutions obtained from (27) and (28), even if they are well defined from a mathematical point of view, will in general violate some fundamental physical principle, such as the Ehrenfest theorem (23), the energy conservation principle or the unitarity condition
$$\frac{}{t}\phi (x_S,x_D,t)|_{x_D=0}\text{d}x_S=0$$
(31)
which may be regarded as an expression of the non-relativistic principle of mass conservation. Therefore, our next step consists in stating clearly the assumptions under which the above cited principles may be derived from the Schrödinger equation; we will then require that our quantum trajectories satisfy these assumptions. The unitarity property (31) may be derived from the following condition:
$$\frac{\phi }{x_D}|_{x_D=0}\mathrm{\hspace{0.17em}0}\text{as}x_S\pm \mathrm{}$$
(32)
A sufficient condition to obtain the Ehrenfest theorem (23) is
$$x_S\frac{\phi }{x_D}|_{x_D=0}\mathrm{\hspace{0.17em}0}\text{and}\frac{^2\phi }{x_D^2}|_{x_D=0}\mathrm{\hspace{0.17em}0}\text{as}x_S\pm \mathrm{}$$
(33)
and finally the energy conservation principle follows from
$$\frac{^3\phi }{x_D^3}|_{x_D=0}\mathrm{\hspace{0.17em}0}\text{as}x_S\pm \mathrm{}$$
(34)
Thus the assumptions (32)-(34) depend only on the field $`\phi `$ and its three first derivatives with respect to $`x_D`$ computed at $`x_D=0`$. Besides, from the expressions (13)-(15) we see that the same is true for the three observable quantities, i.e. center of mass, momentum and energy (in this case only the first two derivatives are involved). If we now set
$$\phi ^{\left(0\right)}=\phi |_{x_D=0}\text{and}\phi ^{\left(n\right)}=\frac{^n\phi }{x_D^n}|_{x_D=0}n=1,2,3$$
(35)
where $`\phi ^{\left(0\right)}`$ and $`\phi ^{\left(n\right)}`$ depend only on $`x_S`$ and $`t`$, we easily derive from the Schrödinger equation (27) the following relations
$`{\displaystyle \frac{\phi ^{\left(1\right)}}{x_S}}`$ $`=`$ $`i{\displaystyle \frac{m}{\mathrm{}}}{\displaystyle \frac{\phi ^{\left(0\right)}}{t}}`$ (36)
$`{\displaystyle \frac{\phi ^{\left(2\right)}}{x_S}}`$ $`=`$ $`i{\displaystyle \frac{m}{\mathrm{}}}{\displaystyle \frac{\phi ^{\left(1\right)}}{t}}+{\displaystyle \frac{m}{\mathrm{}^2}}V^{}\left(x_S\right)\phi ^{\left(0\right)}`$ (37)
$`{\displaystyle \frac{\phi ^{\left(3\right)}}{x_S}}`$ $`=`$ $`i{\displaystyle \frac{m}{\mathrm{}}}{\displaystyle \frac{\phi ^{\left(2\right)}}{t}}+2{\displaystyle \frac{m}{\mathrm{}^2}}V^{}\left(x_S\right)\phi ^{\left(1\right)}`$ (38)
From the above considerations we deduce that we do not need to know the complete quantum trajectory $`\phi (x_S,x_D,t)`$: all properties of physical interest follow from the knowledge of the four functions $`\phi ^{\left(i\right)}(x_S,t)`$, with $`i=0,\mathrm{},3`$. We only need to impose that these four functions satisfy the relations (36)-(38) together with the boundary conditions (32)-(34), and that $`\phi ^{\left(0\right)}(x_S,t)`$ be a linear combination as in (30), with $`x=y=x_S`$, for some choice of the coefficients $`C(k_1,k_2)`$. Surprisingly enough, the relations (36)-(38), obtained from the quantum evolution (19), are exactly the same that we would obtain from the classical evolution (21). Therefore, if a classical trajectory $`\phi ^{\left(0\right)}(x_S,t)=\delta \left(x_Sx_{cl}\left(t\right)\right)`$ may be expressed as a linear combination (30), then it can be exactly reproduced in the quantum domain: this is more likely to happen in the case of open trajectories, for which the quantum energy spectrum is continuous, rather than for closed trajectories, where the quantum energy spectrum is discrete and most classical energy levels are forbidden; let’s examine these two cases in more detail.
An open trajectory represents a particle which in the limit $`t\pm \mathrm{}`$ behaves like a free particle; at some finite time, the particle interacts with the potential $`V\left(x\right)`$ and undergoes a scattering process. In one space dimension, there are only two possible outgoing directions, i.e. the particle may be transmitted or reflected; this is true both in the classical and in the quantum description. However, in the quantum domain we must take into account a new, highly non-classical feature, namely the tunnel effect: particles with energy less than the potential peak $`V_{max}`$ may cross the barrier, and viceversa particles with energy greater than the barrier peak may be reflected. Therefore, in the quantum case we have to remove the classical restriction that all particles with kinetic energy $`mv^2/2>V_{max}`$ are transmitted while all particles with $`mv^2/2<V_{max}`$ are reflected. Thus, knowing the energy of a quantum particle is not enough to determine whether it will be transmitted or reflected; since we believe that the physical laws are deterministic, we deduce that the property of an incoming wave-packet to be transmitted or reflected depends also on the form of the wave-packet, not only on its energy: in other words, besides the center of mass and momentum, the state of a quantum particle must include also some internal, or hidden, degree of freedom even in the case of a spinless uncharged particle.
Furthermore, the existence of non-classical tunnelling particles implies that their wave-packet must lose its sharp localization while crossing the barrier region; if this were not true, then (23) and (24) would yield a quasi-classical trajectory for the center of mass $`x\left(t\right)`$, thus preventing the appearance of non-classical effects. This is the reason why in our localization postulate we explicitly introduced the limit $`t\pm \mathrm{}`$: in general, an open quantum trajectory will be well localized when $`t\mathrm{}`$, i.e. when it is still far away from the potential region; then, while interacting with the external potential, it will spread out to some extent; finally, as $`t+\mathrm{}`$ and the particle leaves the potential region, it will recover its localization. In the next section, we will see a practical application of our approach to the tunnel effect.
Let’s now consider briefly the case of closed trajectories. It is natural to assume that quantum closed trajectories are periodic, i.e. $`\phi _{cl}^Q(x,y,t+T_{cl})=\phi _{cl}^Q(x,y,t)`$ where $`T_{cl}`$ is the revolution time. The following condition must then be fulfilled
$$\frac{1}{\mathrm{}}\left(E_{k1}E_{k2}\right)=n\frac{2\pi }{T_{cl}}$$
(39)
for some integer number $`n`$, where again $`E_{k1}`$, $`E_{k2}`$ are energy eigenvalues of the Schrödinger equation for pure states. Since normally the revolution time $`T_{cl}`$ depends on the energy of the particle, (39) may then be interpreted as a quantization condition for the energy spectrum in the quantum domain; an important exception is provided by the linear harmonic oscillator, where notoriously the revolution time is independent from the energy. In our approach the condition (39) is responsible for the existence of discrete energy spectra.
## 3 Applications
### 3.1 Free particle
In the free particle case, the quantum trajectories are simply given by
$$\phi _Q(x,y,t)=\delta \left(x_Sx_0v_0t\right)e^{\frac{i}{\mathrm{}}mv_0x_D}$$
(40)
that is, the quantum description is equivalent to the classical description. We will briefly show, by means of the Moyal-Wigner transformation, the classical probability densities associated to some well known quantum states.
In the case of a plane wave we have
$$\psi (x,t)=e^{i\frac{p_0}{\mathrm{}}\left(x\frac{p_0}{2m}t\right)}F(x,p,t)=\delta \left(pp_0\right)$$
(41)
Therefore, the plane wave describes an ensemble of classical particles moving at speed $`v_0=\frac{1}{m}p_0`$ uniformly distributed over the $`x`$ axis (the normalization is one particle per unit of length).
In the case of a particle perfectly localized at $`x=0`$, whose initial quantum state is given by $`\psi (x,0)=2\pi \mathrm{}\delta \left(x\right)`$, we have:
$$\psi (x,t)=e^{i\frac{p}{\mathrm{}}\left(x\frac{p}{2m}t\right)}\text{d}pF(x,p,t)=\delta \left(x\frac{p}{m}t\right)$$
(42)
Thus we obtain an ensemble of classical particles following the trajectories $`x=vt`$, with $`v=\frac{1}{m}p`$, uniformly distributed in momentum space (the normalization is one particle per unit of momentum).
Finally, for a gaussian state we have:
$$\psi (x,t)=\sqrt{\frac{1}{\sqrt{2\pi }\mathrm{\Delta }x_0\alpha \left(t\right)}}e^{\frac{x^2}{4\mathrm{\Delta }x_{0}^{}{}_{}{}^{2}\alpha \left(t\right)}}F(x,p,t)=\frac{1}{\pi \mathrm{}}e^{2\mathrm{\Delta }x_{0}^{}{}_{}{}^{2}\frac{p^2}{\mathrm{}^2}\frac{\left(xp/mt\right)^2}{2\mathrm{\Delta }x_{0}^{}{}_{}{}^{2}}}$$
(43)
where $`\alpha \left(t\right)=1+\frac{i\mathrm{}}{2m\mathrm{\Delta }x_{0}^{}{}_{}{}^{2}}t`$. The marginal deviations for the classical joint distribution $`F(x,p,t)`$ at $`t=0`$ are $`\mathrm{\Delta }x_0`$ and $`\mathrm{\Delta }p_0=\frac{\mathrm{}}{2\mathrm{\Delta }x_0}`$, yielding the minimum deviation product allowed by the Heisenberg uncertainty relations. For $`t>0`$ $`\mathrm{\Delta }p`$ remains constant, while $`\mathrm{\Delta }x`$ grows according to the law $`\mathrm{\Delta }x=\sqrt{\mathrm{\Delta }x_{0}^{}{}_{}{}^{2}+\frac{\mathrm{}^2}{4m^2\mathrm{\Delta }x_{0}^{}{}_{}{}^{2}}t^2}`$, in agreement with the quantum predictions.
### 3.2 The linear harmonic oscillator
Let’s consider the case of the linear harmonic oscillator, defined by the quadratic potential
$$V\left(x\right)=\frac{1}{2}m\omega _0^2x^2$$
(44)
In the standard representation of the Schrödinger equation (1), it is well known that the energy spectrum is discrete, and the energy eigenvalues are
$$E_n=\left(n+\frac{1}{2}\right)\mathrm{}\omega _0n=0,1,2,\mathrm{}$$
(45)
On the contrary, in Section 2 we saw that for a quadratic potential the quantum evolution in the density matrix representation is equivalent to the classical evolution: the physical quantum trajectories are then given by
$$\phi _Q(x,y,t)=\delta \left(x_Sx_{cl}\left(t\right)\right)e^{\frac{i}{\mathrm{}}p_{cl}\left(t\right)x_D}$$
(46)
where $`p_{cl}`$ and $`x_{cl}`$ belong to classical trajectories:
$$p_{cl}\left(t\right)=\sqrt{2mE_0}\mathrm{cos}\left(\omega _0t+\varphi _0\right)x_{cl}\left(t\right)=\frac{1}{\omega _0}\sqrt{\frac{2E_0}{m}}\mathrm{sin}\left(\omega _0t+\varphi _0\right)$$
(47)
Therefore, all classical energy levels $`E_0`$ are also acceptable in the quantum domain, yielding a continuous energy spectrum; this seems to be a severe contradiction between our approach and ordinary quantum mechanics. We will now show that this contradiction is produced by an inherent limitation in the formalism of ordinary quantum mechanics.
In the Schrödinger equation for pure states, the hamiltonian operator plays two different roles: on one hand it generates the time evolution of the system, on the other hand it is the energy operator. Therefore its eigenvectors are at the same time the stationary states of the system and the energy dispersion-free states; its eigenvalues are at the same time the natural frequencies of the system and the possible outcomes of an energy measurement. If the natural frequencies form a discrete set, then the energy spectrum must be discrete too.
On the contrary, in the density matrix representation we have two different operators: the hamiltonian $`H`$, as expressed by (3) or by (19), generates the time evolution of the system, and therefore its eigenvalues are the natural frequencies; the operator $`E_S=\frac{1}{2m}\text{P}_S^2+V\left(\text{Q}_S\right)`$, which appears in the definition (15) of the total energy, is the “physical” energy operator, and its eigenvectors are the energy dispersion-free states. Since in the Wigner representation the operator $`E_S`$ is a c-number, its eigenvalues form a continuous set and include all the classical energy levels; this is of course a consequence of the commutation relation $`[Q_S,P_S]=0`$, and is true for all quantum systems, not only for the linear harmonic oscillator. Does this mean that all classical energy levels are always physically realizable in the quantum domain? No, because in general the energy eigenvectors do not belong to physical quantum trajectories, as defined by our localization postulate and by our quantization condition (39). However, in the specific case of the linear harmonic oscillator, the answer is yes: all classical trajectories are also quantum trajectories, and all classical energy levels are physically realizable in the quantum domain; indeed, since the revolution time does not depend on the energy, the condition (39) does not impose any restriction on the admissible energy values.
As for the natural frequencies, they form a discrete spectrum already in the classical statistical description of the linear harmonic oscillator: indeed, the classical probability density $`F(x,p,t)`$ is periodic, i.e. $`F(x,p,t+T_0)=F(x,p,t)`$ where $`T_0=\frac{2\pi }{\omega _0}`$ is the revolution time of the oscillator. Therefore its Fourier expansion contains only the integer harmonics of the fundamental frequency $`\omega _0`$; these are exactly the same observable natural frequencies provided by the standard quantum description, i.e. the differences between two energy levels as defined by (45).
Thus we conclude that standard quantum mechanics predicts the correct natural frequencies but does not in general predict the correct energy levels: the equivalence between natural frequencies and energy levels is not a physical law, rather it is a consequence of an inherent limitation of the standard quantum mechanical formalism. If we switch to the density matrix representation, and eliminate the requirement that a quantum matrix be a product of pure states $`\phi (x,y)=\psi \left(x\right)\psi ^{}\left(y\right)`$, then we should obtain the correct quantum energy levels by imposing our localization postulate and our quantization condition (39).
The validity of our approach may be confirmed by examining the ground state of the linear harmonic oscillator in the usual quantum mechanical description. It is well known that this state, corresponding to the energy level $`E=\frac{1}{2}\mathrm{}\omega _0`$, is described by the gaussian function
$$\psi (x,t)=\left(\frac{m\omega _0}{\pi \mathrm{}}\right)^{\frac{1}{4}}e^{\frac{m\omega _0}{2\mathrm{}}x^2i\frac{\omega _0}{2}t}$$
(48)
If we apply to this state the Moyal-Wigner transformation, we obtain again a gaussian function
$$F(x,p,t)=\frac{1}{\pi \mathrm{}}e^{\frac{2}{\mathrm{}\omega _0}\left(\frac{p^2}{2m}+\frac{1}{2}m\omega _0^2x^2\right)}$$
(49)
which is always positive and therefore may be interpreted as a joint probability distribution in classical phase space: it reproduces the same statistical predictions of the quantum state (48) for measurements of position and momentum. The Wigner function (49) represents indeed a stationary state, since the right hand side does not depend on time; however, it is certainly not an energy eigenvector. On the contrary, it represents an ensemble of classical particles moving with all possible energies; the energy mean value is $`E=\frac{1}{2}\mathrm{}\omega _0`$ as expected, but the standard deviation is $`\mathrm{\Delta }E=\frac{1}{2}\mathrm{}\omega _0`$, while we would expect it to vanish for an energy eigenvector. This result confirms our belief that the stationary states are not necessarily energy dispersion-free: the equivalence between stationary states and energy eigenvectors is just a shortcoming of the usual quantum mechanical formalism, and disappears when we switch to the density matrix representation.
As for the standard energy eigenstates with $`E>\frac{1}{2}\mathrm{}\omega _0`$, it is well known that the corresponding Wigner functions are negative in some phase-space regions. Therefore in our approach they have no physical meaning, i.e. they do not represent neither individual particles nor statistical ensembles.
### 3.3 The delta barrier potential
Let’s now consider a system whose potential is given by:
$$V\left(x\right)=V_0\delta \left(x\right)$$
(50)
A set of orthonormal eigenvectors may be chosen as follows:
$$f_k^1\left(x\right)=\frac{1}{\sqrt{\pi }}\mathrm{sin}kxf_k^2\left(x\right)=\frac{1}{\sqrt{\pi }}\mathrm{cos}\left(k\left|x\right|\varphi _k\right)$$
(51)
where $`k>0`$, while $`\varphi _k`$ is defined by
$$\mathrm{tan}\varphi _k=\frac{mV_0}{\mathrm{}^2k}0<\varphi _k<\frac{\pi }{2}$$
(52)
The energy eigenvalues associated to the above solutions are simply
$$E_k=\frac{\mathrm{}^2k^2}{2m}$$
(53)
as in the free particle case. A linear combination of the two states $`f_k^1`$ and $`f_k^2`$ gives the well known state
$$\psi \left(x\right)=\{\begin{array}{cc}A_0e^{ikx}+A_Re^{ikx}\hfill & \text{if}x<0\hfill \\ A_Te^{ikx}\hfill & \text{if}x>0\hfill \end{array}$$
(54)
where
$$A_R=A_0\frac{mV_0}{ik\mathrm{}^2mV_0}A_T=A_0\frac{ik\mathrm{}^2}{ik\mathrm{}^2mV_0}$$
(55)
Since a term $`e^{\pm ikx}`$ represents a plane wave with momentum $`p=\pm \mathrm{}k`$, the state (54) is usually interpreted as follows: $`A_0e^{ikx}`$ is the incident wave, moving from the left towards the barrier; $`A_Re^{ikx}`$ is the reflected wave, moving away from the barrier with negative momentum $`p=\mathrm{}k`$; $`A_Te^{ikx}`$ is the transmitted wave, moving over the barrier with the same positive momentum as the incident wave. However, classical mechanics predicts that no particle can be transmitted over an infinite potential barrier; therefore, the state (54) is an example of a non-classical feature of quantum mechanics, the so called tunnel effect. In the rest of this section we will examine the tunnel effect in the density matrix representation.
Our goal will be to find quantum trajectories representing transmitted and reflected particles; these quantum trajectories will be obtained as solutions of the simplified equations (36)-(38) and will satisfy the fundamental physical principles defined in Section 2, i.e. unitarity, energy conservation and Ehrenfest theorem. In the case of a reflected particle, there is a very simple solution: we already know that the equations (36)-(38) are equal to their classical counterparts, therefore we will choose the classical solution:
$$\phi ^{\left(0\right)}(x_S,t)=\delta \left(x_Svt\right)\theta \left(t\right)+\delta \left(x_S+vt\right)\theta \left(t\right)$$
(56)
where $`v>0`$ is the speed and $`\theta `$ is the Heaviside step function. If we apply (36) to $`\phi ^{\left(0\right)}`$, we obtain for the momentum density $`𝒫(x_S,t)`$ the following expression:
$$𝒫(x_S,t)=i\mathrm{}\phi ^{\left(1\right)}(x_S,t)=mv\left(\delta \left(x_Svt\right)\theta \left(t\right)\delta \left(x_S+vt\right)\theta \left(t\right)\right)$$
(57)
as expected, i.e. the incident wave-packet has momentum $`P=mv`$ while the reflected wave-packet has momentum $`P=mv`$. Now we integrate (37) to obtain $`\phi ^{\left(2\right)}`$; the energy density $`(x_S,t)`$ is then given by:
$`(x_S,t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}\phi ^{\left(2\right)}(x_S,t)+V\left(x_S\right)\phi ^{\left(0\right)}`$ (58)
$`=`$ $`{\displaystyle \frac{1}{2}}mv^2\left(\delta \left(x_Svt\right)\theta \left(t\right)+\delta \left(x_S+vt\right)\theta \left(t\right)\right)`$
The derivation of (58) is not straightforward, due to the products of distributions appearing in the intermediate calculations. However, it may be easily obtained by replacing (56) with the classical solution associated to the triangular potential $`V_ϵ\left(x\right)=\frac{V_0}{ϵ}\left(1\frac{\left|x\right|}{ϵ}\right)\theta \left(1\frac{\left|x\right|}{ϵ}\right)`$ and then taking the limit as $`ϵ0`$; of course, as $`ϵ0`$ the function $`V_ϵ\left(x\right)`$ approaches the delta barrier potential (50).
Let’s now turn to the case of a transmitted particle. As a first attempt, we choose the classical solution for a free particle:
$$\phi ^{\left(0\right)}(x_S,t)=\delta \left(x_Svt\right)$$
(59)
even if we know that it can’t be the final solution, since it violates the Ehrenfest theorem at $`t=0`$. From (36) we then obtain
$$\phi ^{\left(1\right)}(x_S,t)=i\frac{mv}{\mathrm{}}\delta \left(x_Svt\right)$$
(60)
which is again equal to the free particle case, since equation (36) does not involve the potential $`V\left(x_S\right)`$. The next step would be to integrate (37) and find an expression for $`\phi ^{\left(2\right)}`$; this would produce the following result:
$$\phi ^{\left(2\right)}(x_S,t)=\frac{m}{\mathrm{}^2}\left(mv^2\delta \left(x_Svt\right)+V_0\delta \left(x_S\right)\delta \left(vt\right)+V_0\delta ^{}\left(vt\right)\theta \left(x_S\right)\right)$$
(61)
which does not satisfy the condition (33) at $`t=0`$ and therefore violates the Ehrenfest theorem; responsible for this violation is the step function $`\theta \left(x_S\right)`$ appearing in the r.h.s. of (61). Similar problems arise when we try to integrate (38): the resulting expression for $`\phi ^{\left(3\right)}(x_S,t)`$ does not satisfy condition (34) at $`t=0`$, which means that energy conservation is violated too.
Thus we have shown by direct calculations what we already knew from physical considerations: the classical free field solution (59) may not represent a non-classical tunnelling particle. However, since we now know the mathematical source of the above violations, we may find a way to circumvent them: we simply must add countertems to (59), so that the unacceptable terms appearing in the expressions for $`\phi ^{\left(2\right)}(x_S,t)`$ and $`\phi ^{\left(3\right)}(x_S,t)`$ are canceled by the new ones. The resulting expression is:
$$\phi ^{\left(0\right)}(x_S,t)=\delta \left(x_Svt\right)\frac{g_1^{\left(4\right)}\left(x_S\right)}{g_1^{\left(4\right)}\left(0\right)}\delta \left(vt\right)+\frac{g_2^{\left(3\right)}\left(x_S\right)}{g_2^{\left(4\right)}\left(0\right)}\delta ^{}\left(vt\right)$$
(62)
where $`g_1\left(x_S\right)`$ and $`g_2\left(x_S\right)`$ are finite approximations of the Dirac $`\delta `$ distribution, i.e. even functions localized around $`x_S=0`$; $`g_1^{\left(n\right)}`$ and $`g_2^{\left(n\right)}`$ are their n-th derivatives. From the fact that $`g_1`$ and $`g_2`$ are even we deduce that the derivatives $`g_1^{\left(n\right)}\left(0\right)`$ and $`g_2^{\left(n\right)}\left(0\right)`$ vanish when $`n`$ is odd. The expressions for $`\phi ^{\left(1\right)}(x_S,t)`$ and $`\phi ^{\left(2\right)}(x_S,t)`$ are then:
$`\phi ^{\left(1\right)}(x_S,t)`$ $`=`$ $`i{\displaystyle \frac{mv}{\mathrm{}}}\left(\delta \left(x_Svt\right)+{\displaystyle \frac{g_1^{\left(3\right)}\left(x_S\right)}{g_1^{\left(4\right)}\left(0\right)}}\delta ^{}\left(vt\right){\displaystyle \frac{g_2^{\prime \prime }\left(x_S\right)}{g_2^{\left(4\right)}\left(0\right)}}\delta ^{\prime \prime }\left(vt\right)\right)`$ (63)
$`\phi ^{\left(2\right)}(x_S,t)`$ $`=`$ $`{\displaystyle \frac{m^2v^2}{\mathrm{}^2}}\left(\delta \left(x_Svt\right){\displaystyle \frac{g_1^{\prime \prime }\left(x_S\right)}{g_1^{\left(4\right)}\left(0\right)}}\delta ^{\prime \prime }\left(vt\right)+{\displaystyle \frac{g_2^{}\left(x_S\right)}{g_2^{\left(4\right)}\left(0\right)}}\delta ^{\left(3\right)}\left(vt\right)\right)`$ (64)
The step function has disappeared from (64), so that $`\phi ^{\left(1\right)}(x_S,t)`$ and $`\phi ^{\left(2\right)}(x_S,t)`$ now satisfy both conditions (32) and (33), required respectively by the unitarity principle and by the Ehrenfest theorem. As for $`\phi ^{\left(3\right)}(x_S,t)`$, we do not write its final expression, since it does not correspond to an observable physical quantity; it is enough to say that it now satisfies the condition (34) required by energy conservation. From (63) and (64) we may then obtain the momentum density $`𝒫(x_S,t)=i\mathrm{}\phi ^{\left(1\right)}(x_S,t)`$ and the energy density $`(x_S,t)=\frac{\mathrm{}^2}{2m}\phi ^{\left(2\right)}(x_S,t)+V\left(x_S\right)\phi ^{\left(0\right)}(x_S,t)`$.
We further note that the wave-packet defined by (62) is sharply localized at $`x_S=vt`$ when $`t0`$, while it has a non-vanishing dispersion around $`x_S=0`$ when $`t=0`$, i.e. when it crosses the barrier. This temporary loss of localization was expected, since it has been already recognized in the previous section as a necessary condition for satisfying the Ehrenfest theorem in the case of a non-classical tunnelling particle.
From what we have seen until now, we may say that (56) and (62) represent quantum trajectories associated respectively to reflected and transmitted particles, since they satisfy our localization postulate together with the three fundamental physical principles. It remains to show that they may be extended to complete solutions of the Schrödinger equation, i.e. that they may be expressed as linear combinations
$$\phi ^{\left(0\right)}(x_S,t)=\underset{i,j=1}{\overset{2}{}}\text{d}k_1\text{d}k_2C_{ij}(k_1,k_2)f_{k1}^i\left(x_S\right)f_{k2}^j\left(x_S\right)e^{\frac{i}{\mathrm{}}\left(E_{k1}E_{k2}\right)t}$$
(65)
for some choice of the coefficients $`C_{ij}(k_1,k_2)`$. We will not demonstrate it here; however, in Appendix A we will explicitly build numerical solutions of the Schrödinger equation having a very strong resemblance with (56) and (62).
Thus we have shown that in a typical quantum system, exhibiting the highly non-classical behaviour known as tunnel effect, it is still possible to find solutions which remain sharply localized around their center of mass as $`t\pm \mathrm{}`$; these solutions represent wave-packets which are totally reflected or totally transmitted. It is our belief that they represent the individual physical particles; the usual quantum wave-packets, which split into a reflected part and a transmitted part, are just a statistical mixture of our quantum trajectories. Besides, those quantum states which may not be written as positive superpositions of quantum trajectories must be rejected as physically meaningless.
As an example, we will examine again the state defined by (54). If we consider the density matrix $`\phi (x,y)=\psi \left(x\right)\psi ^{}\left(y\right)`$, we may easily compute the expressions for $`\phi ^{\left(0\right)}(x_S,t)`$, $`𝒫(x_S,t)`$ and $`(x_S,t)`$; in the case $`x_S>0`$ we obtain:
$`\phi ^{\left(0\right)}(x_S,t)`$ $`=`$ $`\left|A_T\right|^2`$ (66)
$`𝒫(x_S,t)`$ $`=`$ $`\left|A_T\right|^2\mathrm{}k`$ (67)
$`(x_S,t)`$ $`=`$ $`\left|A_T\right|^2{\displaystyle \frac{\mathrm{}^2k^2}{2m}}`$ (68)
The physical interpratation of (66)-(68) is straightforward: they describe an ensemble of transmitted particles uniformly distributed to the right of the barrier (the particle density being $`\left|A_T\right|^2`$ per unit of length), moving with speed $`v=\mathrm{}k/m`$, momentum $`p=\mathrm{}k`$ and energy $`E=\mathrm{}^2k^2/2m`$. To the left of the barrier, i.e. in the case $`x_S<0`$, we have:
$`\phi ^{\left(0\right)}(x_S,t)`$ $`=`$ $`\left|A_0\right|^2+\left|A_R\right|^2+2Re\left\{A_0A_R^{}e^{2ikx_S}\right\}`$ (69)
$`𝒫(x_S,t)`$ $`=`$ $`\left|A_0\right|^2\mathrm{}k\left|A_R\right|^2\mathrm{}k`$ (70)
$`(x_S,t)`$ $`=`$ $`\left|A_0\right|^2{\displaystyle \frac{\mathrm{}^2k^2}{2m}}+\left|A_R\right|^2{\displaystyle \frac{\mathrm{}^2k^2}{2m}}`$ (71)
The terms involving $`\left|A_0\right|^2`$ and $`\left|A_R\right|^2`$ represent respectively an ensemble of incoming and of reflected particles, with the right values for momentum and energy, $`\left|A_0\right|^2`$ and $`\left|A_R\right|^2`$ being the associated particle densities; besides, the relation $`\left|A_T\right|^2+\left|A_R\right|^2=\left|A_0\right|^2`$ ensures that the total number of particles is conserved at the two sides of the barrier. However, the last term at the r.h.s. of (69) has no physical explanation: it does not contribute to the momentum and energy densities, so it seems to describe an ensemble of particles at rest; yet, it has both positive and negative values depending on $`x_S`$, and therefore it may not be interpreted as a particle density. For this reason, we are forced to reject the state (54) as physically meaningless, since it is not a positive superposition of quantum trajectories.
On the contrary, starting from our quantum trajectories (56) and (62) it is very easy to build a statistical ensemble having the same physical properties as (54), but with no strange unphysical terms. We simply have to write:
$$\phi ^{\left(0\right)}(x_S,t)=_{\mathrm{}}^+\mathrm{}\left(\left|A_R\right|^2\phi _R^{\left(0\right)}(x_S,tt_0)+\left|A_T\right|^2\phi _T^{\left(0\right)}(x_S,tt_0)\right)v\text{d}t_0$$
(72)
where of course $`\phi _R^{\left(0\right)}`$ is the “reflected” quantum trajectory (56) while $`\phi _T^{\left(0\right)}`$ is the “transmitted” quantum trajectory (62). This is another example of the limitations inherent to the standard formalism of quantum mechanics, which may be overcome by switching to the density matrix representation.
Finally, we point out that the quantum trajectories (56) and (62) must not be considered as the only possible quantum trajectories for the delta barrier potential: specifically, it is well known that standard quantum wave-packets impinging onto a potential barrier may experiment a shift in time, either a time delay or a time advance. For instance Nakazato shows, precisely in the case of the delta barrier potential, that a gaussian wave-packet under certain conditions splits into a reflected and a transmitted wave-packet which, in the limit $`t\mathrm{}`$, are again approximately gaussian; both wave-packets are delayed in time with respect to the ideal case of an instant transmission or reflection at the barrier. Therefore, it is quite possible that quantum trajectories may be found which are shifted in time during the interaction with the barrier; this is again a highly non-classical feature of tunnelling particles, since of course in the classical case the interaction time with a delta barrier vanishes and no time shift may occur.
## 4 Discussion
In the early days of quantum mechanics, Schrödinger tried to interpret the solutions of his equation as “matter waves”, i.e. he thought that the individual particles could be described by sharply localized wave-packets. Unfortunately, this realistic interpretation of the wave-function was soon abandoned, when it became clear that wave-packets lose their localization with time and become indefinitely extended as $`t\mathrm{}`$; this is clearly incompatible with the fact that particles are always detected in small space regions. Then Born proposed his interpretation of the wave-function as “probability waves”, and his statistical postulate, incorporated in the standard Copenhagen interpretation, has survived until today. The theory built upon this postulate, standard quantum mechanics, gives predictions which are in spectacular agreement with all experiments ever performed and was successfully extended to include relativistic phenomena, such as particle creation/annihilation, and to describe all the relevant interactions of the microphysical world.
Nevertheless, the history of quantum mechanics is not only a history of experimental successes; it is also a history of never ending debates about its foundations. Countless interpretations have been proposed, alternative to the standard Copenhagen interpretation, but none of them has been capable of gaining universal acceptance inside the scientific community: the supporters of the various interpretations still debate throughout the specialized literature, emphasizing both the merits of their favourite theories and the flaws of the opponent interpretations. It is not my intention here to give a detailed description of the fundamental problems which still prevent a completely satisfactory interpretation of quantum mechanics; I just want to focus the reader’s attention about a few concepts which represent the basis on which the present paper is built.
Realism. A physical system exists as an objective reality, independently from the presence of an observer performing measurements on it; therefore a sound physical theory should be able to describe the objective properties of the system. On the contrary, standard quantum mechanics only provides statistical predictions about the possible outcomes of experiments, denying the existence of objective properties prior to measurement: the famous Schrödinger cat is described by a superposition of “dead” state and “alive” state, until somebody opens the box and the cat suddenly collapses to either “dead” or “alive”. Transposing the paradox to the description of a single particle, we have a wave-function which may be extended over an arbitrarily large space region; then the particle’s position is measured and the wave-function suddenly collapses to the state $`|x_0`$, where $`x_0`$ is the measurement’s result. A much more realistic view would be to think that immediately before the measurement the particle is already localized at the position $`x_0`$, this position being an objective property of the particle: the measurement simply brings to light this property already possessed by the particle. I am aware that such way of thinking would be labeled as “naïf realism” by the majority of physicists; nevertheless, I am still convinced that it is the only acceptable starting point for building a sound physical theory.
Determinism. Probabilities are a very useful tool for handling systems whose state, for whatever reason, is not completely known: they allow us to make calculations and to obtain interesting results even if we are partially ignorant about the system’s state. Standard quantum mechanics has radically changed this original meaning of the probability concept: on one hand the wave-function is assumed to represent a complete description of an individual physical system, on the other hand the Born postulate implies that we cannot predict with absolute certainty the result of a measurement performed on a system, even if its wave-function is perfectly known. Thus the concept of uncertainty becomes a fundamental principle of the theory, and the Heisenberg uncertainty relations play a central role in the so-called “quantum revolution”.
For many years it was generally believed that a fully deterministic theory reproducing the same statistical predictions of quantum mechanics could not exist; Von Neumann’s proof of the impossibility of “hidden variables” theories was considered as the last word about the subject. However, in 1952 Bohm was able to build exactly this kind of theory: in Bohmian mechanics, the state of a quantum particle is given by its wave-function $`\psi \left(x\right)`$ and by its position $`Q`$. The time evolution of the wave-function follows the usual Schrödinger equation, while the time evolution of the particle’s position is obtained from a new equation, satisfying the following remarkable property: if we consider at time $`t=0`$ a statistical ensemble of particles, all described by the same wave-function $`\psi (x,0)`$, whose positions $`Q`$ are distributed with a probability density $`\rho (x,0)=\left|\psi (x,0)\right|^2`$, then the equality $`\rho (x,t)=\left|\psi (x,t)\right|^2`$ will remain true for all times $`t>0`$, thus reproducing the same probabilities of standard quantum mechanics for measurements of position. Hence Bohmian mechanics is at the same time realistic, the position $`Q`$ being considered as an objective property of the particle, and deterministic, since the result of a measurement performed on the system is uniquely determined by the system’s state (here the assumption is made that only those observables which may be expressed in terms of position measurements have physical meaning).
Of course, the existence of at least one “hidden variables” theory equivalent to standard quantum mechanics is enough to disprove all claims about the impossibility of such theories. In the present paper, as in the previous one, I strongly support the view that the fundamental laws of nature must be deterministic: if there is uncertainty about the results of measurements performed on a physical system, this simply means that the state of the system is not entirely known. For instance, let’s consider two particles impinging onto the same potential barrier, both starting with the same initial position and momentum; if the first particle is reflected while the second one is transmitted, this is not a consequence of some fundamental uncertainty in the laws of nature: on the contrary, this means that position and momentum are not a complete description of the particle’s state.
Nevertheless, there is one serious problem that must be faced by every deterministic theory trying to reproduce the statistical predictions of quantum mechanics: of all possible probability distributions over the system’s state, why only a few may be observed in real experiments, while the majority of them seems to be forbidden by nature? In the case of tunnelling, we can imagine a statistical ensemble of particles all being reflected (or transmitted); however, the experimental evidence tells us that the real ensembles are always formed by both reflected and transmitted particles, and their relative weight in the mixture depends on the energy of the ensemble in a well defined manner. In Bohmian mechanics, we could think of an ensemble of particles having all the same initial wave-function $`\psi (x,0)`$ and the same initial position $`Q_0`$; this ensemble would have a dispersion-free position at all future times, thus contradicting the quantum principle that a particle’s trajectory may not be perfectly known at all times.
To solve this problem, the supporters of “hidden variables” theories usually propose to complete the deterministic equations of motion with some kind of statistical postulate, in analogy with the classical case where the thermodynamic equilibrium hypothesis (Gibbs postulate) allows to obtain the laws of statistical mechanics from the deterministic laws of newtonian mechanics. In the case of Bohmian mechanics, this statistical postulate is the already mentioned hypothesis that $`\rho (x,0)=\left|\psi (x,0)\right|^2`$ at some initial time $`t=0`$; in the literature, this postulate is usually referred to as the “quantum equilibrium hypothesis” .
Simplicity. The history of physics supports the belief that the fundamental laws of nature should be as simple as possible; even if this is not a clearly defined principle, it generally means that the fundamental objects of a theory, together with its fundamental equations, should be reduced to the smallest possible number. For instance, the theory of special relativity unified the electric field and the magnetic field into a single object, the electromagnetic field, at the same time reducing the four Maxwell equations to a single covariant expression; besides, the elimination of absolute space and time brought as a consequence that the laws of physics must be the same in all inertial reference frames (Lorentz invariance). Then came the general theory of relativity, which further simplified the definition of reference frame by eliminating the priviliged role of inertial frames and by establishing the equivalence between accelerated frames and gravitational fields; as a consequence, the laws of physics must be the same in *all* reference frames, inertial or not (general covariance). Besides, general relativity provides an elegant explanation of the gravitational forces as effect of space-time curvature and clarifies the status of space-time, which becomes a true physical object, not only an arena where physical events take place; it is generally believed that the Einstein’s equations for the gravitational field are one of the highest peaks reached by the physical science in its search for simplicity.
Unfortunately, the same level of simplicity has never been reached by quantum mechanics and quantum field theory; on the contrary, there are many controversial questions, ranging from the interpretational problems to the divergencies which plague quantum field theory and which still prevent the unification with gravity. Here I just want to concentrate about one point, which is relevant to the approach that I am proposing in the present paper: the so called wave-particle duality. It is well known that particles manifest wave-like behaviours in some experimental situations: for instance the interference fringes obtained in the famous two-slit experiment appear to be a consequence of the linear superposition of two waves, adding up where the phase is the same and canceling each other where they have opposite phases. At the same time, the individual particles are always detected as (almost) point-like objects, i.e. each of them leaves a well localized spot on the screen: the interference fringes may be seen only after a big number of particles has been detected. This wave-particle duality is not explained by standard quantum mechanics, rather it is considered as a postulate, not requiring further explanations; embarrassing questions such as which slit did the particle pass through are rejected as meaningless, since in standard quantum mechanics there is no room for such an old-fashioned concept as the trajectory of a particle.
On the contrary, Bohmian mechanics provides a clear explanation of wave-particle duality: in the case of the two-slit experiment the wave function passes through both slits, while the particle trajectory, expressed by the function $`Q\left(t\right)`$, passes through but one of the slits. However, the point that I want to underline here is that this explanation does not meet the simplicity criterium defined above: instead of unifying the two concepts (particles and waves) into one single object, capable of manifesting both behaviours, the duality is solved by stating that particles and waves exist as two separate entities: the number of fundamental objects of the theory, together with its fundamental equations, is increased with respect to standard quantum mechanics, while the predictive power of the theory remains the same. This is probably the main reason why Bohmian mechanics has been considered by the majority of physicists more like a mathematical curiosity, rather than a serious alternative to standard quantum mechanics.
As a counteraxample, I will now briefly outline how the wave-particle duality is solved by a recently proposed theory, namely Hasselmann’s “metron model” . In this theory, the only fundamental object is the metric tensor $`g_{LM}`$ defined over a higher-(eight- or nine- )dimensional space, where the first four dimensions define the ordinary space-time, while the last four or five are extra-space dimensions. The only fundamental equation is the Einstein’s gravitational field equation
$$R_{LM}=0$$
(73)
where $`R_{LM}`$ is the Ricci curvature tensor. The author shows that the simple equation (73) has “an extremely rich nonlinear structure which encompasses all the principal interactions of quantum field theory and can be used as the foundation of a unified deterministic theory of fields and particles”. Specifically, it is postulated that soliton-type solutions exist, named “metrons”, consisting of “a localized, strongly non linear core and a set of linear far fields … The core is the origin of the corpuscular properties of matter, while the … far fields give rise to the wave-like interference phenomena”. Thus the metron model unifies gravity with the other forces of nature and solves the wave-particle duality, and all this is achieved starting from equation (73), which is even simpler than the original Einstein’s field equations: indeed the energy-momentum tensor, which appears as an external source term in the original field equations for gravity, is not present in the fundamental equation of the metron model; on the contrary, the author shows that the standard energy-momentum tensor arises from the contraction of the extra-space components of the Riemann curvature tensor. Hasselmann’s model is a striking example of a theory which meets our simplicity criterium; besides, it has more than one point in common with the approach that I propose in the present paper.
So far we have identified three important guidelines for building a sound physical theory: realism, determinism, simplicity; now let’s see how the approach presented in this paper tries to meet these criteria. Basically, what I propose is to revive the original Schrödinger’s idea of the wave-function as “matter waves”. The main obstacle to this idea, i.e. the fact that wave-packets spread out with time, is overcome by choosing as the fundamental equation of the theory the Schrödinger equation for density matrices, instead of the one for pure states on which standard quantum mechanics is based. In Section 2 we have seen that in the density matrix representation the Heisenberg uncertainty principle is not true, i.e. there are fields perfectly localized both in position space and in momentum space; as a consequence, wave-packets do not necessarily spread out with time. We have further postulated that the only physical solutions (here “physical” means “representing individual particles”) are those which remain well localized as $`t\pm \mathrm{}`$, labeled “quantum trajectories”. In Section 3 we have seen that quantum trajectories exist in three important cases: the free particle, the linear harmonic oscillator and the delta barrier potential; I am firmly convinced that they may be proved to exist in all cases of physical interest.
Our approach is deterministic: if we know the state $`\phi (x,y,0)`$ of the particle at $`t=0`$, then the Schrödinger equation (4) determines the state $`\phi (x,y,t)`$ at all future times; besides, the values of the observable quantities (center of mass, momentum and energy) are uniquely determined from the state by means of equations (13)-(15). Note that the particle’s position is not a well defined concept for a wave-packet, and therefore we replace it with its center of mass; however, since we postulate that the real wave-packets remain sharply localized as $`t\mathrm{}`$, the wave-packet’s center of mass is practically indistinguishable from the position of a pointlike particle.
Our approach is also realistic: no “Schrödinger cat” state is allowed and no wave-function collapse is needed; the position of the particle, i.e. the center of mass of a sharply localized wave-packet, is known prior to measurement and therefore may be considered as an objective property of the particle. Of course, since the state $`\phi (x,y)`$ depends upon two position coordinates, our approach postulates the existence of an objective extra-space coordinate: in addition to the “physical” position $`x_S=(x+y)/2`$, we have the “auxiliary” position coordinate $`x_D=xy`$, having observable effects only around the point $`x_D=0`$. The idea of extra spacetime dimensions has a long lasting tradition in physics, which was initiated by the original Kaluza-Klein theories and has been recently revived by the already cited Hasselmann’s metron model.
As for simplicity, we have only one fundamental object, the field $`\phi (x,y)`$, and one fundamental equation, the Schrödinger equation (4). The wave-particle duality is solved by postulating that the physical wave-packets are sharply localized as $`t\pm \mathrm{}`$, i.e. when the particles are produced and detected: this explains why the individual particles are always experienced as (almost) pointlike objects. At the same time, the wave-packets may spread out to some extent while interacting with the potential $`V\left(x\right)`$; this temporary loss of localization is responsible for non-classical phenomena such as tunnelling or interference, i.e. for the wave-like behaviour of quantum particles. We may easily imagine the form of a quantum trajectory in the case of the two-slit experiment: it starts as a localized wave-packet as $`t\mathrm{}`$, then when approaching the potential region it splits into two separate wave-packets; each wave-packet passes through one of the slits; then, while leaving the potential region, the two wave-packets merge again and as $`t+\mathrm{}`$ a single localized wave-packet is detected, producing a pointlike spot on the screen. There is an important difference between our approach and standard QM concerning the concept of interference: in standard QM, interference takes place between two plane (or spherical) waves defined over wide space regions; in our approach, interference is localized in space and time, i.e. it is confined to the potential region and to the time interval when the physical wave-packet interacts with the potential. Outside this space-time region our wave-packets move along classical trajectories, behaving as classical free particles; however, the interaction with the potential bends the outgoing trajectories with respect to the classical case, producing statistical effects which may not be explained classically, like the interference fringes in the two-slit experiment.
The superposition principle of standard QM does not apply to our quantum trajectories, i.e. a linear combination of quantum trajectories does not in general represent an individual particle. However, it may be accepted in its statistical meaning: a positive superposition of quantum trajectories may be taken to represent a statistical ensemble (or mixture) of particles. This statistical generalization, of course, is not enough to reproduce the predictions of standard QM; we still need some kind of statistical postulate (similar to the quantum equilibrium hypotesis in Bohmian mechanics) to explain, for instance, the relative frequency of transmitted and reflected particles in the case of tunnelling, or the exact form of the interference fringes in the case of the two-slit experiment. This remains an open issue in our approach; nevertheless, it seems to me that the solution to this problem should be made easier by the fact that our fundamental equation, both for individual particles and for statistical ensembles, is exactly the same equation which describes the time evolution of density matrices in standard QM: we did not change the equation, we only adopted a different criterium to decide which solutions are physically meaningful and which are not.
There is a second open issue in our approach: the definition of quantum trajectory is slightly vague, since we did not specify exactly what we mean by “well localized wave-packet”; should it be a Dirac delta function, with vanishing dispersion, or is it allowed to have a finite extension (for instance a sharp gaussian)? Moreover, the Schrödinger equation (4) admits solutions which do not have physical meaning, i.e. they are neither quantum trajectories, nor positive superpositions of quantum trajectories; we had to introduce a specific postulate to define which are the physical solutions, but clearly a truly fundamental theory should contain in itself a mechanism for selecting the physical solutions and discarding the unphysical ones. The origin of this problem is evident: it is rooted in the linearity of the Schrödinger equation. It is linearity which allows, by means of the superposition principle, to obtain finite dispersion wave-packets starting from delta functions; it is again linearity which allows both positive and negative superpositions, thus producing solutions for which no physical interpretation is possible, not even at the statistical level.
Linearity is a much celebrated feature of standard quantum mechanics: it is often claimed that without the superposition principle the wave-like features of elementary particles could not be explained. The mathematical apparatus of the theory is totally based on linear objects (Hilbert spaces, matrix operators, commutation relations, …); since this linear apparatus happens to work pretty well in the non-relativistic case, it has been extended with few modifications to describe relativistic phenomena. Even if the experimental successes of quantum field theory may not be denied, I do not share all this enthusiasm about linearity; on the contrary my opinion is that linearity is responsible for some of the problems which plague quantum mechanics and quantum field theory. At the interpretational level, it is the main origin of confusion about the meaning of the wave-function: does it represent an individual particle or a statistical ensemble of particles? Of course this question may not find an answer inside a linear theory, where both quantum trajectories (representing individual particles) and positive superpositions of quantum trajectories (representing statistical mixtures) are solutions of the fundamental equation of the theory! At the formal level, I am convinced that the divergences encountered in quantum field theory, with the consequent need of cumbersome renormalization procedures (not always effective, as in the case of gravity), are originated by the attempt of forcing a linear structure upon a phenomenology which is essentially non-linear.
As a consequence, I don’t believe that an attempt to refine our definition of quantum trajectory would be much useful, if confined to the linear Schrödinger equation. A much more useful effort would be trying to add non-linear terms to the equation, with the goal of eliminating the unphysical solutions: the only stable solutions of the new equation should be those which represent individual particles. The premises for a non-linear modification of the Schrödinger equation have been studied by many authors (see for instance ); these explorations have been motivated mainly by the general observation that “all linear equations describing the evolution of physical systems are known to be approximations of some nonlinear theories, with only one notable exception of the Schrödinger equation” . This is undoubtedly a well founded observation; unfortunately, our motivation is much more specific, and the existing literature seems to be of little help in our case, especially because the starting point of the above cited authors is usually the Schrödinger equation for pure states, not the one for density matrices. Besides, we must remember that the non-relativistic Schrödinger equation is just an approximation for some relativistic equation, the Dirac equation for fermions or the Klein-Gordon equation for bosons; it seems to me that the search for a fundamental equation should start already at the relativistic level, i.e. the non-linear modifications should be applied directly to the Dirac or Klein-Gordon equations. An example of a non-linear modification of the Dirac equation has already been proposed in my first paper , where the non-linearity was introduced by coupling the Dirac field with the classical electromagnetic field.
The attentive reader may have noticed that our main effort has been devoted to the subject of open quantum trajectories: closed trajectories have been introduced only briefly at the end of Section 2 and examined in detail only for the linear harmonic oscillator (Section 3). The main conclusion of the present paper about closed quantum trajectories is the distinction between physical energy levels and natural frequencies, two concepts which in standard quantum mechanics are taken to coincide; in our approach, the natural frequencies are the eigenvalues of the hamiltonian operator $`H`$, defined by (3), while the energy levels are the eigenvalues of the physical energy operator $`E_S=\frac{1}{2m}\text{P}_S^2+V\left(\text{Q}_S\right)`$. However, since the spectrum of the operator $`E_S`$ is continuum, thus including all classical energy levels, the fact that the outcomes of energy measurements are quantized is explained by our approach in a different way than in standard QM, namely by means of the quantization condition (39). Of course this subject needs further investigation; here I just want to point out a simple observation, which may have some interesting physical consequences. The observation is, trivially, that nobody has ever seen what happens inside an isolated atom: what we see is how the atom reacts to some external excitation, for instance the freqeuncy of the light emitted by the atom after it has been excited by the collision with an accelerated electron; the hypotesis that the emitted light frequencies correspond to jumps of the internal electrons between two different energy levels, as defined by standard QM, may not be experimentally tested. It is the only possible explanation if we confine ourselves to the Schrödinger equation for pure states; however, in the density matrix representation there may be other explanations, in terms of closed quantum trajectories whose total energy does not necessarily coincide with the energy levels of standard QM. In the end, the only observable effect is that the emitted light frequencies correspond to the natural frequencies of the system under examination, and these are the same both in standard QM and in the density matrix representation.
The present paper has been entirely devoted to the case of a single particle; in Appendix B we will briefly sketch a possible way of dealing with many-particle systems in the density matrix representation. Here I just want to point out that every realistic interpretation of a many-particle quantum system must face the well known dilemma of quantum non-locality: the result of a given measure performed on one particle seems to be influenced by what kind of measurement is performed on a second particle, distant from the first. However, this weird feature of quantum mechanics is not fully established: there are controversies still open both at the theoretical and at the experimental level. On one hand, there are authors who believe that the quantum probabilities should be considered as conditional probabilities and therefore they have no right to enter the Bell inequalities; on the other hand there are still doubts about the experimantal violation of Bell inequalities, due to the well known loopholes . Therefore the impossibility of reproducing the experimental predictions of standard QM (and the observed experimental results) by means of local realistic theories has not yet been proved beyond doubt, and maybe never will be; my attitude about this point is not different from the one I have taken in the single particle case, where those quantum states which could not be expressed as positive superpositions of quantum trajectories were rejected as physically meaningless: in the many-particle case, I am inclined to reject as physically meaningless those quantum states whose statistics may not be reproduced by local realistic theories.
## Appendix A Numerical quantum trajectories for the delta barrier potential
As a first step, we normalize the Schrödinger equation by imposing $`m=\mathrm{}=V_0=1`$; this choice does not result in a loss of generality, since we may get back the original Schrödinger equation through a suitable rescaling of the coordinates $`x`$, $`t`$ and of the field $`\phi `$. As a second step, we must choose a finite set of eigenvectors out of (51), since it is impossible to deal numerically with an infinite continuum of eigenvectors; therefore we will consider only a limited space region $`L<x<L`$ and only those eigenvectors $`f\left(x\right)`$ which behave smoothly at the borders of such region, i.e. $`f\left(L\right)=f\left(L\right)`$ and $`f^{}\left(L\right)=f^{}\left(L\right)`$. Our set of eigenvectors will then be defined by:
$`f_n\left(x\right)`$ $`=`$ $`\mathrm{cos}\left(k_n\left|x\right|\varphi _n\right)\text{for}n\text{even}`$ (74)
$`f_n\left(x\right)`$ $`=`$ $`\mathrm{sin}k_nx\text{for}n\text{odd}`$ (75)
For $`n`$ odd, $`k_n`$ is simply given by
$$k_nL=\frac{n+1}{2}\pi $$
(76)
while for $`n`$ even it may be obtained by solving the coupled equations:
$`\mathrm{tan}\varphi _n`$ $`=`$ $`{\displaystyle \frac{1}{k_n}}0<\varphi _n<{\displaystyle \frac{\pi }{2}}`$ (77)
$`k_nL`$ $`=`$ $`\varphi _n+{\displaystyle \frac{n}{2}}\pi `$ (78)
For our computations, we will choose $`L=100`$.
Our numerical quantum trajectories will then have the form
$$\phi (x,y,t)=\underset{i,j=0}{\overset{N}{}}C_{ij}\phi _{ij}(x,y,t)=\underset{i,j=0}{\overset{N}{}}C_{ij}f_i\left(x\right)f_j\left(y\right)e^{i\left(\omega _j\omega _i\right)t}$$
(79)
where $`\omega _i=k_i^2/2`$; our goal will be to find coefficients $`C_{ij}`$ such that (79) provides good approximations for the simplified solutions (56) and (62), which represent reflected and transmitted particles. To accomplish this task, we first associate to every quantum trajectory $`\phi (x_S,x_D,t)`$ the three functions
$`\rho (x_S,t)`$ $`=`$ $`\phi (x_S,x_D,t)|_{x_D=0}`$ (80)
$`\mathrm{}(x_S,t)`$ $`=`$ $`i{\displaystyle \frac{}{x_D}}\phi (x_S,x_D,t)|_{x_D=0}`$ (81)
$`\epsilon (x_S,t)`$ $`=`$ $`\left({\displaystyle \frac{1}{4}}{\displaystyle \frac{^2}{x^2}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{^2}{y^2}}+\delta \left(x_S\right)\right)\phi (x_S,x_D,t)|_{x_D=0}`$ (82)
which represent respectively the mass density, the momentum density and the energy density; the reader may notice that the definition (82) for the energy density differs from the usual one obtained from (15). The reason is that, while the total energy is the same, the choice made in (82) eliminates from the beginning the possible appearance of delta functions at $`x_S=0`$.
Then we define the scalar product
$`\phi _1\phi _2`$ $`=`$ $`{\displaystyle _T^T}\text{d}t{\displaystyle _L^L}\text{d}x_S(w_0\rho _1(x_S,t)\rho _2(x_S,t)+`$ (83)
$`+w_1\mathrm{}_1(x_S,t)\mathrm{}_2(x_S,t)+w_2\epsilon _1(x_S,t)\epsilon _2(x_S,t))`$
where $`w_0`$, $`w_1`$ and $`w_2`$ are appropriate weights, while $`[T,T]`$ is a time interval to be defined. To obtain our numerical quantum trajectories we will project, by means of the scalar product (83), the simplified solutions (56) and (62) over the linear space defined by (79). For instance, if $`\phi _R`$ is the solution associated to a reflected particle as defined by (56)-(58), the coefficients $`C_{ij}`$ will be obtained by solving the system of $`(N+1)^2`$ linear equations
$$\underset{i,j=0}{\overset{N}{}}C_{ij}\phi _{kl}\phi _{ij}=\phi _{kl}\phi _Rk,l=0,\mathrm{},N$$
(84)
which requires the inversion of a $`(N+1)^2\times (N+1)^2`$ square matrix. The choice of $`N`$ is then dictated by the computational resources of the machine used for performing this matrix inversion; in our case we will set $`N=100`$.
For both reflected and transmitted particles we will choose the speed $`v=1`$; in the standard quantum mechanical treatment this is the speed for which the transmission probability and the reflection probability are the same. Besides, due to numerical considerations, we will not use delta functions in the definitions of $`\phi _R`$ and $`\phi _T`$. On the contrary, we will use the smooth function
$$f\left(x\right)=\frac{8}{3\pi \mathrm{\Delta }x}\mathrm{cos}^4\left(\frac{x}{\mathrm{\Delta }x}\right)\theta \left(\frac{\pi }{2}\frac{\left|x\right|}{\mathrm{\Delta }x}\right)$$
(85)
which in the limit $`\mathrm{\Delta }x0`$ behaves like a Dirac delta function. In (85), $`\theta `$ is the Heaviside step function, while $`\mathrm{\Delta }x`$ is a finite dispersion to be defined; for our computations we will use $`\mathrm{\Delta }x=3`$. Besides, we are now able to choose the time $`T`$ appearing in the definition (83) of our scalar product: we will set $`T=L\frac{\pi }{2}\mathrm{\Delta }x`$, so that the wave-packet $`f\left(xvt\right)`$ is entirely included in the space region $`[L,L]`$ when $`t[T,T]`$.
Now we are ready to show the results of our computations. In the case of a reflected particle, with a suitable choice of the coefficients $`w_0`$, $`w_1`$ and $`w_2`$, we obtain the quantum trajectory described by figure 1. In figure 1.a, we show the wave-packet representing the mass density $`\rho \left(x_S\right)`$ at different times $`t=\pm 80,\pm 40,0`$. In figures 1.b and 1.c, we show the wave-packets associated to the momentum and energy densities at the same times; the momentum density at $`t=0`$ may not be displayed, since it vanishes identically for all $`x_S`$. Finally, in figures 1.d and 1.e we show the mean value $`x_M`$ and the standard deviation $`\sigma _x`$ of the wave-packet as a function of time; these two quantities have been obtained using as weight the function $`\rho ^4`$, for instance the mean value is given by
$$x_M=\frac{_L^{+L}x\rho ^4\left(x\right)\text{d}x}{_L^{+L}\rho ^4\left(x\right)\text{d}x}$$
(86)
This choice has been made to reduce the effect of the low amplitude noise appearing in figure 1.a. From figure 1, it is clear that our quantum trajectory is a very good approximation of the corresponding classical trajectory. Specifically, the standard deviation is approxmately constant when $`t`$ is different from zero, but is reduced for $`t=0`$, due to the overlapping of the incoming part and the reflected part of the wave-packet; this is qualitatively the same behaviour that may be seen in the classical case.
Turning now to the case of a transmitted particle, we make a slight change to the definition of our scalar product (83) and multiply the integrand term by the factor
$$p(x_S,t)=1\theta \left(\pi \mathrm{\Delta }x\left|x_S\right|\right)\theta \left(\frac{\pi }{2}\mathrm{\Delta }x\left|t\right|\right)$$
(87)
that is, we eliminate from the integral the region defined by $`x_S[\pi \mathrm{\Delta }x,\pi \mathrm{\Delta }x]`$ and $`t[\frac{\pi }{2}\mathrm{\Delta }x,\frac{\pi }{2}\mathrm{\Delta }x]`$. The reason is that we do not want to force any form to the wave-packet while it is tunnelling through the barrier; after all, this is a non-classical behaviour, and we do not have any a priori knowledge of the deformations produced on the wave-packet by this non-classical interaction. With a suitable choice of the coefficients $`w_0`$, $`w_1`$ and $`w_2`$, we then obtain the quantum trajectory described by figure 2. The time evolution of the mean value $`x_M`$ (figure 2.d) is not different from the case of a classical free particle; however, from figure 2.a we see that at $`t=0`$ the wave-packet is heavily deformed, while from figure 2.e we see that the standard deviation increases during the interaction time, as expected.
Thus we have found two complete solutions of the Schrödinger equation which approximate pretty well the simplified solutions (56) and (62). Our numerical solutions have some limitations: their range of validity is restricted to the space-time region defined by $`x_S[L,L]`$ and $`t[T,T]`$, while outside this region the wave-packet looses abruptly its localization; besides, the localization of our wave-packets is not so sharp even inside the region of validity. However, it should be clear that these limitations arise from the fact that we are working with a finite set of eigenvectors: if we were able to perform our calculations in the limit $`N\mathrm{}`$, the validity of our quantum trajectories could be extended to larger space-time regions and to sharper wave-packets, thus approaching the limit $`L\mathrm{}`$, $`T\mathrm{}`$ and $`\mathrm{\Delta }x0`$.
## Appendix B Brief introduction to the many-particle case
In standard QM, a system of $`N`$ particles is described by a wave-function $`\psi (x_1,\mathrm{},x_N)`$ and the fact that $`\psi `$ is defined over a configuration space, whose dimension depends on the particles’ number, is a further obstacle to a realistic interpretation of the wave-function; therefore this undesired feature must disappear in the density matrix representation. To be simple, let’s consider a system of two particles moving in one space dimension under the effect of an external potential $`V(x_1,x_2)`$; to avoid possible complications which may arise in the case of identical particles, we will further suppose that the two particles have different masses $`m_1`$ and $`m_2`$. The Schrödinger equation for the density matrix $`\phi (x_1,x_2,y_1,y_2)`$ is then:
$$i\mathrm{}\frac{\phi }{t}=\underset{i=1}{\overset{2}{}}\frac{\mathrm{}^2}{2m_i}\left(\frac{^2\phi }{x_i^2}\frac{^2\phi }{y_i^2}\right)+\left[V(x_1,x_2)V(y_1,y_2)\right]\phi $$
(88)
In our approach, the solutions of equation (88) are taken to represent individual physical systems instead of statistical mixtures. To eliminate the dependence of $`\phi `$ upon the configuration variables $`(x_1,y_1)`$ and $`(x_2,y_2)`$, we would like to impose the following separability condition:
$$\phi (x_1,x_2,y_1,y_2,t)=\phi _1(x_1,y_1,t)\phi _2(x_2,y_2,t)$$
(89)
so that the two functions $`\phi _1`$ and $`\phi _2`$ would describe separately the time evolution of the two particles and could be thought to depend upon just one coordinate pair $`(x,y)`$, as in the single particle case. Unfortunately, the condition (89) is in general too strong, and may be satisfied only in some special cases, for instance when $`V(x_1,x_2)=V_1\left(x_1\right)+V_2\left(x_2\right)`$. However, since we know that the observable quantities depend only upon the values of $`\phi `$ in the region where $`x_1y_1`$ and $`x_2y_2`$, we may impose a weaker separability condition: therefore we first define
$$x_{iS}=\frac{1}{2}\left(x_i+y_i\right)x_{iD}=x_iy_ii=1,2$$
(90)
and then require that (89) be satisfied only far small values of $`x_{1D}`$ and $`x_{2D}`$.
For $`x_{1D}=x_{2D}=0`$ we then obtain the condition
$$\phi |_{x_{1D}=x_{2D}=0}=\rho _1(x_{1S},t)\rho _2(x_{2S},t)$$
(91)
while at first order in $`x_{1D}`$ and $`x_{2D}`$ we obtain:
$`i\mathrm{}{\displaystyle \frac{\phi }{x_{1D}}}|_{x_{1D}=x_{2D}=0}`$ $`=`$ $`𝒫_1(x_{1S},t)\rho _2(x_{2S},t)`$ (92)
$`i\mathrm{}{\displaystyle \frac{\phi }{x_{2D}}}|_{x_{1D}=x_{2D}=0}`$ $`=`$ $`\rho _1(x_{1S},t)𝒫_2(x_{2S},t)`$ (93)
If we now extend the definitions (13) and (14) of center of mass and momentum to the two-particle case, we easily obtain
$`Q`$ $`=`$ $`Q_1+Q_2={\displaystyle x_S\rho _1\left(x_S\right)\text{d}x_S}+{\displaystyle x_S\rho _2\left(x_S\right)\text{d}x_S}`$ (94)
$`P`$ $`=`$ $`P_1+P_2={\displaystyle 𝒫_1\left(x_S\right)\text{d}x_S}+{\displaystyle 𝒫_2\left(x_S\right)\text{d}x_S}`$ (95)
where we supposed that both $`\rho _1`$ and $`\rho _2`$ satisfy the unitarity condition $`\rho _i\left(x_S\right)\text{d}x_S=1`$. As for the energy, it is clearly impossible to define two separate energy densities, since the potential energy depends on the position of both particles; this was true already in the classical case.
At first order in $`x_{1D}`$ and $`x_{2D}`$, the Schrödinger equation (89) may be decoupled in the two separate equations:
$$\frac{𝒫_i}{x_S}=m\frac{\rho _i}{t}i=1,2$$
(96)
which are the natural extensions of the condition (36) to the two-particle case. This decoupling of the Schrödinger equation confirms that solutions satisfying to both (91) and (92)-(93) may indeed exist.
At this point it is clear what should be the definition of quantum trajectories in the two-particle case: we will choose those solutions which satisfy both (91) and (92)-(93) and which remain well localized around their center of mass as $`t\pm \mathrm{}`$; this localization condition applies separately to both functions $`\rho _1`$ and $`\rho _2`$. These quantum trajectories represent individual physical systems, while positive superpositions of quantum trajectories represent statistical ensembles; in our approach, all other solutions are devoid of physical meaning. Of course, to be physically acceptable, our quantum trajectories must also satisfy the three fundamental principles, i.e. unitarity, Ehrenfest theorem and energy conservation.
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# Diblock Copolymer Ordering Induced by Patterned Surfaces Above the Order-Disorder Transition
## I Introduction
The bulk properties of diblock copolymers (BCP) are now well understood . These long linear macromolecules composed of two incompatible sub-chains, or blocks, cannot phase separate because of the covalent bond between them. This connectivity, together with the incompatibility between the two blocks, gives rise to the appearance of microphase separated phases. The state of segregation is controlled by $`N\chi `$ and $`f`$, where $`\chi `$ is the Flory parameter, $`N=N_A+N_B`$ is the total number of constituents monomers per chain and $`f`$ is the fraction of the $`A`$ block, $`f=N_A/\left(N_A+N_B\right)`$. For large enough $`N\chi `$ one of the ordered phases, such as the lamellar, hexagonal or cubic phases is preferred, depending on the degree of asymmetry $`f`$.
Less understood is the interfacial behavior of copolymer melts near solid surfaces or at the free surface with air. Surface phenomena of BCP may enable creating and controlling technologically important devices of characteristic size comparable to the wavelength of light. As examples we mention waveguides, light-emitting diodes and other optoelectronic device, anti-reflection coating for optical surfaces and dielectric mirrors .
The presence of a wall in a BCP system leads to new energy and length scales, depending on the specific chemical interaction of the polymers with the surface. In a semi-infinite system in contact with a single planar wall, the morphology near the surface can be very different from the bulk morphology. Fredrickson has considered BCP in contact with a surface having a uniform preferential adsorption to one of the two blocks. Above the order-disorder transition (ODT), where $`\chi <\chi _c`$ ($`\chi _c`$ is the critical point value of $`\chi `$ above which an ordered phase appears), he used mean-field theory and found that the order parameter (being the concentration difference between the two blocks) has decaying oscillations. He showed that the oscillation periodicity depends on $`\chi `$, and tends to the bulk lamellar periodicity as $`\chi \chi _c`$. In the same $`\chi \chi _c`$ limit, the correlation length $`\xi `$ of the oscillations was found to diverge. Further investigations showed that the inclusion of higher order, nonlinear corrections to the mean-field theory results in a non-diverging $`\xi `$. For the same system cooled below the ODT, modulated sinusoidal behavior was found. In a related work a Ginzburg-Landau free energy was used to describe the propagation of a surface-induced lamellar ordering into a bulk hexagonal phase. In the strong-segregation limit a lamellar region of finite thickness close to the surface becomes stable, provided that the surface field is larger than some critical value.
The situation is even more complicated in thin films, where the distance between the two boundaries, associated with the film thickness, is comparable to the periodicity of modulations in the bulk, and the surface induced morphology can be of different symmetry than that of the bulk. For a system taken in one of its ordered phases (below the ODT), the free energy has a local minimum when the spacing between the surfaces is an integer multiple of the bulk repeat period. The mean-field behavior of BCP close to surfaces and for BCP films was calculated using a method applicable in both the strong and weak segregation limits. It was found that confinement of lamellar phase BCP may lead to parallel layering, or in some cases even to a perpendicular arrangement of the lamellae. Self-consistent field theory (SCF) was used to study the stability of these parallel, perpendicular and mixed lamellar phases in thin films of BCP. The latter phase consists of parallel lamellae near one surface and perpendicular lamellae near the opposite surface, but it was found to be unstable for symmetric A-B ($`f=1/2`$) diblock copolymers.
So far, we mentioned situations where the surfaces have a uniform preference to one of the two blocks. More complex, chemically patterned surfaces break the lateral translation symmetry. Different surface regions will now have a different preference for the A/B blocks, thereby inducing a lateral structured morphology near the surface. Very few works took into account this possibility of a non-uniform surface. In particular, Petera and Muthukumar have investigated the effect of a one dimensional sinusoidal surface pattern on BCP morphologies close to the surface in the weak-segregation limit, both below and above the ODT.
In this paper we consider a BCP melt above the ODT near a surface, whose pattern is truly arbitrary in two dimensions, generalizing the results of Refs. . A Ginzburg-Landau free energy is expressed in term of the polymer concentration is presented in Sec. II. In Sec. III we consider a melt close to one surface or confined between two surfaces whose chemical pattern has one dimensional symmetry. Minimization of the free energy expansion gives rise to an Euler-Lagrange equation for the order parameter. A natural generalization to two-dimensional surface patterns is then considered in Sec. IV. We are able to give a complete description of the order parameter in terms of all the $`q`$–modes of the surface pattern. Finally, conclusions and some future prospects are presented in Sec. V.
## II The model
The copolymer melt is described by the order parameter $`\varphi (𝐫)`$, defined as $`\varphi (𝐫)=\varphi _A(𝐫)f`$, the difference in local A monomer concentration from its average value. Hereafter we restrict the treatment to the symmetric $`f=1/2`$ case, following the same coarse-grained free energy as was used by Fredrickson and Binder :
$$\frac{N}{k_BT}F=\left\{\frac{1}{2}\varphi \left[\tau +h\left(^2+q_0^2\right)^2\right]\varphi +\frac{u}{4!}\varphi ^4\right\}\mathrm{d}^3𝐫$$
(1)
Where $`k_B`$ is the Boltzmann constant and $`T`$ is the temperature. The other parameters are:
$`q_01.9456/\sqrt{R_g^2}`$ (2)
$`\tau =2\rho N\left(\chi _c\chi \right)`$ (3)
$`\chi _c=10.495/N`$ (4)
$`h=1.5\rho c^2R_g^2/q_0^2`$ (5)
The fundamental wavelength of the system, $`q_0`$, is expressed by $`R_g`$, the radius of gyration of the chains. The chain density $`\rho `$ is equal to $`1/Na^3`$ for an incompressible melt, and $`u/\rho `$ and $`c`$ are of order unity. More details can be found in Ref. and extensions for asymmetric BCP, $`f1/2`$ are possible as well. The use of (1) limits our treatment to a region of the phase diagram close enough to the critical point where the expansion in powers of $`\varphi `$ and its derivatives is valid, but not too close to it, because then critical fluctuation effects may be important .
This and similar types of free energy has been used to describe bulk and surface phenomena in amphiphilic systems , diblock copolymers , Langmuir films and magnetic (garnet) films . The $`\varphi ^2`$ and $`\varphi ^4`$ terms appear in the usual Landau expansion. The added $`\varphi ^2\varphi `$ and $`\varphi ^2^2\varphi `$ terms compete to produce modulated phases below the order-disorder temperature. This free energy describes a system in the disordered phase ($`\varphi =0`$, $`f=1/2`$) for $`\chi <\chi _c`$, and in the lamellar phase for $`\chi >\chi _c`$. The $`q=q_0`$ mode goes critical first, and the lamellar phase is described by $`\varphi =\varphi _q\mathrm{cos}(𝐪_\mathrm{𝟎}𝐫)`$, of repeat period $`d_02\pi /q_0`$. This single-mode approximation is accurate to order $`(\chi \chi _c)^{1/2}`$ and can be justified near the critical point . Far from the critical point higher harmonics are needed to describe the lamellar phase. As the asymmetry in composition is increased, other ordered phases of hexagonal and cubic symmetries become more stable than the lamellar phase.
As stated above, block copolymers exhibit complex surface behavior characterized by the strength and range of the interaction between the polymer chains and the surface, the typical size of chemical heterogeneities of the surface, and the distance between the two surfaces, in case of a thin film.
The presence of chemically interacting confining walls is modeled by an added short-range surface coupling term in the free energy,
$$F_s=\mathrm{d}^2𝐫_𝐬\left(\sigma (𝐫_𝐬)\varphi (𝐫_𝐬)+\tau _s\varphi ^2(𝐫_𝐬)\right)$$
(6)
The vector $`𝐫_𝐬`$ define the position of the confining surfaces. The $`\sigma \varphi `$ term expresses the preferential interaction of the surface with the A and B blocks. For example, if $`\sigma >0`$ then the B block ($`\varphi <0`$)is attracted to the surface more than the A block ($`\varphi >0`$). Control over the specificity of this surface term can be achieved by coating the substrate with carefully prepared random copolymers . The coefficient of the $`\varphi ^2`$ term in (6), $`\tau _s`$, is a surface correction to the Flory parameter $`\chi `$ . $`\tau _s>0`$ corresponds to a suppression of surface segregation of the A and B monomers.
We first consider systems in which the polymer melt is confined by a flat, rigid wall at $`y=0`$, with the $`x`$-axis chosen in the plane of the wall, and is translational invariant along the $`z`$-direction. Extension to the system of two parallel surfaces located at $`y=\pm L`$ is straightforward and will be considered later. The order parameter $`\varphi `$ is expected to vanish in regions where the interfacial interactions can be neglected,
$$\underset{y\mathrm{}}{lim}\varphi (x,y)=0$$
(7)
recovering the value $`\varphi =0`$ of the bulk phase far from the surface. In the next section we find profile solutions $`\varphi (x,y)`$ for a BCP system at temperatures above the bulk ODT.
## III One dimensional surface patterns
For high enough temperatures, or equivalently, for $`\chi <\chi _c`$, the phase of lowest free energy is the homogeneous disordered phase, with $`\varphi (𝐫)=0`$ in the bulk. The presence of a patterned surface induces ordering in the copolymer melt. If the chemical surface composition is uniform, one of the monomers, A or B, will be attracted to the surface, resulting in a parallel orientation of the lamellae (a perpendicular orientation of the chains). If the pattern is modulated, say sinusoidally, then different blocks are attracted to different regions of the surface. We will start with this case of a semi-infinite system bounded by one rigid, flat surface, and then proceed to describe thin film systems between two surfaces.
### A One patterned surface
Consider the semi-infinite BCP melt at $`y>0`$ bounded by a flat surface given by $`y=0`$. We assume a one dimensional periodic surface pattern and write it in terms of the Fourier components of the surface field $`\sigma (x)`$
$$\sigma (x)=\underset{q}{}\sigma _q\mathrm{e}^{iqx}$$
(8)
where $`\sigma _q`$ set the amplitude of the respective $`q`$–modes. The order parameter $`\varphi (x,y)`$ satisfies the boundary conditions on the surface and approaches the bulk solution far from the surface. It is convenient to decompose $`\varphi `$ in terms of its $`q`$–modes in the $`x`$-direction
$$\varphi (x,y)=\underset{q}{}f_q(y)\mathrm{e}^{iqx}$$
(9)
The requirement (7) leads to the bulk boundary condition
$$\underset{y\mathrm{}}{lim}f_q(y)=0$$
(10)
The form (9) is substituted in (1). Above the order-disorder transition (ODT) temperature, the theory is stable to second order in $`\varphi `$, and therefore the $`\varphi ^4`$ term is neglected. Using the explicit $`x`$-dependence of $`\varphi `$ in (9) we perform the $`x`$ and $`z`$ integration, yielding the free energy $`F`$:
$`F`$ $`=`$ $`{\displaystyle }{\displaystyle \underset{q}{}}\{(\tau +hq_0^4)f_qf_q^{}+hq_0^2[f_q(f_q^{^{\prime \prime }}q^2f_q)^{}+c.c.]`$ (11)
$`+`$ $`{\displaystyle \frac{1}{2}}h[f_q(f_q^{^{\prime \prime \prime \prime }}2q^2f_q^{^{\prime \prime }}+q^4f_q)^{}+c.c.]\}\mathrm{dy}`$ (12)
$`+`$ $`{\displaystyle \underset{q}{}}\left(\sigma _qf_q^{}(0)+\tau _sf_q(0)f_q^{}(0)\right)+c.c.`$ (13)
where $`(\mathrm{})^{}`$ indicates complex conjugation ( $`c.c.`$ ) operation. A standard minimization technique is carried on and yields the governing linear ordinary differential equation for the functions $`\{f_q\}`$ for $`y>0`$:
$$\left(\tau /h+\left(q^2q_0^2\right)^2\right)f_q+2(q_0^2q^2)f_q^{^{\prime \prime }}+f_q^{^{\prime \prime \prime \prime }}=0$$
(14)
This equation possesses four independent solutions in the form of an exponential $`\mathrm{e}^{k_qy}`$, with $`k_q`$ found from the characteristic equation
$$\left(\tau /h+\left(q^2q_0^2\right)^2\right)+2(q_0^2q^2)k_q^2+k_q^4=0$$
(15)
Thus, with the semi-infinite geometry, the solution is
$$f_q(y)=A_q\mathrm{e}^{k_qy}+B_q\mathrm{e}^{k_q^{}y}$$
(16)
and
$$k_q^2=q^2q_0^2+i\left(\tau /h\right)^{1/2}$$
(17)
Each $`q`$-mode solution $`f_q`$ is characterized by two complex amplitudes $`\{A_q,B_q\}`$. From the solutions of Eq. (15) wavevectors $`\{k_q\}`$ with negative real value, $`\mathrm{Re}(k_q)<0`$, are discarded, to comply with the boundary condition (10). Note that $`\mathrm{Re}(k_q)`$ is increasing monotonously as a function of $`q`$. A large value of $`\mathrm{Re}(k_q)`$ means short decay length, and hence the smallest surface $`q`$–mode decays the least. This behavior is demonstrated on Fig. 1 showing the $`q`$ and $`\chi `$ dependence of the real and imaginary parts of the wave-vector $`k_q`$. For a fixed value of $`\chi `$ the decay length, proportional to $`1/\mathrm{Re}(k_q)`$, decreases as $`q`$ increases, while the wavelength $`2\pi /\mathrm{Im}(k_q)`$ of the modulations in $`f_q(y)\mathrm{e}^{k_qy}`$ increases.
The boundary conditions for the functions $`f_q`$ can be determined by considering the Euler-Lagrange equation for $`\{f_q\}`$ in the range that includes $`y=0`$. In this case a term proportional to the Dirac delta function $`\delta (y)`$ appears in (14). There are two conditions relating $`f_q`$ and its derivatives at $`y=0`$:
$`f_q^{^{\prime \prime }}(0)+2\left(q_0^2q^2\right)f_q(0)`$ $`=`$ $`0`$ (18)
$`{\displaystyle \frac{2\sigma _q}{h}}+{\displaystyle \frac{4\tau _s}{h}}f_q(0)+2(q_0^2q^2)f_q^{^{}}(0)+f_q^{^{\prime \prime }}(0)`$ $`=`$ $`0`$ (19)
Recently, it has been found that surface states exist even in the absence of a surface field $`\sigma `$. This effect can be attributed to a loss of entropy close to the surfaces. However, in our linear theory this does not happen, and the copolymer response is proportional to the surface field $`\sigma `$. The case where $`\sigma _0`$ is a non-zero constant and $`\sigma _{q0}=0`$, corresponds to the special case of uniform interfacial interactions. The system exhibits a decaying lamellar layering of the polymers, with the B–polymer adsorbed to the surface if $`\sigma _0>0`$.
Close to the ODT, the complex wavevector $`k_0`$ can be approximated by
$$k_0\frac{\left(\tau /h\right)^{1/2}}{2q_0}+iq_0\left(1\frac{\tau /h}{8q_0^4}\right)$$
(20)
This expression shows that $`f_0\mathrm{e}^{k_0y}\mathrm{e}^{y/\xi }`$ has a diverging characteristic length $`\xi \left(\chi _c\chi \right)^{\frac{1}{2}}`$, while the oscillatory part has a wavelength slightly longer than that of the bulk lamellar phase, in agreement with the results of Fredrickson for chemically uniform surfaces . The correlation length $`\xi `$ diverges for small composition oscillations because a linear theory is employed; addition of the $`\varphi ^4`$ term in $`F`$ would give a finite value of $`\xi `$. As a result of the assumed short-range surface interactions, the periodicity and decay length of $`f_q`$ depend only on properties of the bulk, and not on surface details.
Using the notation $`k_q=k_q^{^{}}+ik_q^{^{\prime \prime }}`$ the real part of the form (16) can be rewritten as:
$`2\mathrm{R}\mathrm{e}\left(f_q\right)=\left(A_q+B_q^{}\right)\mathrm{e}^{k_qy}+c.c.`$ (21)
$`=2|A_q+B_q|\mathrm{e}^{k_q^{^{}}y}\mathrm{cos}(k_q^{^{\prime \prime }}y+\alpha _q)`$ (22)
where $`\alpha _q`$ is the phase of the $`q`$-mode modulation. It determines the value of the $`f_q`$ solution at the boundary, $`y=0`$.
The phase $`\alpha _0`$ for the $`q=0`$ mode is found to satisfy the following relation
$$\mathrm{tan}\alpha _0=\frac{q_0^2}{\sqrt{\tau /h}}$$
(23)
and therefore is determined by the degree of segregation $`\chi `$, but not by the pattern amplitude $`\sigma _q`$. A plot of $`\alpha _q`$ as a function of $`q`$ for several values of $`\chi `$ is shown in Fig. 2 (a). Far from the ODT point and deep into the disordered phase, $`\chi \chi _c`$, we find that all $`q`$-modes have a phase angle $`\alpha _q=0`$. As the ODT is approached, the $`q=q_0`$ mode retains its value but larger $`q`$-modes have a negative phase while smaller $`q`$-modes have a positive phase. At $`\chi =\chi _c`$ this becomes a step function, with $`\alpha _q=\pi /2`$ for $`q<q_0`$ and $`\alpha _q=\pi /2`$ for $`q>q_0`$. The amplitude behavior is shown in Fig. 2 (b), for the same series of segregation values $`\chi `$. As the ODT is approached, the $`q=q_0`$ mode becomes critical first, with a diverging amplitude.
An interesting limit occurs when $`\sigma _0=0`$, that is, the average surface interaction is zero (no net adsorption). No lamellar ordering parallel to the surface is expected. Indeed, the resulting checkerboard behavior is illustrated for a surface pattern chosen for simplicity to contain only one mode: $`\sigma (x)=\sigma _q\mathrm{cos}(qx)`$. Fig. 3 depicts alternating A-rich (white) and B-rich (black) regions. In (a) the decay length $`\xi `$ is smaller than in (b), because in the former case the surface periodicity is twice as large. The oscillatory behavior, characterized by $`\mathrm{Im}(k_q)`$, has a very long wavelength, diverging as $`(\chi _c\chi )^{1/2}`$ close to the ODT point.
Usually, if no special measures are taken , there is a net preference to one of the monomers: $`\sigma _00`$. The BCP morphology where the surface interactions were chosen to have both a non-zero average preference and undulatory character, namely $`\sigma =\sigma _0+\sigma _q\mathrm{cos}(qx)`$, is shown in Fig. 4. A smooth crossover from surface-induced ordering at small distance to the bulk disorder occurs. The parallel lamellae resulting from the $`\sigma _0`$ term persist farther from the surface than the bulges resulting from the $`\sigma _q`$ term, as $`f_0`$ decays slower than $`f_q`$. For a given $`\sigma _0`$, having a higher $`q`$-mode or reducing the modulation strength $`\sigma _q`$ will enhance the lamellar features far from the surface.
### B Two patterned surfaces
Until now the BCP melt was assumed to be bounded by one surface at $`y=0`$. In this section we extend our analysis to a thin-film system confined between two parallel surfaces located at $`y=L`$ and $`y=L`$, shown in Fig. 5. When the distance $`2L`$ between the surfaces is comparable to the natural bulk periodicity, the two surfaces interact via the BCP and the resulting film morphology can be very different from that of the one-surface case (Sec. III A). However, the mathematical analysis is almost the same; one only has to apply different boundary conditions on the BCP order parameter $`\varphi `$.
The surfaces at $`y=\pm L`$ are assumed to carry different surface fields of the form $`\sigma ^\pm (x)=_q\sigma _q^\pm \mathrm{e}^{iqx}`$. Only small modifications must be included to adjust the results of the previous section. The same ansatz (9) for $`\delta \varphi `$ is used. The functions $`f_q`$ that minimize the appropriate $`x`$-averaged free energy (1) obey
$$f_q(y)=A_q\mathrm{e}^{k_qy}+B_q\mathrm{e}^{k_q^{}y}+C_q\mathrm{e}^{k_qy}+D_q\mathrm{e}^{k_q^{}y}$$
(24)
with $`\{k_q\}`$ given by the same relation (17). However, unlike the semi-infinite bulk one-surface case (16), both signs of the $`k`$ vectors are used because the system is finite in the $`y`$-direction. In addition, repeating the procedure outlined for the one-surface case, gives the boundary conditions for $`f_q`$:
$`f_q^{\prime \prime }(\pm L)+2\left(q_0^2q^2\right)f_q(\pm L)`$ $`=`$ $`0`$ (25)
$`{\displaystyle \frac{2\sigma _q^\pm }{h}}+{\displaystyle \frac{4\tau _s}{h}}f_q(\pm L)`$ (26)
$`2\left(q_0^2q^2\right)f_q^{}(\pm L)f_q^{\prime \prime \prime }(\pm L)`$ $`=`$ $`0`$ (27)
We consider now several specific surfaces. In the first $`\sigma ^+=1`$ is a constant, while $`\sigma ^{}=\mathrm{cos}(qx)`$ is purely sinusoidal and average to zero, as depicted in Fig. 6. As expected, the B polymer (in black) is attracted to the upper surface, while the bottom surface exhibits modulated adsorption pattern. Although lamellar features are seen near the top surface, the overall apparent phase in the sample cannot be classified as such. The corresponding plots of the functions $`f_0(x)`$ and $`f_q(x)`$ are shown in Fig. 7 (same parameters as in Fig. 6). In general $`f_q`$ is nonzero even at the $`y=L`$ surface, although the surface does not induce modulations by itself, $`\sigma ^+=const`$. Thus modulations propagate from one surface to the other by the copolymer melt. This is an interesting observation which may have relevance in applications. It relies on the relative small thickness of the BCP film.
The situation as depicted in Fig. 6 represents a competition between two mechanisms. The modulated pattern at the bottom surface induces a laterally modulated pattern of the BCP, while the top surface uniform interaction induces a lamellar-like layering of the copolymers. As the modulated adsorption pattern strongly depends on the modulation wavenumber, so does the resulting morphology. This effect is shown explicitly in Fig. 8, where the top surface is uniform and the bottom is modulated, for a series of $`q/q_0`$ values. The transition from a locally perpendicular (bottom patterned surface) to a locally parallel orientation (at the top uniform surface) is seen in (a), similar to the so-called T-junctions between grains of different orientations . Similar behavior was found by the SCF calculation in Ref. .
Figure. 9 (a) shows the spatial dependence of the BCP order parameter when the two surfaces contain only one $`q`$-mode and are patterned in phase with each other (symmetric arrangement), but with opposite signs, $`\sigma ^\pm =\pm \sigma _q\mathrm{cos}(qx)`$. The copolymer patterns create a perfect checkerboard arrangement and are related to each other at the surfaces by an interchange of monomers A$``$B. The surface pattern (8) contains only $`\mathrm{cos}(qx)`$ terms. A generalization that includes $`\mathrm{sin}(qx)`$ sinusoidally varying modes is straightforward. In this case the patterns at the surfaces can be out–of–phase with each other. Figure. 9 (b) shows such a morphology, for $`\sigma ^+=\sigma _q\mathrm{cos}(qx)`$, $`\sigma ^{}=\sigma _q\mathrm{sin}(qx)`$, where there is a $`\pi /2`$ phase shift between the two surface fields. The perfect checkerboard arrangement of 9 (a) is now distorted to accommodate this phase shift.
## IV Two-dimensional surface patterns
So far we considered a melt in contact with a surface or confined between two surfaces of one dimensional symmetry. In our approximation, like in any linear response theory, there is no $`q`$-mode coupling proportional to $`\sigma _{q_1}\sigma _{q_2}`$. This fact allows us to go further and introduce a two-dimensional generalization of the surface pattern, which so far was taken to be independent on $`z`$. The surface now is assumed to carry a chemical pattern $`\sigma (x,z)`$ which can be written as:
$$\sigma (x,z)=\underset{q_x,q_z}{}\sigma _{q_x,q_z}\mathrm{e}^{i\left(q_xx+q_zz\right)}$$
(28)
The “linear response” function is then
$$\delta \varphi (x,y,z)=\underset{q_x,q_z}{}f_{q_x,q_z}(y)\mathrm{e}^{i\left(q_xx+q_zz\right)}$$
(29)
Because $`f`$ and $`\sigma `$ are real functions, $`f_{q_x,q_z}=f_{q_x,q_z}^{}`$ and similarly for $`\sigma `$. Using the above form it is possible to carry out the integration of the free energy in the $`xz`$ plane. Denoting $`\mathrm{}_{xz}`$ as the average in the $`xz`$ plane, it can be checked, for example, that
$$\varphi ^2\varphi _{xz}=\underset{q_x,q_z}{}f_{q_x,q_z}\left(f_{q_x,q_z}^{\prime \prime }\left(q_x^2+q_z^2\right)f_{q_x,q_z}\right)^{}$$
(30)
Defining $`𝐪_{_{}}(q_x,q_z)`$ and performing the free energy minimization with respect to $`f_{q_x,q_z}^{}`$, the functions $`f_q_{}=f_{q_x,q_z}`$ obey the same master equation (14) that $`f_q`$ previously obeyed, with the only change that $`q^2`$ is replaced by $`q_{_{}}^2`$. For a BCP in contact with a single surface, the appropriate boundary conditions are:
$`f_q_{}^{\prime \prime }(0)+2\left(q_0^2q_{_{}}^2\right)f_q_{}(0)`$ $`=`$ $`0`$ (31)
$`{\displaystyle \frac{2\sigma _q_{}}{h}}+{\displaystyle \frac{4\tau _s}{h}}f_q_{}(0)++2(q_0^2q_{}^2)f_q_{}^{}(0)+f_q_{}^{\prime \prime \prime }(0)`$ $`=`$ $`0`$ (32)
The solution for $`f_q_{}`$ is analogous to (16),
$$f_q_{}(y)=A_q_{}\mathrm{e}^{k_q_{}y}+B_q_{}\mathrm{e}^{k_q_{}^{}y}$$
(33)
$$k_q_{}^2=q_{_{}}^2\frac{1}{2}+i\left(\chi _c\chi \right)^{1/2}$$
(34)
Having found the response of the polymers to the surface modes $`\sigma _q_{}`$, one is able to deduce the concentration profiles for any given two-dimensional surface pattern. In order to illustrate this, we take a system of chemical affinity in the shape of the letters “BCP” on the $`y=0`$ surface \[Fig. 10 (a)\], and calculate the polymer concentration in the planes parallel and above it. All sizes are expressed in terms of $`d_0`$, the lamellar fundamental periodicity. The shape of the letters continuously deforms as one moves away from the surface. Contour plots corresponding to planes parallel to the $`xz`$ surface and separated by a distance $`(n+1/2)d_0`$, for integer $`n`$, are approximately given by an $`AB`$ interchange of monomers. This is the characteristic distance at which the polymers flip. Note that Fig. 10 (c) and (e) are approximately the inverse image (“negative”) of (b), (d) and (f). The original features are completely washed away as the distance $`y`$ from the surface is further increased. In our case for $`5d_0y6d_0`$ where the surface pattern size was roughly $`d_0`$.
Figure 10 also illustrates circumstances where a certain surface pattern is transferred via the bulk BCP to another, distant surface. It may be important to know, for example, if the contrast of the distant image can be experimentally detected. This reduction of the contrast is clearly seen by comparing Fig. 10 (b) and (f), and can easily be calculated from our expressions. The lamellar order created parallel to the edges of the letters in Fig. 10 (b) is the result of the undulatory nature of the block copolymers: order extending perpendicular to the surface induces order in the direction parallel to it.
The copolymer melt can follow the surface pattern when its size is larger than the polymer length-scale $`d_0`$. The effect of reducing the size of the surface structure is seen as a blurred morphology in Fig. 11 (a), where the “BCP” pattern was chosen to have dimensions $`4d_0\times 4d_0`$, compare to $`20d_0\times 20d_0`$ of Fig. 10 (b). The effect of raising the temperature (further away from the ODT) is seen in Fig. 11 (b). It is similar to Fig. 10 (b), only that the temperature is higher, $`N\chi =9`$, and the lamellar features along the edges of the letter are less prominent. Again, using our order parameter expressions one can quantify the $`q`$-mode spectrum and contrast of the distant image, as a function of the original surface pattern $`\sigma (x,z)`$, temperature and distance from the $`y=0`$ surface.
In a thin film, creation of truly three dimensional, complex morphologies between the two surfaces can be achieved by using only one-dimensional surface patterns. As an example we choose a simple sinusoidal pattern on each of the $`y=\pm L`$ surfaces, rotated 90 degrees with respect to one another: $`\sigma ^{}=\mathrm{cos}(qx+qz)`$ and $`\sigma ^+=\mathrm{cos}(qxqz)`$. In Fig. 12(b) the resulting morphology in the $`y=0`$ mid-plane is shown and is a superposition of the two surface patterns. Because $`y=0`$ is a symmetric plane, the pattern has a square symmetry. More complex patterns can be created at different $`y`$ planes.
In Sec. III we showed that the $`q=0`$ mode of the surface pattern is the slowest decaying mode, resulting in a lamellar layering parallel to the surface as $`y\mathrm{}`$, no matter what the surface pattern is. We demonstrate this in Fig. 13, where in (a) we choose a simple one-dimensional structure in the shape of a stripe of width $`d_0`$. Inside the stripe of width $`d_0`$, $`\sigma (x,z)=0.5`$ while outside it, the surface area is neutral, $`\sigma =0`$. Thus, the B-polymer is preferentially adsorbed onto the stripe. The order parameter contour plot in the $`xy`$ plane is shown in (b). It can be seen that the “surface disturbance” is enclosed with alternating lamellae. The distorted lamellae close to the $`y=0`$ surface appear curved, and slowly fade away as the distance from the surface is increased.
A different scenario is presented in Fig. 14, where inside the stripe of thickness $`d_0`$, $`\sigma =0.5`$ as above, but outside the stripe the surface is not neutral: $`\sigma =0.5`$. We find that the adsorption on the surface is quite different than in Fig. 13. Far from the stripe, the A-polymer is adsorbed onto the surface and induces stacking of the BCP in a direction parallel to the surface. Close to the surface perturbation (the stripe) the behavior is altered as the lamellae are strongly deformed in order to optimize their local interaction with the surface stripe.
## V Conclusions
We have employed a simple Ginzburg-Landau expansion of the BCP free energy to study analytically the confinement effects of block copolymers between two patterned surfaces as well as the interfacial behavior of a BCP close to a patterned surface. Our approach consists of finding the governing equation for a presumably small perturbation to the bulk order parameter, by retaining second-order terms in the free energy. This approach can be justified in the vicinity of the critical point. Above the ODT it gives rise to a simple linear equation with fixed coefficients . A generalization to two-dimensional surface patterns is presented, where a complete spatial description of the polymer concentration is given in terms of an arbitrary surface pattern. However, this approach applies to systems below the ODT as well, where a linearization is to be taken around an ordered phase .
The assumption that the surface interactions are strictly local means that the length scales of the polymer morphology are determined by bulk properties. Moreover, each of the surface $`q`$-modes in $`\sigma (x)=_q\sigma _q\mathrm{cos}qx`$ gives rise to a corresponding mode $`f_q\mathrm{cos}qx`$ in the local polymer concentration $`\varphi (x,y)`$. This “response” mode is characterized by a single wavevector $`k_q`$. The wavevector $`k_q`$ is determined by $`\chi `$ and the surface wavenumber $`q`$. In Fig. 1 we show the dependence of $`k_q`$ on these parameters. The high $`q`$-modes of the surface pattern $`\sigma (x,z)`$ decay more rapidly than those of low $`q`$. For high $`q`$-modes of characteristic length scale much smaller than the polymer chains $`d_0`$, the BCP melt cannot follow the surface modulations, and feels just the average of those modulations (which is zero for $`q>0`$). This dependence of $`k_q`$ on $`q`$ and $`\chi `$ is very similar to the results found by Petera and Muthukumar using a different free-energy functional.
Moreover, we generalized surface patterns to any two-dimensional patterns as can be seen in Fig. 10. Even within a mode decoupled (linear response) theory, many interesting effects follow for a single surface as well as for films confined between two surfaces. Tuning a few surface parameters can lead to controlled micro-structures of the BCP film.
For a BCP melt in contact with a homogeneous surface, a decaying lamellar order appears. The phase $`\alpha _0`$ of these sinusoidally damped oscillations obeys $`\mathrm{tan}\alpha _0=q_0^2/\sqrt{\tau /h}`$, and it is independent of the surface pattern amplitude $`\sigma _0`$ . For high temperatures all $`q`$-modes have the same phase $`\alpha _q=0`$. As $`\chi \chi _c`$, the $`q=q_0`$ retains this value, while higher $`q`$-modes tend to $`\pi /2`$, and lower $`q`$ tend to $`\pi /2`$. At the same limit the $`q=q_0`$ mode gets critical first, with a diverging amplitude.
Our expressions for the spatial dependence of the order parameter on a general patterned surface gives a complete description of the system, and allows for the calculation of free energy, pressure, etc. It may also help in tuning the required distance between the two surfaces in various applications. Using a strong enough surface field or fixing the conditions close to the ODT, one can hope to transform a pattern from one surface to the other surface. We also demonstrate in Fig. 12 how the superposition of simple one dimensional patterns can bring about a three dimensional behavior in a thin film system. A desired complex phase can then be achieved by tuning the Flory parameter and the relevant distances.
Possible extension to this work will be to calculate the phase diagram of the $`L_{}`$ phase (a confined lamellar phase where lamellae are perpendicular to the confining surfaces) vs. the $`L_{}`$ phase (where the lamellae are parallel to the surfaces), by calculating the surface contribution to the bulk free energy in $`F`$ (1). In the weak segregation limit this contribution is important and may lead to a completely different diagram than that of the strong segregation regime.
###### Acknowledgements.
We would like to thank S. Herminghaus, G. Krausch, M. Muthukumar, R. Netz, G. Reiter, T. Russell, M. Schick and U. Steiner for useful discussions. Partial support from the U.S.-Israel Binational Foundation (B.S.F.) under grant No. 98-00429 and the Israel Science Foundation founded by the Israel Academy of Sciences and Humanities — centers of Excellence Program is gratefully acknowledged. * Fig. 1: The real (a) and imaginary (b) parts of the wavevector $`k_q`$ as a function of the modulation $`q`$–mode and the Flory parameter $`N\chi `$. For values of $`N\chi `$ close to its critical value, $`N\chi _c=10.495`$, and for small $`q`$, $`\mathrm{Re}(k_q)`$ is small. As $`q`$ increases $`\mathrm{Re}(k_q)`$ starts to increase rapidly and $`\mathrm{Im}(k_q)`$ decreases. The value and magnitude of this sharp change in $`k_q`$ are determined by the proximity to ODT. Farther from the critical point (smaller $`\chi <\chi _c`$) the variation of $`k_q`$ with $`q`$ are smoothed out.
* Fig. 2: (a) A plot of the phase angle $`\alpha _q`$ from Eq. (21), as a function of the surface $`q`$-mode, for different Flory parameters $`N\chi `$. Circles, dotted, diamond and dashed lines correspond to $`N\chi =6.9`$, $`8.5`$, $`9.3`$ and $`10.2`$, respectively. The solid line is for $`N\chi =N\chi _c`$. Far from the ODT point (high temperatures, $`\chi \chi _c`$), all $`q`$-modes have phases equal to zero, creating a quarter-lamella region of adsorption near the surface. In the opposite limit, i.e. when $`\chi \chi _c`$, $`q`$-modes with $`q<q_0`$ have $`\alpha _q\pi /2`$, while the $`q`$-modes with $`q>q_0`$ have $`\alpha _q\pi /2`$.
In (b) are shown the surface amplitudes $`|A_q+B_q^{}|`$ from Eq. (21), as a function of the $`q`$-mode, for the same series of $`\chi `$ values as in (a). As $`\chi \chi _c`$, the $`q=q_0`$ mode goes critical first, with a diverging amplitude.
* Fig. 3: A contour plot of the BCP order parameter $`\varphi (x,y)`$, where the surface pattern (bottom line, $`y=0`$) contains only one mode: $`\sigma (x)=\sigma _q\mathrm{cos}(qx)`$, with $`\sigma _q=1`$. A-rich regions are black while B-rich are white. In (a), $`q=q_0`$ while in (b) $`q=0.5q_0`$. The decay length $`\xi `$ is smaller in (a) than in (b) because the surface $`q`$ mode is larger. The Flory parameter was set to $`N\chi =10.4<N\chi _c`$.
* Fig. 4: A contour plot of the BCP order parameter close to the critical point ($`N\chi =10.4`$). The surface pattern at $`y=0`$ is $`\sigma (x)=\sigma _0+\sigma _q\mathrm{cos}(qx)`$, where $`\sigma _0`$ is the average preference and $`q=\frac{2}{3}q_0`$ is the modulation $`q`$-mode, with amplitudes $`\sigma _0=\sigma _q=0.1`$. The $`q=0`$ surface mode has a longer range effect than the $`q>0`$ surface mode, and induces parallel lamellar arrangement farther away from the surface. At yet larger distances the order parameter decays to its bulk $`\varphi =0`$ value.
* Fig. 5: A sketch of a thin-film system confined between two surfaces. The $`y`$ axis is perpendicular to the two parallel surfaces, located at $`y=\pm L`$. The $`z`$ axis is out of the plane of the paper.
* Fig. 6: Polymer order parameter for a system of homogeneous interactions $`\sigma ^+=1`$ at $`y=L=1.5d_0`$ surface and modulated interactions at the opposite surface $`\sigma ^{}=\mathrm{cos}(qx)`$, $`y=L`$. The surface modulation wavenumber was chosen to be $`q=0.5q_0`$. As expected, the B polymer (shown in black) is preferentially attracted to the upper surface, while the bottom surface exhibits modulated adsorption pattern. This pattern propagates to the top surface. The Flory parameter was chosen $`N\chi =10.4`$.
* Fig. 7: The two amplitude functions $`f_0(y)`$ (dashed line) and $`f_q(y)`$ (solid line) from Fig. 6 plotted against $`y/d_0`$. $`f_0`$ is negative at $`y=L`$ (top uniform surface of Fig. 6), and $`f_q`$ is negative at the opposite modulated surface, $`y=L`$. Notice that the pattern at $`y=L`$ induces order at the vicinity of the other surface, as $`f_q(L)0`$, although $`\sigma ^+=const.`$
* Fig. 8: A system of modulated surface $`\sigma ^{}=\mathrm{cos}(qx)`$ at $`y=L=1.5d_0`$ and of uniform $`\sigma ^+=1`$ at the opposite surface, $`y=L=1.5d_0`$, for a series of different values of $`q/q_0`$. The effect of changing the repeat period $`q`$ is clearly seen when $`q/q_0`$ varies: $`1`$ in (a), $`2/3`$ in (b), $`1/3`$ in (c). In all cases, $`N\chi =10.2`$.
* Fig. 9: The crystalline-like checkerboard character of polymer order parameter. In (a) two patterned surfaces in phase with one another, but opposite in sign: $`\sigma ^+=\sigma ^{}=\sigma _q\mathrm{cos}(\frac{1}{2}q_0x)`$. In (b) the patterns are $`\pi /4`$ out of phase: $`\sigma ^+=\sigma _q\mathrm{cos}(\frac{1}{2}q_0x)`$, $`\sigma ^{}=\sigma _q\mathrm{sin}(\frac{1}{2}q_0x)`$. The amplitude is $`\sigma _q=0.2`$, $`N\chi =10.4`$, and the top and bottom surfaces are located at $`y=\pm 1.25d_0`$.
* Fig. 10: Propagation of surface order into the bulk. (a) is the original chemical pattern (the letters “BCP”) at the $`y=0`$ surface, whose size is $`20\times 20`$ in units of $`d_0`$. White corresponds to A-block preferring regions. A sequence of contour plots for $`y=0.5`$, $`2d_0`$, $`3.5d_0`$, $`5d_0`$ and $`6.5d_0`$ are shown in (b), (c), (d), (e) and (f), respectively. The original pattern is gradually fading (small features, high $`q`$-modes first) as $`y`$ is increased, until it is completely washed out. For $`y(n+\frac{1}{2})d_0`$ with $`n`$ integer, there is an inversion of the original pattern, as the A (white) and B block (black) are interchanged relatively to the original pattern. The Flory parameter is taken as $`N\chi =9.5`$.
* Fig. 11: Contour plots as in Fig. 10, but in (a) the surface pattern is reduced to smaller size of about $`4d_0\times 4d_0`$, while in (b) size is as in Fig. 10 but the temperature is higher, $`N\chi =9`$. The lamellar features along the letter are less prominent than in Fig. 10 (b). Note that in (b) the bulk ordering cannot tightly follow the surface pattern when the pattern size becomes comparable to $`1d_0`$, as in (a).
* Fig. 12: Creation of a complex three dimensional morphology by superposition of two one-dimensional surface patterns: $`\sigma ^{}=\mathrm{cos}(qx+qz)`$ and $`\sigma ^+=\mathrm{cos}(qxqz)`$. Shown is the thin film BCP morphology, where (a) is the surface pattern at the $`y=L=d_0`$ surface and (c) is the pattern at $`y=L=d_0`$. A contour plot of the order parameter in the mid-plane, $`y=0`$, is shown in (b), depicting A-rich and B-rich regions with square symmetry. The Flory parameter is taken as $`N\chi =9`$ and $`q=q_0/3`$.
* Fig. 13: Appearance of curved lamellae as a result of a one-dimensional surface pattern along the $`z`$ surface axis. In (a) is surface stripe is shown in the $`xz`$ plane. The (white) stripe has a surface field of $`\sigma =0.5`$ inducing preferential adsorption of the B-polymer. The rest of the surface (denoted by a grey color outside the stripe) has $`\sigma =0`$ and is indifferent to A/B adsorption. (b) is a contour plot in the $`xy`$ plane, depicting curved lamellae surrounding the “disturbance” at the middle. As the distance from the stripe is increased more than $`10d_0`$ shown here, the lamellae gradually fade away. The Flory parameter is taken to be $`N\chi =10`$.
* Fig. 14: Same as in Fig. 13, but with $`\sigma =0.5`$ inside the stripe (white), while the rest of the surface (black) has $`\sigma =0.5`$. In (a) the $`xz`$ surface is shown while in (b) the contour plots are shown in the $`xy`$ plane. Far from the stripe the B-polymer (in black) is adsorbed to the surface, and overall a lamellar morphology parallel to the surface is seen. Close to the stripe disturbance these lamellae are modified, distorted locally by the presence of the stripe.
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# DFF 361–7–00 Density fluctuations and multifragmentation of nuclear matter
## I Introduction
Semiclassical kinetic equations for the one–body phase–space density provide a powerful tool for studying the dynamics of complex processes occurring in heavy ion collisions . However, these equations in their original version give a deterministic description for the evolution of the one–body phase–space density and their solution represents the mean value of this density at each time. Thus, they are not able to account for phenomena such as the nuclear multifragmentation observed in heavy ion collisions. In this process fluctuations about the mean phase–space density are believed to play an essential role (for a review on nuclear multifragmentaion see, e.g., Ref.). In the last decade, to remedy for this drawback, an extension of the transport theory has been proposed . This approach, that incorporates a fluctuating stochastic term into the kinetic equation, is usually known as the Boltzmann–Langevin equation, and was originally applied to the treatment of hydrodynamic fluctuations in the theory of classical fluids . In Refs. collisions between nucleons in the nuclear medium are regarded as random processes and the diffusion coefficient of the Langevin (fluctuating) term is ultimately related to the amplitude of the nucleon–nucleon scattering. The last step is a particular case of the fluctuation–dissipation theorem. More recently, a new method to take into account fluctuations has also been proposed . In the approach of Ref. the statistical fluctuations of the one–body phase–space density are directly introduced by assuming local thermodynamic equilibrium. The white–noise nature of the stochastic term is the basic assumption, generally shared by all works on this subject. The authors of Ref., instead, have introduced an extension of the Boltzmann–Langevin theory by including a coloured–noise term in the stochastic force, the occurrence of such a term has been ascribed to the finite nucleon–nucleon collision time. Actual applications to nuclear problems of this interesting approach still have to be made.
In the present paper, we study the density fluctuations and their time evolution by introducing a self–consistent stochastic field acting on the constituents of the system. The self–consistency condition is provided by the fluctuation–dissipation theorem. The evolution of the fluctuations is treated within a linear approximation for the stochastic field. For simplicity we consider an infinite homogeneous system. First, without introducing any particular assumption, we prove that a withe–noise stochastic field cannot satisfy the self–consistency condition in general. Then, with reference to infinite nuclear matter, and within a collisionless mean–field approximation, we specify the particular conditions under which the withe–noise assumption for the stochastic field can be retained. These conditions are fulfilled for values of density and temperature lying in the proximity of the boundary of the spinodal region in the phase diagram, both inside and outside this region. Thus we consider nuclear matter in this physical situation and are able to solve the stochastic equation for the density fluctuations in a closed form.
With respect to the previous works on this subject , here we consider a different source of fluctuations: Landau damping. This source is present even when collisions have a negligible role in the evolution of the system.
In Sec. II we propose a procedure to determine the structure of domains formed within the system during a spinodal decomposition. If the fragmentation phenomenon observed in heavy ion collisions can be ascribed to a spinodal decomposition of the bulk of nuclei , we are allowed to identify the pattern according to which nuclear matter is decomposing, with the fragmentation pattern, and can compare the results of our nuclear matter calculations with the fragment distribution observed in heavy ion collisions. This comparison is made in Sec. III. Finally in Sec. IV a brief summary and conclusions are given.
Many papers, both theoretical and experimental, have been devoted to the multifragmentation problem. Here we mention only a few theoretical ones representing different approaches. In the statistical models of Refs., a complete statistical equilibrium of all degrees of freedom is assumed in a freeze–out volume and the various exit channels are sorted according to their statistical weight in the microcanonical ensemble. In Ref. instead nuclear multifragmentation has been described in terms of ”reducibility” and ”thermal scaling”. This means that fragments are emitted practically independently of each other and the one–fragment probability is given by a Boltzmann factor. In the dynamical models of Refs. clusters are constructed from the one–body phase–space density governed by the Boltzmann-Langevin equation . In order to take into account the quantal nature of the system and the requirement of antisymmetrization, the Quantum Molecular Dynamics model and the more sophisticated Fermionic (Antisymmetrized) Molecular Dynamics model have been developed. In addition, percolation and lattice–gas models have also been introduced. These models are particularly suitable to deal with the critical phenomena which can be expected to occur in multifragmantation. To conclude this non–exhaustive survey, we mention the calculations of Ref. that are based on the classical nucleation theory.
Even if it should eventually turn out that multifragmentation must be ascribed to very complicated processes, we think that in any case our present approach could give some insight into the underlying mechanism.
## II Formalism
### A White–noise assumption
The mean–field approximation allows us to obtain a self–consistent equation for the time evolution of the one–body density. We assume that the time scale of the terms neglected in the mean–field approximation is shorter than the characteristic times of mean–field dynamics. In order to take into account thermodynamic fluctuations, quantum effects and short–range correlations, we add to the mean–field a stochastic term similar to the random force in the Langevin equation. We assume that this additional field is a gaussian white noise with vanishing mean. In this case the time–evolution of the density is a markovian process. We will determine the conditions in which the white–noise assumption can be considered reasonable.
The additional stochastic mean field will induce density fluctuations with respect to the mean density. To be more specific, we assume that at the time $`t=0`$ in the system is present a density fluctuation $`\delta \varrho (𝐫,t=0)`$, with $`\delta \varrho (𝐫,t)=\varrho (𝐫,t)\varrho _0`$, ( $`\varrho _0`$ is the density of the reference homogeneous state i.e the state towards which the system relaxes ). Within a linear approximation for the stochastic mean–field the Fourier coefficients of $`\delta \varrho (𝐫,t)`$ for $`t>0`$ are given by ( see for example Ref., Sec. 15 I )
$$\delta \varrho _𝐤(t)=\delta \varrho _𝐤(t=0)\frac{\delta \varrho _𝐤(t=0)}{D_k(\omega =0)}_0^tD_k(tt^{})𝑑t^{}+_0^tD_k(tt^{})B_𝐤(t^{})𝑑W_𝐤(t^{}),$$
(1)
where $`D_k(tt^{})`$ is the response function of the nuclear medium and $`D_k(\omega )`$ its time Fourier transform. For symmetry reasons $`D_k(tt^{})`$ and its Fourier transform depend only on the magnitude of the wawe vector. In the second integral $`B_𝐤(t^{})dW_𝐤(t^{})`$ gives the contribution of the stochastic field in the interval $`dt^{}`$. The real and imaginary parts of the Fourier coefficients $`W_𝐤(t^{})`$ are indipendent components of a multivariate Wiener process . The fact that the stochastic field is real requires $`B_𝐤^{}(t)=B_𝐤(t)`$ and $`W_𝐤^{}(t)=W_𝐤(t)`$.
The stochastic part of the mean field is completely determined once the coefficients $`B_𝐤(t)`$ are known. In order to gain information about these coefficients we concentrate on the correlations of density fluctuations at equilibrium. Due to the independence of the components of the multivariate Wiener process $`W_𝐤(t)`$, only the terms with $`𝐤^{}=𝐤`$ survive. Within a linear approximation, these correlations can be expressed by means of the same quantities that appear in Eq.(1). This does, in a sense, correspond to the Onsanger hypothesis about the decay of deviations from equilibrium . The equation for the equilibrium fluctuations is obtained from Eq.(1) by moving the initial time to $`\mathrm{}`$, without including any particular condition at finite times. Then the correlations are given by the equation ( the brackets denote ensemble averaging )
$$<\delta \varrho _𝐤(t)\delta \varrho _𝐤(t^{})>=_{\mathrm{}}^{min(t,t^{})}𝑑t_1D_k(tt_1)D_k(t^{}t_1)B_𝐤(t_1)B_𝐤(t_1).$$
(2)
The time–translation invariance of the left–hand side of Eq.(2) requires that the coefficients $`B_𝐤(t)`$ must be constant. Taking the Fourier transform, Eq.(2) gives
$`<\delta \varrho _𝐤(\omega )\delta \varrho _𝐤(\omega ^{})>`$ $`=`$ $`2\pi \delta (\omega +\omega ^{})<(\delta \varrho _𝐤\delta \varrho _𝐤)(\omega )>`$ (3)
$`=`$ $`2\pi \delta (\omega +\omega ^{})D_k(\omega )D_k(\omega ^{})|B_𝐤|^2.`$ (4)
By exploiting the fluctuation–dissipation theorem
$$<(\delta \varrho _𝐤\delta \varrho _𝐤)(\omega )>=\frac{2}{1e^{\beta \omega }}\mathrm{Im}D_k(\omega ),$$
(5)
where $`\beta =1/T`$ is the inverse temperature, ( we use units such that $`\mathrm{}=c=k_B=1`$ ), we obtain for the coefficients $`B_𝐤`$ the equation
$$|B_𝐤|^2=\frac{2}{1e^{\beta \omega }}\frac{\mathrm{Im}D_k(\omega )}{|D_k(\omega ))|^2}.$$
(6)
We have used the relation $`D_k(\omega )=D_k^{}(\omega )`$. Equation (6) can be satisfied only if the right–hand side does not depend on $`\omega `$. This can occur only in particular situations, thus the original white–noise assumption about the stochastic mean–field is not correct in general. This result is quite general, so we can conclude that for a perturbed system approaching an equilibrium state, fluctuations about the average trajectory cannot usually be accounted for by a white–noise stochastic force.
Now, with reference to symmetric nuclear matter, we discuss particular physical situations in which Eq.(6) can have a solution. Only in such conditions the assumption of a withe–noise stochastic field is valid. The relevant quantity is the linear–response function $`D_k(\omega )`$. We evaluate $`D_k(\omega )`$ within a self–consistent mean–field approximation. In order to derive compact analytical expressions, here we use the linearized Vlasov equation for calculating the response function. This equation can be regarded as a semiclassical approximation to the random phase approximation, valid in the longwavelength limit. We also use a Skyrme–like form of the nucleon–nucleon effective interaction. Our self–consistent mean–field potential is given by
$$U=a\frac{\varrho }{\varrho _{eq}}+b\left(\frac{\varrho }{\varrho _{eq}}\right)^{\alpha +1}d^2\varrho ,$$
(7)
where $`\varrho _{eq}`$ is the saturation density of nuclear matter. For the parameters in Eq.(7) we take the values:
$`a=356.8\mathrm{MeV},b=303.9\mathrm{MeV},\alpha ={\displaystyle \frac{1}{6}},`$
$`d=130\mathrm{MeV}\mathrm{fm}^5.`$The values of $`a`$, $`b`$ and $`\alpha `$ reproduce the binding energy ( $`15.75\mathrm{MeV}`$ ) of nuclear matter at saturation ($`\varrho _{eq}=0.16\mathrm{fm}^3`$) and give an incompressibility modulus of $`201\mathrm{MeV}`$. For the values of $`d`$ we follow the prescriptions of Ref..
The response function is given by
$$D_k(\omega )=\frac{D_k^{(0)}(\omega )}{1𝒜_kD_k^{(0)}(\omega )},$$
(8)
where $`D_k^{(0)}(\omega )`$ is the non–interacting particle–hole propagator, and
$$𝒜_k=a\frac{1}{\varrho _{eq}}+\frac{b}{\alpha +1}\frac{1}{\varrho _{eq}^{\alpha +1}}\varrho _0^\alpha +dk^2$$
(9)
are the Fourier coefficients of the effective interaction. Here $`\varrho _0`$ is the density of the reference homogeneous state.
By substituting the expression (8) for $`D_k(\omega )`$ into Eq.(6), we obtain
$$|B_𝐤|^2=\frac{2}{1e^{\beta \omega }}\frac{\mathrm{Im}D_k^{(0)}(\omega )}{|D_k^{(0)}(\omega ))|^2}.$$
(10)
This equation shows that the coefficients $`B_𝐤`$ do not explicitely depend on the nucleon–nucleon effective interaction. However, we remark that $`D_k^{(0)}(\omega )`$ is the propagator of independent particles that are moving in the mean–field of the reference homogeneous state,
$$U_0=a\frac{\varrho _0}{\varrho _{eq}}+b\left(\frac{\varrho _0}{\varrho _{eq}}\right)^{\alpha +1},$$
(11)
thus the interaction between constituents does enter, although not explicitly, into the expression of $`B_𝐤`$.
We shall now show that in the classical limit $`\omega /T1`$, the right–hand side of Eq.(10) does not depend on $`\omega `$, thus the assumption of a white–noise stochastic field can be considered valid in that limit.
In the actual physical situations considered in this paper the values of the temperature are small enough with respect to the Fermi temperature so that the Pauli principle is still operating. Therefore the strength of the particle–hole excitations having energies much higher than $`kv_F`$$`v_F`$ is the Fermi velocity ) can be considered negligible. Moreover the relevant values of the wave vector $`k`$ turn out to be such that the quantity $`kv_F`$ is of the same order of magnitude as $`T`$. Thus the limit $`\omega /T1`$ also implies $`\omega /kv_F1`$.
The non interacting particle–hole propagator $`D_k^{(0)}(\omega )`$ acquires a very simple form in the longwavelength (Vlasov) limit. The imaginary part is
$$\mathrm{Im}D_k^{(0)}(\omega )=\frac{1}{4\pi }m^2\frac{\omega }{k}𝑑\epsilon _p\frac{n_p}{\epsilon _p}\theta \left(1\frac{\omega }{kv}\right),$$
(12)
where
$`n_p={\displaystyle \frac{4}{e^{\beta (ϵ_p\stackrel{~}{\mu })}+1}}`$is the mean occupation number of nucleons with kinetic energy $`ϵ_p=p^2/2m`$, and $`v=p/m`$. The effective chemical potential $`\stackrel{~}{\mu }`$ is measured with respect to the uniform mean field $`U_0`$.
For $`\omega /kv_F1`$ the imaginary part of $`D_k^{(0)}(\omega )`$ is given by
$$\mathrm{Im}D_k^{(0)}(\omega )=\frac{1}{\pi }m^2\frac{\omega }{k}\frac{1}{e^{\beta \stackrel{~}{\mu }}+1}+O\left((\frac{\omega }{kv_F})^3\right),$$
(13)
while the real part of $`D_k^{(0)}(\omega )`$ in the longwavelength limit takes the form
$$\mathrm{Re}D_k^{(0)}(\omega )=\frac{1}{2\pi ^2}𝑑pp^2\frac{n_p}{\epsilon _p}\left(1+\frac{1}{2}\frac{\omega }{kv}\mathrm{ln}\frac{1+\omega /kv}{|1\omega /kv|}\right).$$
(14)
For T sufficiently low with respect to $`\stackrel{~}{\mu }`$, the most important contribution to the integral in Eq.(14) comes from a small domain of $`\epsilon _p`$ around $`\stackrel{~}{\mu }`$. So we can take $`\omega /kv1`$ in evaluating the integral, and obtain
$$\mathrm{Re}D_k^{(0)}(\omega )=\frac{\varrho _0}{\stackrel{~}{\mu }}+O\left((\frac{\omega }{kv_F})^2\right).$$
(15)
With $`D_k^{(0)}(\omega )`$ given by Eqs.(13) and (15), the right–hand side of Eq.(10) is independent of $`\omega `$ to the lowest significant order in $`\omega /T`$. Thus, for $`\omega /T1`$ the magnitude of the coefficients $`B_𝐤`$ is given by
$$|B_𝐤|^2=\frac{2}{\pi }m^2(\frac{\stackrel{~}{\mu }}{\varrho _0})^2\frac{T}{e^{\beta \stackrel{~}{\mu }}+1}\frac{1}{k}.$$
(16)
The phases of $`B_𝐤`$, instead, remain unknown. However, we will see that only the quantities $`|B_𝐤|^2`$ are needed to determine the probability distribution of density fluctuations. Finally, we remark that $`|B_𝐤|`$ for a given $`𝐤`$, is determined solely by the density and temperature of nuclear matter.
The white–noise assumption is justified if the excitation strength is concentrated in a narrow range of energy close to zero. This condition requires that $`\mathrm{Im}D_k(\omega )`$ is a sharply peaked function in the proximity of $`\omega =0`$ and is negligile elsewhere. The imaginary part of $`D_k(\omega )`$ displays this feature for values of temperature and density near the borders of the spinodal region, since the pole of $`D_k(\omega )`$ lying on the imaginary axis, moves towards $`\omega =0`$ as the system approaches the mechanical instability. This is shown in Fig.1, where we report $`\mathrm{Im}D_k(\omega )`$ calculated with Eq.(8) using the complete expression of $`D_k^{(0)}(\omega )`$. With our effective interaction, for $`T=5\mathrm{MeV}`$ the spinodal region starts at $`\varrho _c=0.617\varrho _{eq}`$. The values of $`\varrho _0`$ used in Fig.1 are close to this critical value.
### B Distribution of fluctuations.
We now derive from Eq.(1) the probability distribution for $`\delta \varrho _𝐤(t)`$ in the limit $`\omega /T1`$, and for values of temperature and density in the proximity of the spinodal region. The response function $`D_k(\omega )`$ has a pole in the lower part of the imaginary $`\omega `$–axis, at a position given by
$$i\mathrm{\Gamma }_k=i\frac{\pi }{m^2}(1+e^{\beta \stackrel{~}{\mu }})\frac{\varrho _0}{\stackrel{~}{\mu }}\frac{\left({\displaystyle \frac{^2f}{\varrho _0^2}}|_T+dk^2\right)}{𝒜_k}k.$$
(17)
We have used the relation
$$\frac{\stackrel{~}{\mu }}{\varrho _0}|_T=\frac{^2f}{\varrho _0^2}|_T𝒜_0,$$
(18)
where $`f`$ is the free–energy density and $`𝒜_0=𝒜_{k=0}`$.
In Eq.(17) the relevant quantity is the isothermal stiffness $`{\displaystyle \frac{^2f}{\varrho _0^2}}|_T`$, which vanishes on the boundary of the spinodal region. Since we limit our calculations to the proximity of the spinodal region, we neglect $`{\displaystyle \frac{^2f}{\varrho _0^2}}|_T`$ with respect to $`𝒜_0`$ in evaluating $`D_k(t)`$. Furthermore, in actual calculations the typical values of $`k`$ which come into play are such that the term $`dk^2`$ is smaller than $`|𝒜_0|`$, thus we also neglect this term with respect to $`|𝒜_0|`$. This approximation is consistent with the longwavelength limit adopted in the calculation of $`D_k(\omega )`$.
Substituting into Eq.(1) the response function $`D_k(tt^{})`$ calculated with these approximations, the equation for the fluctuations $`\delta \varrho _𝐤(t)`$ becomes:
$$\delta \varrho _𝐤(t)=\delta \varrho _𝐤(t=0)e^{\mathrm{\Gamma }_kt}+\stackrel{~}{B}_k_0^te^{\mathrm{\Gamma }_k(tt^{})}𝑑W_k(t^{}),$$
(19)
where
$$|\stackrel{~}{B}_k|=\frac{1}{|𝒜_0|}\sqrt{\frac{2\pi T}{m^2}(1+e^{\beta \stackrel{~}{\mu }})k},$$
(20)
and $`\mathrm{\Gamma }_k`$ is given by Eq.(17), neglecting the term $`dk^2`$ in $`𝒜_k`$. We recall that $`\mathrm{\Gamma }_k`$ is negative, so that $`|\mathrm{\Gamma }_k|`$ represents the damping rate of fluctuations, that vanishes for long wavelengths when $`{\displaystyle \frac{^2f}{\varrho _0^2}}|_T0`$.
Equation (19) represents an Ornstein–Uhlenbeck process with $`|\mathrm{\Gamma }_k|`$ as drift coefficient and $`\stackrel{~}{B}_k`$ as diffusion coefficient. The corresponding Fokker–Planck equation for the probability distribution $`P[\delta \varrho _𝐤(t)]`$ reads
$$\frac{}{t}P[\delta \varrho _𝐤(t)]=|\mathrm{\Gamma }_k|\frac{}{\delta \varrho _𝐤(t)}\delta \varrho _𝐤(t)P[\delta \varrho _𝐤(t)]+\frac{1}{2}|\stackrel{~}{B}_k|^2\frac{^2}{\delta \varrho _𝐤^2(t)}P[\delta \varrho _𝐤(t)].$$
(21)
For simplicity we assume the state of the system at $`t=0`$ to be homogeneous on average ( $`<\delta \varrho _𝐤(t=0)>=0`$ for $`k0`$ ). Equation (19) says that this property holds during time evolution. In this case the solution of Eq.(21) is a gaussian distribution with zero mean value. Whenever it is necessary, a non vanishing mean value can easily be introduced. The explicit expression of the distribution $`P[\delta \varrho _𝐤(t)]`$ is
$$P[\delta \varrho _𝐤(t)]=N_1e^{{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐤}{}}\delta \varrho _𝐤^{}\left(t\right){\displaystyle \frac{1}{\sigma _k^2(t)}}\delta \varrho _𝐤\left(t\right)},$$
(22)
with the variance $`\sigma _k^2(t)`$ given by
$$\sigma _k^2(t)=\sigma _k^2(t=0)e^{2\mathrm{\Gamma }_kt}+\frac{T}{f^{\prime \prime }+dk^2}(1e^{2\mathrm{\Gamma }_kt}).$$
(23)
Here the constant $`N_1`$ is a normalization factor and we have introduced the abbreviation
$`f^{\prime \prime }={\displaystyle \frac{^2f}{\varrho _0^2}}|_T.`$For $`t\mathrm{}`$ Eq.(22) reproduces the usual gaussian approximation with variance
$$\sigma _k^2(t=\mathrm{})=\frac{T}{f^{\prime \prime }+dk^2}$$
(24)
for the equilibrium thermodynamical fluctuations . We furthermore remark that Eq.(23) for the time evolution of the variance is similar to that obtained with different approaches in previous works on this subject .
For later purpose we report also the distribution of the fluctuations in ordinary space:
$$P[\delta \varrho (𝐫,t)]=N_2e^{{\displaystyle \frac{1}{2}}{\displaystyle 𝑑𝐫𝑑𝐫^{}\delta \varrho (𝐫,t)M(𝐫,𝐫^{},t)\delta \varrho (𝐫^{},t)}},$$
(25)
where
$`M(𝐫,𝐫^{},t)={\displaystyle \frac{1}{V}}{\displaystyle \underset{𝐤}{}}e^{i𝐤\dot{(}𝐫𝐫^{})}{\displaystyle \frac{1}{\sigma _k^2(t)}},`$and $`N_2`$ is an appropriate normalization factor.
The diffusion coefficients of Eq.(20) are derived by means of the fluctuation–dissipation theorem, which concerns only fluctuations about equilibrium. In Ref. a way has been suggested to extend the treatment of stable cases to processes where instabilities can develop. Following that suggestion we include in our approach the case of nuclear matter merged in the spinodal region. In practice, we still assume the validity Eq.(22) for the probability distribution of the density fluctuations, with the variance $`\sigma _k^2(t)`$ of Eq.(23) calculated with the values of temperature and density of the new situation. This amounts to treating the diffusion coefficients for the unstable case as an analytic continuation of the stable–case coefficients in the $`(\varrho ,T)`$ plane. The reliability of such a procedure lies in the fact that both the growing rate $`\mathrm{\Gamma }_k`$ and the diffusion coefficient $`\stackrel{~}{B}_k`$ change smoothly when the system crosses the stability boundary and enters the spinodal region. The pole of $`D_k(\omega )`$ in turn, continuously moves along the imaginary axis from the lower part to the upper part of the complex $`\omega `$–plane ( see Eq.(17) ). In order to preserve causality, the integral for calculating the Fourier anti transform $`D_k(t)`$ must be performed along a path which cuts the imaginary axis above the pole.
In the unstable case, the time behaviour of the variance $`\sigma _k^2(t)`$ in Eq.(23) is similar to that predicted by linear theories of the spinodal decomposition of alloys and fluids ( for an extensive review on this subject see Ref. ). The variance grows exponentially for the fluctuations with wave number
$$k<k_c=\sqrt{\frac{|f^{\prime \prime }|}{d}},$$
(26)
while it tends to the aymptotic value $`\sigma _k(t=\mathrm{})`$ of Eq.(24) for $`k>k_c`$. In particular the growth rate $`\mathrm{\Gamma }_k`$ presents a maximum for $`k=k_M=k_c/\sqrt{3}`$. This means that the pattern of the regions which contain coherently correlated fluctuations is asymptotically characterized by the wavelength $`\lambda _M=2\pi /k_M`$. These features for the growth rate of unstable modes are analogous to those obtained in Ref. within a different scheme.
### C Size of fragments
Starting from the probability distribution for density fluctuations given by Eq.(25), we can determine the corresponding distribution for the size of the correlation domains. It has already been recalled that the stable and unstable cases can be treated within the same scheme. Thus we shall invesigate two different situations which could be explored by nuclear matter during a nucleus–nucleus collision: in one case the system is in the metastable region and relaxes towards a local minimum of the free energy, while in the other case the system is merged in the spinodal zone and develops density fluctuations which grow with time and will eventually lead to decomposition. According to our approximations, we limit our analysis in both cases to values of temperature and density in the proximity of the borders of the spinodal region. Moreover, we consider homogenous nuclear matter in both cases, and still assume that $`<\delta \varrho _𝐤(t=0)>=0`$ for $`k0`$.
Before performing explicit calculations, we make a few remarks. It is known that linear theories are unable to describe the late stages of the spinodal decomposition of alloys and fluids ( see Ref. and references quoted therein ). In particular, they predict a limiting value for the length scale that characterizes the pattern of the correlation domains. This value is given by the wavelength $`\lambda _M`$ for which the growth rate of fluctuations has a maximum. Instead, Monte Carlo simulations and experimental results show a continuous coarsening of the domains with increasing time. However it has been argued in Ref. that the early–time Monte Carlo results are consistent with a linear theory, provided that a stochastic force is included.
In the physical situations considered in the present paper (heavy ion collisions), the value of the characteristic wavelength $`\lambda _M`$ is larger than $`10\mathrm{fm}`$, beyond the size of the nuclear system involved. Moreover, the corresponding growth time $`1/\mathrm{\Gamma }_{k_M}`$ is of the same order of magnitude as the characteristic times of the nucleus–nucleus collisions in the energy range considered here. Thus the fluctuations with wave number $`k_M`$ are still far from being the predominant ones in this time interval. This means that the processes that we are investigating correspond to an early stage of the spinodal decomposition. Then we can expect reliable results from our approach, at least at a qualitative level.
From Eq.(25) we obtain the usual expression for the equilibrium correlation function
$$G(|𝐫𝐫^{}|)=\frac{1}{4\pi }\frac{T}{d}\frac{e^{|𝐫𝐫^{}|/\xi }}{|𝐫𝐫^{}|},$$
(27)
where
$`\xi =\sqrt{{\displaystyle \frac{d}{f^{\prime \prime }}}}`$is the correlation length. This quantity, which represents the average extension of the correlation domains, can be obtained by an appropriately weighted integral of the correlation function:
$$\xi =𝑑𝐫𝑑𝐫^{}F(𝐫,𝐫^{})G(|𝐫𝐫^{}|).$$
(28)
The function $`F(𝐫,𝐫^{})`$ is a suitable weight function. Here we extend this relation between averaged quantities to fluctuating quantities, for systems both at equilibrium and out of equilibrium. We then assume that the size of correlation domains at time $`t`$ is given by a quadratic functional of the fluctuations $`\delta \varrho (𝐫,t)`$:
$$b=L(t)\frac{𝑑𝐫𝑑𝐫^{}\delta \varrho (𝐫,t)F(𝐫,𝐫^{})\delta \varrho (𝐫^{},t)}{𝑑𝐫𝑑𝐫^{}F(𝐫,𝐫^{})G(|𝐫𝐫^{}|,t)},$$
(29)
where $`L(t)=<b>`$ is the length scale that characterizes the pattern of the domains, and $`G(|𝐫𝐫^{}|,t)`$ is the correlation function for systems out of equilibrium. The latter quantity is the space Fourier transform of the variance of Eq.(23).
In order to simplify calculations, we further choose for $`F(𝐫,𝐫^{})`$ a separable form. The requirement that $`b`$ should be positive for any function $`\delta \varrho (𝐫,t)`$, enforces a symmetric form
$$F(𝐫,𝐫^{})=f(𝐫)f(𝐫^{}).$$
(30)
of the weight function. This form allows us to obtain a closed expression for the probability distribution of $`b`$. In addition, with this choice the final results are entirely independent of the function $`f(𝐫)`$.
Now we derive the probability distribution for $`b`$ as a function of the length scale $`L(t)`$. Later we shall give a procedure for determining $`L(t)`$.
For a given probability distribution $`P[\delta \varrho (𝐫,t)]`$, the related probability distribution for $`b`$, at a given time $`t`$, can be obtained by means of the functional integral
$$P(b,t)=d[\delta \varrho (𝐫,t)]\delta \left(b\frac{L(t)}{C}𝑑𝐫\delta \varrho (𝐫,t)F(𝐫)\delta \varrho (0,t)\right)P[\delta \varrho (𝐫,t)],$$
(31)
where we have put
$$C=𝑑𝐫𝑑𝐫^{}f(𝐫)G(|𝐫𝐫^{}|,t)f(𝐫^{})$$
(32)
in order to simplify the notation.
With the distribution of Eq.(25) and using the integral representation of the $`\delta `$–function
$`\delta (x)={\displaystyle \frac{1}{2\pi }}{\displaystyle 𝑑\eta e^{i\eta x}}`$the equation for the distribution $`P(b,t)`$ takes the form
$`P(b,t)`$ $`=`$ $`{\displaystyle \frac{N_2}{2\pi }}{\displaystyle 𝑑\eta e^{i\eta b}}`$ (33)
$`\times `$ $`{\displaystyle d[\delta \varrho (𝐫,t)]e^{{\displaystyle \frac{1}{2}}{\displaystyle 𝑑𝐫𝑑𝐫^{}\delta \varrho (𝐫,t)\left(M(𝐫,𝐫^{},t)+2i\eta \frac{L(t)}{C}f\left(𝐫\right)f\left(𝐫^{}\right)\right)\delta \varrho (𝐫^{},t)}}}.`$ (35)
The functional integral is of gaussian type and allows us to express the result of the integration in closed form:
$$P(b,t)=\frac{N_2}{2\pi }𝑑\eta e^{i\eta b}\frac{1}{(det\frac{1}{2\pi }[\widehat{M}+2i\eta \frac{L(t)}{C}\widehat{F}])^{\frac{1}{2}}}.$$
(36)
The quantities $`\widehat{M}`$ and $`\widehat{F}`$ are infinite–dimensional operators, with matrix elements $`M(𝐫,𝐫^{},t)`$ and $`f(𝐫)f(𝐫^{})`$ respectively, in the coordinate representation.
The determinant in the last equation can be factorized as
$$det\frac{1}{2\pi }[\widehat{M}+2i\eta \frac{L(t)}{C}\widehat{F}]=det[\frac{\widehat{M}}{2\pi }]det[\mathrm{𝟏}+2i\eta \frac{L(t)}{C}\widehat{M}^1\widehat{F}],$$
(37)
where $`\mathrm{𝟏}`$ is the unit matrix. The square root of the first factor on the right–hand side and the normalization constant $`N_2`$ of Eq.(36) coincide and cancel. What remains to be evaluated is the inverse of the square root of the second determinant. For this purpose we write the determinant in exponential form and expand the exponent in a power series. Thus, we obtain the following formal expression
$`(det[\mathrm{𝟏}+2i\eta {\displaystyle \frac{L(t)}{C}}\widehat{M}^1\widehat{F}])^{\frac{1}{2}}`$ $`=`$ $`e^{{\displaystyle \frac{1}{2}}Tr\mathrm{ln}\left(\mathrm{𝟏}+2i\eta {\displaystyle \frac{L(t)}{C}}\widehat{M}^1\widehat{F}\right)}`$ (38)
$`=`$ $`e^{{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{n}\left(1\right)^{1+n}\left(2i\eta \frac{L(t)}{C}\right)^nTr\left(\widehat{M}^1\widehat{F}\right)^n}}.`$ (40)
We recall that the matrix element $`M^1(𝐫,𝐫^{},t)`$ and the correlation function $`G(|𝐫𝐫^{}|,t)`$ coincide. Because of the separable form chosen for the function $`F(𝐫,𝐫^{})`$, Eq.(30), the trace operation on the generic $`n`$–term of Eq.(LABEL:det) simply yields $`C^n`$. Thus the series can be resummed and gives $`\mathrm{ln}(1+2i\eta L(t))`$. Then the probability distribution $`P(b,t)`$ acquires the form
$`P(b,t)={\displaystyle \frac{1}{2\pi }}{\displaystyle 𝑑\eta \frac{e^{i\eta b}}{(1+2i\eta L(t))^{\frac{1}{2}}}}.`$A simple integration in the complex $`\eta `$ plane gives the final result
$$P(b,t)=\frac{1}{\sqrt{\pi }}\frac{1}{\sqrt{2L(t)b}}e^{b/\left(2L\left(t\right)\right)}.$$
(41)
From the probability distribution of the domain size we can derive the distribution of the number of nucleons $`A`$ that are contained in a correlation domain, assumed to be spherical. For a homogeneous liquid the relation between $`A`$ and the size $`b`$ is $`b=\mathrm{\hspace{0.17em}2}r_0A^{1/3}`$, where $`r_0`$ is determined by the actual density. With a simple transformation of variables we obtain for the probability distribution of $`A`$, $`P(A,t)`$, the equation
$$P(A,t)=\frac{1}{3}\frac{1}{\sqrt{\pi }}\sqrt{\frac{r_0}{2L(t)}}A^{5/6}e^{{\displaystyle \frac{r_0}{L(t)}}A^{1/3}}.$$
(42)
Further, to take into account that $`A`$ is a discrete variable we express the probability of finding a correlation domain containing $`A`$ nucleons, $`Y(A)`$, through the integral
$$Y(A)=_{A1}^A𝑑AP(A,t).$$
(43)
For large $`A`$, $`Y(A)`$ tends to coincide with $`P(A,t)`$.
## III Results
The distribution $`P(A,t)`$ and the probability $`Y(A)`$ are completely determined once the ratio between the length scale $`L(t)`$ and the mean interparticle spacing $`r_0`$ is fixed. The parameter $`L(t)`$ sets the scale for the decrease of the correlation function $`G(r,t)`$ with increasing $`r`$. We can obtain an estimate of $`L(t)`$ by analyzing the behaviour of $`G(r,t)`$ as a function of $`r`$ at a given $`t`$. The correlation function is initially determined by the variance $`\sigma _k^2(t=0)`$, then, in the stable case, it asymptotically assumes the form given in Eq.(27), with the appropriate correlation length $`\xi =L(t=\mathrm{})`$, while in the unstable case, it acquires a damped oscillatory behaviour carachterized by the asymptotic wavelength $`\lambda _M`$. In order to illustrate the general features of the function $`L(t)`$, we simply assume that the initial fluctuations are negligible, $`\sigma _k^2(t=0)\mathrm{\hspace{0.17em}0}`$. In this case the function $`G(r,t)`$ is completely determined by the density and temperature of nuclear matter. Here we consider two sample values for the density ( $`\varrho _0=\mathrm{\hspace{0.17em}0.65}\varrho _{eq}`$ and $`\varrho _0=0.58\varrho _{eq}`$ ) and a single value for the temperature ( $`T=5\mathrm{MeV}`$ ). This temperature is in the range of values expected for the nuclear multifragmentation process . The two corresponding points in the phase diagram $`(\varrho ,T)`$ lie in the metastable region and in the spinodal region respectively, and are sufficiently close to the boundary of the spinodal zone to justify our assumption of a white–noise stochastic field. In Figs. 2 and 3 we show the behaviour of $`G(r,t)`$ as a function of $`r`$ at three different values of time, both in the stable and unstable situations. In the stable case of Fig. 2, a simple inspection of the behaviour of $`G(r,t)`$ shows that it is reasonably well reproduced by a function like that on the right–hand side of Eq.(27) (obviously with $`\xi `$ replaced by $`L(t)`$ ). We adopt such a form for $`G(r,t)`$, then, by comparison with its true behaviour shown in Fig.2, we can determine $`L(t)`$. For the unstable case shown in Fig. 3, the situation is slightly more involved because the asymptotic regime is reached only after a very long time. For this case, we simply assume that $`L(t)`$ does coincide with the distance at which the value of $`G(r,t)`$ is reduced by $`80\%`$ with respect to its value at $`r=1\mathrm{fm}`$ (because of our approximations we cannot expect the present approach to be reliable for distances shorter than $`1\mathrm{fm}`$).
At a given time $`t`$ the value of the length scale $`L(t)`$ depends strongly on the distance from the boundary of the spinodal zone, the shorter this distance, the larger is $`L(t)`$. In Fig. 4 the calculated length $`L(t)`$ is displayed as a function of $`t`$ for the two chosen sets of parameters. The values of $`t`$ are in the range that is relevant for nuclear fragmentation . Figure 4 shows that for $`t200\mathrm{fm}/c`$, $`L(t)`$ pratically reaches its asymptotic value ($`L(\mathrm{})3.0\mathrm{fm}`$) in the metastable situation, whereas in the unstable case $`L(t)`$ is still much smaller than $`L(\mathrm{})12\mathrm{fm}`$.
In the two physical situations considered here, two different processes could drive nuclear matter towards a spinodal decomposition. In the metastable case, if the density fluctuations are large enough, the nuclear system can explore the unstable region for a time sufficiently long to move towards a phase separation. In the unstable case instead, fluctuations grow with time until they cause the decomposition of the nuclear system. In both cases we expect that the pattern of domains containing the liquid phase is determined by the probability distribution $`P(b,t)`$ or $`P(A,t)`$ of Eqs.(41) and (42).
In order to assess the degree of validity of our approach, we compare the results of our calculations with the corresponding experimental data by identifying the probability $`Y(A)`$ of Eq.(43) with the distribution of the fragment yield. Since experimentally the fragments are detected according to their charge, we have to transform $`P(A,t)`$ and $`Y(A)`$ into the corresponding functions of $`Z`$. We assume a homogeneous distribution also for the charge $`Z={\displaystyle \frac{(1\alpha )}{2}}A`$, with $`\alpha =(NZ)/A`$, and use $`\alpha =\mathrm{\hspace{0.17em}0.2}`$, which corresponds to the average asymmetry of the nuclear systems considered.
In Fig. 5 the probability $`Y(Z)`$ is displayed as a function of $`Z`$ on a double logarithmic scale for three different values of the ratio $`L(t)/r_0`$. The range of values for $`L(t)/r_0`$ has been chosen in accordance with that of $`L(t)`$ in Fig. 4. Figure 5 shows that $`Y(Z)`$ can be fit with good accuracy by a power law $`Y(Z)=Y_0Z^{\tau _{eff}}`$. The values of the effective exponent, $`\tau _{eff}`$, lie between $`1.17`$ for $`L(t)/r_0=\mathrm{\hspace{0.17em}4}`$ and $`1.42`$ for $`L(t)/r_0=\mathrm{\hspace{0.17em}2}`$.
The power–law behaviour of the fragment yield and the determination of the exponent have been the subject of several experimental studies of multifragmentation ( see for example the recent papers ). The observed values of the exponent are in the interval $`1.21.5`$ for nuclear rections with beam energies lower than $`40A\mathrm{MeV}`$, whereas they exceed the value of $`2`$ at higher energies . A value of the exponent $`\tau _{eff}2`$ can be unlikely reproduced by our calculations because we would need an unreasonably low value for the ratio $`L(t)/r_0`$. However, in various papers it has been remarked that the effects of collective motions, that have not been taken into account by our present approach, should become more important with increasing beam-energy.
Figures 6 and 7 show a comparison between the charge distributions predicted by our approach, $`Y(Z)`$, and recent experimental data obtained by the Multics/Miniball collaboration for $`Au+Au`$ collisions at an incident energy of $`E=35A\mathrm{MeV}`$ and by the INDRA Collaboration for $`{}_{}{}^{129}Xe+Sn`$ and $`{}_{}{}^{155}Gd+^{238}U`$ collisions at $`E=32A\mathrm{MeV}`$ and $`E=36A\mathrm{MeV}`$ respectively . The calculations have been performed for three values of the parameter $`L(t)/r_0`$. We have normalized the experimental distributions to one in order to perform the comparison on an absolute scale. We can see that the agreement between experimental data and the calculated charge distributions is quite satisfactory for $`Z<30÷35`$ and that for the lighter fragments the experimental points are better reproduced with larger values of the ratio $`L(t)/r_0`$. For $`Z>30÷35`$ the observed distribution presents a slope steeper than that predicted by our calculations. This faster decrease should be ascribed to finite–size effects which have not been included in our nuclear matter treatment.
## IV Summary and conclusions
We have studied the density fluctuations associated with a one–body treatment of nuclear dynamics. In our approach the fluctuations are generated by adding a stochastic term to the mean field. This additional random force is determined by a self-consistency condition required by the fluctuation–dissipation theorem. We have treated the effects of the stochastic field in linear approximation and this has allowed us to express the time evolution of the fluctuations in a closed form.
First we have analyzed the nature of the stochastic field and have shown that in general a white–noise assumption for the stochastic field is not consistent with the fluctuation–dissipation theorem. Then we have studied the particular physical conditions in which the white–noise nature of the stochastic term can be retained. These conditions include hot nuclear matter at a temperature $`T5\mathrm{MeV}`$, where the system can be still considered degenerate. We have found that for a Fermi system the treatment of density fluctuations by means of a white–noise stochastic term is justified when the limit $`\omega /T1`$ gives a reasonable approximation to the density–density response. This condition is better satisfied when the density and temperature of the system are close to the borders of the spinodal region in the ($`\varrho ,T`$) plane. Thus, in the limit $`\omega /T1`$ the equilibrium fluctuations can be adequately described by means of thermodynamic functions and we can expect that in this limit the purely quantum fluctuations will play a negligible role also for systems not too far from equilibrium. We have extended the results obtained for the probability distribution of a metastable system to unstable situations. This has been achieved by extrapolating the relevant quantities across the boundary of the spinodal region. Because of the linear approximation used for evaluating the response of the system to the stochastic force, the fluctuations have a gaussian probability distribution.
In the final part of this paper we have introduced a procedure to determine the size and mass distributions of the domains containing correlated density fluctuations, then we have compared the obtained mass distribution to the yield of light fragments observed in the multifragmentation of heavy nuclei. The procedure proposed here is quite general and can be applied to any gaussian fluctuation distribution
Our approach can account both for the observed power–law distribution and for the value of the effective exponent found experimentally, but for the exponent the agreement is limited to collisions with beam energies lower than $`40A\mathrm{MeV}`$. This discrepancy between our predictions and the observed values of the effective exponent in collisions of higher energies deserves further investigations. A detailed comparison with experiment has shown that our approach fairly reproduces the measured charge distributions for $`Z<30÷35`$. Since we are dealing with infinite nuclear matter, we expect to overestimate the number of fragments having a large fraction of the mass of the emitting source.
Finally, we remark that the obtained mass distribution contains only one parameter, the ratio between the time–dependent length scale of domains $`L(t)`$ and the mean interparticle spacing $`r_0`$. This ratio can become large. A more detailed comparison of the present model with experimental data could also give an estimate of the time required by the system to break up.
Figure captions:
Fig.1 Imaginary part of the response function $`D_k(\omega )`$ of Eq.(8) for hot nuclear matter ($`T=5\mathrm{MeV}`$) at different densities approaching the critical value $`\varrho _c=0.617\varrho _{eq}`$ (dotted line: $`\varrho _0=0.70\varrho _{eq}`$, dashed line: $`\varrho _0=0.65\varrho _{eq}`$, full line: $`\varrho _0=0.63\varrho _{eq}`$). The value of $`k`$ is $`0.1k_F`$.
Fig.2 Spatial behaviour of time–dependent correlation function $`G(r,t)`$ for nuclear matter on the stable side of the spinodal curve ($`\varrho _0=0.65\varrho _{eq}`$, $`T=5\mathrm{MeV}`$). The three curves correspond to different values of $`t`$ (full line: $`t=50\mathrm{fm}/c`$, dashed line: $`t=100\mathrm{fm}/c`$, dotted line: $`t=200\mathrm{fm}/c`$).
Fig.3 Same as Fig.2, but for the unstable case ($`\varrho _0=0.58\varrho _{eq}`$).
Fig.4 Behaviour of the length scale $`L(t)`$ within the spinodal region (full line: $`\varrho _0=0.58\varrho _{eq}`$) and outside it (dashed line: $`\varrho _0=0.65\varrho _{eq}`$).
Fig.5 Fragment distribution $`Y(Z)`$ calculated for different values of the ratio $`L(t)/r_0`$. Full line: $`L(t)/r_0=4`$, dashed line: $`L(t)/r_0=3`$, dotted line: $`L(t)/r_0=2`$.
Fig.6 Comparison of fragment distribution $`Y(Z)`$ calculated for $`L(t)/r_0=6`$ (full line), $`4`$ (dashed line), $`2`$ (dotted line) with experimental distribution for the reaction $`Au+Au`$ at $`35A\mathrm{MeV}`$. Data from Ref. have been normalized to one.
Fig.7 Same as Fig.6, but for the reactions $`{}_{}{}^{129}Xe+Sn`$ at $`E=32A\mathrm{MeV}`$ (triangles) and $`{}_{}{}^{155}Gd+^{238}U`$ at $`E=36A\mathrm{MeV}`$ (circles). Data from Ref. have been normalized to one.
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# 1 Introduction
## 1 Introduction
Quantum field theory on a noncommutative space has been proved to be useful in understanding various physical phenomena, like as various limits of M-theory compactification , low energy effective field theory of D-branes with constant Neveu-Schwarz $`B`$-field background , and quantum Hall effect . Although noncommutative field theories are non-local, they appear to be highly constrained deformation of local field theory. Thus it may help understanding non-locality at short distances in quantum gravity.
Noncommutative field theory means that fields are thought of as functions over noncommutative spaces. At the algebraic level, the fields become operators acting on a Hilbert space as a representation space of the noncommutative space. Since the noncommutative space resembles a quantum phase space (with noncommutativity $`\theta `$ playing the role of $`\mathrm{}`$), it exhibits an interesting spacetime uncertainty relation, which causes a $`UV/IR`$ mixing and a teleological behavior . Also, for nonzero $`\theta `$, there can be nonperturbative effects in the form of soliton solutions even at the classical level and it could not have a smooth limit. Indeed, several such solutions have recently appeared .
In this Letter we will show the existence of nonperturbative solutions in the $`CP(n)`$ model on noncommutative two plane. The $`CP(n)`$ model, even though it consists only of scalar fields, enjoys local gauge invariance and exhibits many similarities to instantons in four-dimensional Yang-Mills theory, such as the existence of self-dual soliton solution with scale and orientation parameters . In addition, it has many applications to condensed matter systems . However, since the $`CP(n)`$ model is relatively simpler than four-dimensional noncommutative Yang-Mills theory, it will be very useful “toy model” to investigate various questions of noncommutative gauge theory if the ordinary $`CP(n)`$ model can be generalized to noncommutative space. We will demonstrate here it is the case.
In Section 2, we construct the consistent $`CP(n)`$ model on noncommutative plane. In Section 3, the Bogomolny bound on the energy is considered and it is shown that it is saturated by (anti-)self-dual solutions. In Section 4, we solve the soliton solutions for the (anti-)self-dual equations and show that their topological charges, which are independent of their scaling and orientation, are quantized. However, we point out that this integer quantization is satisfied only for the field configurations without any singularity in the commutative sense. Also we argue that our solution is the most general BPS solution. In Section 5, we summarize our results with a brief discussion of the possible implication to the instanton physics in Yang-Mills theories on noncommutative $`𝐑^4`$ , including the problem for noncommutative instanton solutions discussed in .
## 2 CP(n) Model
We consider the (2+1)-dimensional field theory whose space is noncommutative two plane. The coordinates $`x,y`$ of this noncommutative plane satisfy the relation
$$[x,y]=i\theta $$
(1)
with $`\theta >0`$. This noncommutative plane has not only the translation symmetry but also rotational symmetry. One can see that the parity operation $`(x,y)(x,y)`$ is broken on noncommutative plane. The classical field on this noncommutative space is an element $`\mathrm{\Phi }(t,x,y)`$ in the algebra $`𝒜_\theta `$ defined by $`x,y`$ with the relation $`(\text{1})`$.
Introduce the complex coordinates
$$z=\frac{x+iy}{\sqrt{2}},\overline{z}=\frac{xiy}{\sqrt{2}},$$
(2)
which satisfy
$$[z,\overline{z}]=\theta >0.$$
(3)
This commutation relation is that of the creation and annihilation operators for a simple harmonic oscillator and so one may use the simple harmonic oscillator Hilbert space $``$ as a representation of (1). The ground state is $`|0>`$ such that $`z|0>=0`$ and $`|n>=\overline{z}^n/\sqrt{\theta ^nn!}|0>`$ so that
$$z|n>=\sqrt{\theta n}|n1>,\overline{z}|n>=\sqrt{\theta (n+1)}|n+1>.$$
(4)
The integration over noncommutative two plane becomes the trace over its Hilbert space, which respect the translation symmetry:
$$d^2x𝒪\mathrm{Tr}𝒪=2\pi \theta \underset{n0}{}<n|𝒪|n>.$$
(5)
The $`CP(n)`$ manifold is defined by an $`(n+1)`$-dimensional complex vector $`\mathrm{\Phi }=(\varphi _1,\varphi _2,\mathrm{},\varphi _{n+1})`$ of unit length with the equivalence relation under the overall phase rotation $`\mathrm{\Phi }e^{i\alpha }\mathrm{\Phi }`$ . This complex projective space of real dimension $`2n`$ is equivalent of the coset space $`U(n+1)/U(1)\times U(n)`$. (It is quite straightforward to generalize our consideration here to the general Grassmanian models with the manifold $`G(n,m)=U(n)/U(m)\times U(nm)`$ .)
Here we are interested in the $`CP(n)`$ model on the noncommutative plane. Thus the field variable $`\mathrm{\Phi }(t,x,y)`$ becomes an operator acting on $``$. The spatial derivatives are
$$_x\mathrm{\Phi }=i\theta ^1[y,\mathrm{\Phi }],_y\mathrm{\Phi }=i\theta ^1[x,\mathrm{\Phi }].$$
(6)
The natural Lagrangian for the $`CP(n)`$ model turns out to be
$$L_1=\mathrm{Tr}\left(_\mu \mathrm{\Phi }^{}^\mu \mathrm{\Phi }+(\mathrm{\Phi }^{}_\mu \mathrm{\Phi })(\mathrm{\Phi }^{}^\mu \mathrm{\Phi })\right)$$
(7)
with the constraint
$$\mathrm{\Phi }^{}\mathrm{\Phi }1=0.$$
(8)
This theory has a global $`U(n+1)`$ symmetry and a local $`U(1)`$ symmetry
$$\mathrm{\Phi }(x)\mathrm{\Phi }(x)g(x)$$
(9)
which removes the degrees of freedom for the overall phase of $`\mathrm{\Phi }`$. The $`U(1)`$ gauge transformation acts on the right side, which leaves the constraint (8) invariant. This ordering of the gauge transformation is the key point which makes the whole model work.
This Lagrangian with the constraint (8) can be rewritten as
$$L_2=\mathrm{Tr}\left(D_\mu \mathrm{\Phi }^{}D^\mu \mathrm{\Phi }+\lambda (\mathrm{\Phi }^{}\mathrm{\Phi }1)\right)$$
(10)
with
$$D_\mu \mathrm{\Phi }=_\mu \mathrm{\Phi }i\mathrm{\Phi }A_\mu ,$$
(11)
where $`A_\mu (x)`$ is the $`U(1)`$ gauge field without its kinetic term and $`\lambda (x)`$ is a Lagrangian multiplier to incorporate the constraint (8). This Lagrangian is invariant under the local gauge transformation defined by (9) and
$$A_\mu (x)g^{}A_\mu gig^{}_\mu g.$$
(12)
As there is no gauge kinetic term, one can solve the $`A_\mu `$ equation to get
$$A_\mu =i\mathrm{\Phi }^{}_\mu \mathrm{\Phi }$$
(13)
which shows that the gauge transformations (9) and (12) are consistent with one another. After using Eqs. (8) and (13), the second Lagrangian (10) becomes the first Lagrangian (7), as it should be.
Note that $`\mathrm{\Phi }^{}D_\mu \mathrm{\Phi }=0`$ and the field strength is
$`F_{\mu \nu }`$ $`=`$ $`_\mu A_\nu _\nu A_\mu +i[A_\mu ,A_\nu ]`$ (14)
$`=`$ $`i(D_\mu \mathrm{\Phi }^{}D_\nu \mathrm{\Phi }D_\nu \mathrm{\Phi }^{}D_\mu \mathrm{\Phi })`$ (15)
which is the curvature tensor $`[D_\mu ,D_\nu ]\mathrm{\Phi }=i\mathrm{\Phi }F_{\mu \nu }`$. From the field equation for $`\mathrm{\Phi }`$,
$$D_\mu D^\mu \mathrm{\Phi }\mathrm{\Phi }\lambda =0,$$
(16)
we can deduce
$$\lambda =\mathrm{\Phi }^{}D_\mu D^\mu \mathrm{\Phi }=D_\mu \mathrm{\Phi }^{}D^\mu \mathrm{\Phi },$$
(17)
and the field equation becomes
$$D_\mu D^\mu \mathrm{\Phi }+\mathrm{\Phi }(D_\mu \mathrm{\Phi }^{}D^\mu \mathrm{\Phi })=0.$$
(18)
## 3 Energy Bound
As in the commutative case, the $`CP(n)`$ model on the noncommutative plane has the Bogomolny energy bound. The conserved energy functional becomes
$`E`$ $`=`$ $`\mathrm{Tr}(D_0\mathrm{\Phi }^{}D_0\mathrm{\Phi }+D_i\mathrm{\Phi }^{}D_i\mathrm{\Phi }))`$ (19)
$`=`$ $`\mathrm{Tr}\left(|D_0\mathrm{\Phi }|^2+|D_z\mathrm{\Phi }|^2+|D_{\overline{z}}\mathrm{\Phi }|^2\right).`$
Similar to the commutative case, let us consider an inequality
$$\mathrm{Tr}\left\{(D_i\mathrm{\Phi }\pm iϵ_{ij}D_j\mathrm{\Phi })^{}(D_i\mathrm{\Phi }\pm iϵ_{ij}D_j\mathrm{\Phi })\right\}0.$$
(20)
Expending this we obtain
$$\mathrm{Tr}(D_i\mathrm{\Phi }^{}D_i\mathrm{\Phi })iϵ_{ij}\mathrm{Tr}(D_i\mathrm{\Phi }^{}D_j\mathrm{\Phi }).$$
(21)
The BPS bound on the energy is then
$$E\mathrm{Tr}(D_0\mathrm{\Phi }^{}D_0\mathrm{\Phi })+2\pi |Q|,$$
(22)
where the $`U(1)`$ gauge invariant ‘topological charge’ is
$$Q=\frac{i}{2\pi }ϵ_{ij}\mathrm{Tr}D_i\mathrm{\Phi }^{}D_j\mathrm{\Phi }=\frac{\mathrm{Tr}F_{12}}{2\pi }.$$
(23)
Contrast to the commutative case, there exists no conserved topological current. Instead, the current
$$J^\mu =\frac{1}{4\pi }ϵ^{\mu \nu \rho }F_{\nu \rho }$$
(24)
is covariantly conserved, $`D_\mu J^\mu =0`$. However, this implies that for the localized configurations, the topological charge $`Q=\mathrm{Tr}J^0`$ is conserved. In the complex coordinate,
$$Q=\frac{1}{2\pi }\mathrm{Tr}\left(|D_z\mathrm{\Phi }|^2|D_{\overline{z}}\mathrm{\Phi }|^2\right).$$
(25)
The energy bound is saturated by the configuration which is static in time and satisfies the (anti-)self-dual equation , $`D_i\mathrm{\Phi }\pm iϵ_{ij}D_j\mathrm{\Phi }=0`$, or in the complex notation,
$`D_{\overline{z}}\mathrm{\Phi }=0(\mathrm{for}\mathrm{self}\text{-}\mathrm{dual}\mathrm{case}Q>0),`$ (26)
$`D_z\mathrm{\Phi }=0(\mathrm{for}\mathrm{anti}\text{-}\mathrm{self}\text{-}\mathrm{dual}\mathrm{case}Q<0).`$ (27)
## 4 (Anti-)Self-dual Solitons
To find the (anti-)self-dual configurations, let us try to parameterize the field as follows
$$\mathrm{\Phi }=W/\sqrt{W^{}W},$$
(28)
where $`W`$ is an $`(n+1)`$-dimensional vector. Since $`\mathrm{\Phi }^{}\mathrm{\Phi }=1`$, locally one can choose finite $`W`$ such that $`\sqrt{W^{}W}(x)`$ is invertible. We also introduce an $`(n+1)`$-dimensional projection operator
$$P=1W\frac{1}{W^{}W}W^{},$$
(29)
whose kernel is one-dimensional space generated by $`W`$ vector. In terms of this field variable, the Lagrangian (7) becomes
$$L=\mathrm{Tr}\left(\frac{1}{\sqrt{W^{}W}}_\mu W^{}P^\mu W\frac{1}{\sqrt{W^{}W}}\right),$$
(30)
and the topological number is
$$Q=\frac{1}{2\pi }\mathrm{Tr}\left\{\frac{1}{\sqrt{W^{}W}}\left(_{\overline{z}}W^{}P_zW_zW^{}P_{\overline{z}}W\right)\frac{1}{\sqrt{W^{}W}}\right\}.$$
(31)
The operator in the trace can be regarded as the topological density operator on the noncommutative space. In terms of $`W`$ variable, the above Lagrangian in the commutative case has a local scaling symmetry, $`WW\mathrm{\Delta }(x)`$ with an arbitrary scalar $`\mathrm{\Delta }(x)`$. What is remarkable about the noncommutative space case is that this scaling symmetry on the Lagrangian still holds. The projection operator $`P`$ is independent of $`\mathrm{\Delta }`$ and so the Lagrangian and the topological number are still invariant. However, the multicomponent field $`\mathrm{\Phi }`$ is not invariant under the local scaling on noncommutative case, contrast to the commutative case. As far as the classical field theory is concerned, we could choose the primary field to be $`W`$ instead of $`\mathrm{\Phi }`$, and regard the classical theory is invariant under the local scaling. This would be crucial in finding the most general solution.
The self-dual equation then becomes
$$D_{\overline{z}}\mathrm{\Phi }=P(_{\overline{z}}W)(W^{}W)^{1/2}=0,$$
(32)
which is equivalent to $`_{\overline{z}}W=WV`$ for arbitrary scalar $`V`$ both for commutative and noncommutative cases. For either cases the most general solution is
$$W=W_0(z)\mathrm{\Delta }(z,\overline{z})$$
(33)
with $`(n+1)`$-dimensional holomorphic vector $`W_0(z)`$ and arbitrary scalar function $`\mathrm{\Delta }(z,\overline{z})`$. As we just argued in the previous paragraph, this arbitrariness is a local scaling and can be scaled away.
Let us consider the self-dual solutions in commutative case, which is well studied before. We choose the scaling so that the $`(n+1)`$th component of $`W`$ is chosen to be unity. Then, we get the standard self-dual equation for the $`n`$-dimensional vector $`w`$ such that $`W=(w,1)`$
$`_{\overline{z}}w=0(\mathrm{for}\mathrm{self}\text{-}\mathrm{dual}),`$ (34)
$`_zw=0(\mathrm{for}\mathrm{anti}\text{-}\mathrm{self}\text{-}\mathrm{dual}).`$ (35)
The most general solution of the above self-dual equation should be a $`n`$-dimensional vector whose components are holomorphic functions. These solutions are characterized by its topological charge $`k`$: the self-dual solutions carry positive integer charges and the anti-self-dual solutions do negative integers. The general self-dual solutions in commutative case are given in the meromorphic form,
$$w=\frac{1}{P_{n+1}(z)}(P_1(z),P_2(z),\mathrm{},P_n(z)),$$
(36)
where $`P_i(z)`$ are $`k`$th order polynomial of $`z`$. The meromorphic function is not holomorphic at poles as
$$_{\overline{z}}\frac{1}{z}=4\pi \delta ^2(z).$$
(37)
However, $`w`$ blows up at poles and so the self-dual equation (32) still holds, making the solutions (36) the most general one.
For the commutative case, the solution (32) is equivalent to the smooth solution
$$\mathrm{\Phi }=(P_1(z),P_2(z),\mathrm{},P_{n+1}(z))\frac{1}{1+_iP_i^{}P_i}$$
(38)
by a singular $`U(1)`$ gauge transformation $`|P_{n+1}|/P_{n+1}(z)`$. In this case the vector
$$W=(P_1(z),P_2(z),\mathrm{},P_{n+1}(z))$$
(39)
has components which are $`k`$th order polynormials of $`z`$ only. This solution has $`2(n+1)k+2n`$ real parameters, among which $`2n`$ are the vacuum moduli parameter for $`CP(n)`$ space and the rest of which $`2(n+1)k`$ parameters account for the size and scale parameters of $`k`$ solitons. Note that the $`W`$ vector in (39) is holomorphic everywhere.
When we goes to the noncommutative case, we should be more careful. As $`z^1=(\overline{z}z)^1\overline{z}=\overline{z}(\overline{z}z+\theta )^1`$, we get
$$zz^1=I,z^1z=I|0><0|.$$
(40)
Since $`_{\overline{z}}f(z,\overline{z})=\theta ^1[z,f(z,\overline{z})]`$,
$$_{\overline{z}}z^1=\theta ^1|0><0|,$$
(41)
and the solution of type $`1/z`$ is not holomorphic on noncommutative space. This has the analogue of (37) on noncommutative space.
In addition, we will see later that the solution $`1/z`$ will have fractional topological charge. On the commutative case, two types of solutions (36) and (38) are gauge equivalent, but that is not true in general on the noncommutative case. While $`z^1`$ is nonholomorphic, the solutions given in Eq. (38) are polynomial so holomorphic, and so they are solutions of the self-dual equation. This is the most general solution even in the noncommutative space, modulo the local scaling we considered before. Not only they satisfy the self-dual equation (32), these solutions also carry the integer topological numbers.
Let us start with a simple solution in the $`CP(1)`$ model,
$$W=(az,1),$$
(42)
and so
$$_{\overline{z}}W^{}P_zW=\frac{a^2}{1+a^2(\overline{z}z+\theta )},$$
(43)
where we have used $`zf(\overline{z}z)=f(\overline{z}z+\theta )z`$ and $`\overline{z}f(\overline{z}z)=f(\overline{z}z\theta )\overline{z}`$. Thus the topological charge is
$`Q_s`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\mathrm{Tr}\left\{{\displaystyle \frac{1}{1+a^2\overline{z}z}}\left(_{\overline{z}}W^{}P_zW\right)\right\}`$ (44)
$`=`$ $`{\displaystyle \frac{1}{2\pi }}\mathrm{Tr}\left\{{\displaystyle \frac{a^2}{(1+a^2\overline{z}z)\left(1+a^2(\overline{z}z+\theta )\right)}}\right\}.`$
With the dimensionless parameter $`s=a^2\theta `$, the trace becomes
$`Q_s`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}(2\pi \theta ){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{a^2}{(1+a^2\theta n)(1+a^2\theta (n+1))}}`$ (45)
$`=`$ $`s{\displaystyle \frac{1}{(1+sn)(1+s(n+1))}}`$
$`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{1}{1+sn}}{\displaystyle \frac{1}{1+s(n+1)}}\right)=1.`$
The scale parameter $`a`$ of the soliton can be arbitrary but the topological number does not change. Especially for the zero size soliton $`a=\mathrm{}`$, topological charge density does not vanish only for the $`|0>`$ state.
If we have used the unacceptable singular solution
$$W=(z^1,1),$$
(46)
then its topological charge becomes
$`Q_s`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\{1/(1+s(n+1))1/(1+s(n+2))\}`$ (47)
$`=`$ $`1/(1+s).`$
As we argued before, this solution is not acceptable. Notice that one can see that this solution has the topological charge less than 1. This fractional topological number contrasts with the commutative case. As on the noncommutative case $`z^1`$ is as singular as $`z`$ in the operator sense, one see that there can be a configuration with a fractional topological charge. One thus wonder whether one should include this configuration in the family of classically acceptable configurations. We do not know the answer for this. This may be answerable by considering what kind of soliton and anti-soliton pairs are created when some amount of energy is put to the vacuum.
For a single anti-soliton solution in $`CP(1)`$ model,
$$W=(a\overline{z},1),$$
(48)
we get
$$_zW^{}P_{\overline{z}}W=\frac{a^2}{1+a^2\overline{z}z}.$$
(49)
Its topological charge is then
$`Q_{\overline{s}}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\mathrm{Tr}\left\{{\displaystyle \frac{a^2}{(1+a^2(\overline{z}z+\theta ))(1+a^2\overline{z}z)}}\right\}`$ (50)
$`=`$ $`s{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(1+sn)(1+s(n+1))}}=1,`$
with $`s=a^2\theta `$.
Thus the topological index works fine for these solitons (38). This suggests that the topological charge can be calculated for arbitrary (anti-)self-dual solutions. This is indeed true as we will see now. For the general solution of Eq. (39), we can say
$$W=az^ku+𝒪(z^{k1}),$$
(51)
where $`u`$ is a $`n+1`$ dimensional unit vector. To calculate its topological charge, we first note that Eq.(31) can be rewritten as
$$Q_s=\frac{1}{2\pi }\mathrm{Tr}\left\{\sqrt{W^{}W}_{\overline{z}}\left(\frac{1}{W^{}W}W^{}_zW\right)\frac{1}{\sqrt{W^{}W}}\right\}.$$
(52)
Now we can insert the complete set of states between operators to get
$`Q_s`$ $`=`$ $`\theta {\displaystyle \underset{n,m,l}{}}<n|\sqrt{W^{}W}|m><m|_{\overline{z}}\left\{{\displaystyle \frac{1}{W^{}W}}_z(W^{}W)\right\}|l><l|{\displaystyle \frac{1}{\sqrt{W^{}W}}}|n>`$ (53)
$`=`$ $`\theta {\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}<n|_{\overline{z}}\left({\displaystyle \frac{1}{W^{}W}}W^{}_zW\right)|n>.`$
This is the analogue of the total derivative on noncommutative plane. Noting $`\theta _{\overline{z}}𝒪(z,\overline{z})=[z,𝒪]`$, we can find the integration of the total derivative as
$`{\displaystyle \frac{1}{2\pi }}\mathrm{Tr}_{\overline{z}}𝒪`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}<n|[z,𝒪]|n>`$ (54)
$`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left\{\sqrt{\theta (n+1)}<n+1|𝒪|n>\sqrt{\theta n}<n|𝒪|n1>\right\}`$
$`=`$ $`\underset{N\mathrm{}}{lim}{\displaystyle \underset{n=0}{\overset{N}{}}}\left\{<n+1|𝒪z|n+1><n|𝒪z|n>\right\}`$
$`=`$ $`\underset{N\mathrm{}}{lim}<N+1|𝒪z|N+1>,`$
assuming that $`𝒪z|0>=0`$. If $`𝒪`$ is singular and so that $`𝒪z|0>0`$, there would be additional boundary terms. (For example, the singular $`w=u/(az^k)`$ solution there is an additional boundary term.) If the large $`N`$ limit vanishes, say like $`1/N`$, then there is no boundary term and so the sum vanishes. If the limit is of order one, the limit is finite. If it diverges like a power of $`N`$, then the limit is not well defined. In this case, the series should be treated more carefully.
Using the method in (54), we get
$$Q_s=\underset{n=0}{\overset{\mathrm{}}{}}\left\{<n+1|\frac{1}{W^{}W}W^{}(_zW)z|n+1><n|\frac{1}{W^{}W}W^{}(_zW)z|n>\right\}.$$
(55)
The expectation $`<N+1|(W^{}W)^1W^{}(_zW)z|N+1>`$ is of order one. More concretely, we see that
$$\frac{1}{W^{}W}W^{}(_zW)z=\frac{1}{W^{}W}(W^{}(_zW)zkW^{}W)+k.$$
(56)
Defining that
$$\mathrm{\Omega }\frac{1}{W^{}W}\left(W^{}(_zW)zkW^{}W\right),$$
(57)
we see
$$\underset{N\mathrm{}}{lim}<N+1|\mathrm{\Omega }|N+1>=\underset{N\mathrm{}}{lim}\frac{N^{k1}}{N^k}=0,$$
(58)
as $`W^{}(_zW)zkW^{}W=𝒪((\overline{z}z)^{k1}).`$ Thus the charge becomes
$`Q_s`$ $`=`$ $`\underset{N\mathrm{}}{lim}\left\{<N+1|\mathrm{\Omega }|N+1>+k\right\}`$ (59)
$`=`$ $`k.`$
For general anti-self-dual soliton solution,
$$W=a\overline{z}^ku+𝒪(\overline{z}^{k1}),$$
(60)
similar argument leads to the topological charge $`k`$.
## 5 Concluding Remarks
In this Letter we have shown that the $`CP(n)`$ model can be also well-defined on noncommutative two plane. There exist the (anti-)self-dual solitons that saturate the BPS energy bounds, which are regular and carry integer topological charge $`k`$ with $`2(n+1)k+2n`$ real parameters. We found that (anti-)self-dual solitons carry integer topological charge regardless their orientation and size when the field configurations are regular. We have also shown that the singular solutions, which are acceptable and related to the regular solutions by gauge transformations on the commutative plane, are not acceptable on noncommutative plane. Not only they do not satisfy the self-dual equations, but also are not related to the regular solution by the gauge transformation on the noncommutative plane. As we have seen, the topological number does not change when the soliton shrinks to zero size. This should remain true after going to the commutative variables by using the Seiberg-Witten map . While it is not clear how that is achieved in our case, there is no natural way to evoke something like freckled instantons and make two space to blow up at some points, contrast to Braden and Nekrasov’s work .
The low energy dynamics of the solitons will be described on the moduli space. To do this, one has to know the metric of the moduli space . Our general solutions for $`k`$ solitons has $`2(n+1)k+2n`$ real parameters in $`CP(n)`$ model. The vacuum moduli space has $`2n`$ real parameters and their kinetic energy diverges due to the volume factor. In addition, the total scale and orientations with $`2n`$ real parameters have infinite inertia. So a single soliton with $`k=1`$ has only two parameters with finite inertia, corresponding to the position of the soliton. For $`k`$ solitons, the moduli space $`_{k,n}`$ of finite inertia has $`2(n+1)k2n`$ real dimension. The low energy dynamics of these $`k`$ solitons can be described by the metric of the moduli space. The solitons would not feel the noncommutativity of the underlying space directly. However, the moduli space of solitons on noncommutative space would be no longer singular when a single soliton collapses to a point. It would be interesting to study in detail the moduli space dynamics of solitons and compare that with those on the commmutative plane.
The solitons in the gauge theories are more complicated than the scalar theory like $`CP(n)`$. Recently there has been many works on the instantons on the noncommutative $`𝐑^4`$ . The soliton properties of the noncommutative $`CP(n)`$ model may play an important role in figuring out the subtle issues in the noncommutative instantons of four dimensional Yang-Mills theory. One of the key observation from our work is that the topological number on the noncommutative space is somewhat tricky quantity, which needs a careful treatment. Clearly one see a possible solution to the recently discovered quandary where the instanton number of a single $`U(2)`$ instanton depends on the size of the instanton.
Acknowledgment BHL is supported by the Ministry of Education, BK21 Project No. D-0055 and by grant No. 1999-2-112-001-5 from the Interdisciplinary Research Program of the KOSEF and would like to thank the hospitality in KIAS where the part of this work is done. KL is supported in part by KOSEF 1998 Interdisciplinary Research Grant 98-07-02-07-01-5 and would like to thank the hospitality in Niels Bohr Institute where the part of this work is done.
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# 𝐾_{𝐿,𝑆}→𝜋𝜋𝜈𝜈̄ Decays Within and Beyond the Standard Model
## I Introduction
Historically, rare kaon decays have provided a crucial testing ground in which to study flavor-changing neutral current (FCNC) and CP-violating phenomena. Within the context of the Standard Model, the observation of such decays gives us further information on the Cabibbo-Kobayashi-Maskawa (CKM) matrix elements and provides independent checks on the consistency of our understanding drawn from other measurements, such as those of semileptonic decays and B-meson decays. They also serve as a good place to look for physics beyond the Standard Model. Even if new physics is first observed outside the kaon sector, we will want to know the footprint it leaves in rare kaon decays and specifically on the $`sd`$ weak transition.
In this paper we focus on the decays $`K_{L,S}\pi ^o\pi ^o\nu \overline{\nu }`$ and $`K_{L,S}\pi ^+\pi ^{}\nu \overline{\nu }`$. As noted previously, at the quark level these FCNC processes involve the $`sd\nu \overline{\nu }`$ transition. In the Standard Model, this transition is dominated by short-distance contributions involving loop diagrams that contain W and Z bosons and heavy quarks. With the top quark mass much greater than that of the charm quark, the imaginary part of the amplitude for this transition arises almost entirely from loop diagrams with top quarks and the resulting amplitude is then proportional to the CKM factor Im($`V_{td}^{}{}_{}{}^{}V_{ts}`$).
The decays, $`K^+\pi ^+\nu \overline{\nu }`$ and $`K_L\pi ^o\nu \overline{\nu }`$ are both governed by the same $`sd\nu \overline{\nu }`$ transition. The latter decay is CP-violating and dominated to high accuracy by the short-distance contribution. As noted above, its amplitude is consequently proportional to Im($`V_{td}^{}{}_{}{}^{}V_{ts}`$) in the Standard Model. For both decays, the hadronic matrix elements of the relevant weak current can be related to ones which enter charged current semileptonic decays, whose direct measurement bypasses any theoretical uncertainty in the hadronic matrix element as well. The branching ratios that are predicted in the Standard Model lie between roughly $`10^{11}`$ and $`10^{10}`$. With negligible long-distance contributions and little hadronic uncertainty, these decays have been pointed to as crucial for precision experimental tests of the Standard Model and correspondingly as places to look for new physics if the Standard Model fails. Because of this, they are being pursued experimentally in spite of very difficult experimental backgrounds.
The rare decays $`K_{L,S}\pi \pi \nu \overline{\nu }`$ that we study in this paper have the same advantage of allowing a theoretical clean study of the $`sd\nu \overline{\nu }`$ transition, although some of the relevant semileptonic matrix elements are not as accurately measured experimentally. As shown previously in the Standard Model, they unfortunately have the serious disadvantage that their predicted branching ratios are several orders of magnitude smaller than for the decays involving a single pion. It is nevertheless interesting to pursue them because they involve different combinations of the CP-conserving and CP-violating parts of the $`sd\nu \overline{\nu }`$ transition, and they thus provide additional handles on both the real and imaginary parts of amplitude. Furthermore, from an experimental point of view, some or all of these decays may prove to be more susceptible to the extraction of a signal from the background.
In this paper, we significantly refine the predictions for $`K\pi \pi \nu \overline{\nu }`$ decays in the Standard Model using recent knowledge of the CKM matrix elements and measurements of decay rates for related processes. We also examine how large these branching ratios could be for physics beyond the Standard Model and find that there are significant, model-independent limits from other measurements. The paper is organized as follows: In Section 2, we provide the general framework for studying the $`sd\nu \overline{\nu }`$ transition and the contributions to it both within and beyond the Standard Model. Section 3 sets forth its relationship to the $`K\pi \pi \nu \overline{\nu }`$ decays under discussion, the CP properties of the amplitudes involved in those decays, and the connection of the relevant hadronic matrix elements to measured semileptonic decays. Numerical results are given in Section 4, followed by some conclusions in Section 5.
## II General Framework
The effective Hamiltonian for $`sd\nu \overline{\nu }`$ transitions takes the form
$$=\frac{G_F}{\sqrt{2}}\frac{\alpha }{2\pi \mathrm{sin}^2\theta _W}W_{ds}\left[\left(\overline{s}\gamma _\mu (1\gamma _5)d\right)\left(\overline{\nu }\gamma ^\mu (1\gamma _5)\nu \right)\right]+h.c.,$$
(1)
where the short-distance physics is lumped in $`W_{ds}`$. In the Standard Model, one-loop contributions to $`W_{ds}`$ are dominated by penguin and box diagrams with intermediate charm and top quarks:
$$W_{ds}^{SM}=\lambda _{sd}^cX(x_c)+\lambda _{sd}^tX(x_t),$$
(2)
where $`\lambda _{sd}^iV_{is}^{}V_{id}`$, with $`V_{ij}`$ the appropriate CKM matrix element, and $`x_i=\overline{m}_{i}^{}{}_{}{}^{2}/M_{W}^{}{}_{}{}^{2}`$. The QCD corrections to the short-distance contributions $`X(x_i)`$ have been calculated some time ago in leading order and then in next-to-leading order. Since the top-quark mass is comparable to the weak scale, these corrections are very small for $`X(x_t)`$, as can be seen explicitly in the values given for $`X(x_t)`$ when written as $`X(x_t)=\eta _tX_0(x_t)`$, with the QCD-uncorrected top quark contribution
$$X_0(x_t)=\frac{x_t}{8}\left[\frac{x_t+2}{x_t1}+\frac{3x_t6}{(x_t1)^2}\mathrm{log}(x_t)\right],$$
(3)
and the QCD correction factor $`\eta _t=0.994`$. On the other hand, these corrections have considerable importance for $`X(x_c)`$.
The quantity $`X(x_t)`$ is roughly three orders of magnitude larger than $`X(x_c)`$, and since $`\mathrm{I}m\lambda _{sd}^c=\mathrm{I}m\lambda _{sd}^t`$, the top contribution completely dominates in the imaginary part of $`W_{ds}^{SM}`$. However, Re$`\lambda _{sd}^t<<`$ Re$`\lambda _{sd}^c`$, allowing the charm contribution, although still smaller in magnitude than that from top, to be roughly comparable and to interfere constructively in the real part of $`W_{ds}^{SM}`$.
As illustrative examples of physics that lies beyond the Standard Model, we consider two very different possibilities:
* Effective Flavor-Changing Neutral Current (FCNC) interaction. Such an interaction, as formulated by Nir and Silverman , takes the form of an extra term in the effective Lagrangian of the form:
$$^{(Z)}=\frac{g}{4\mathrm{cos}\theta _W}U_{ds}\overline{d}\gamma _\mu (1\gamma _5)sZ^\mu .$$
(4)
When combined with the coupling of the Z boson to neutrino-antineutrino pairs, one finds that
$$W_{ds}^{NP}=\frac{\pi ^2}{\sqrt{2}G_FM_W^2}U_{ds}=0.93\times 10^2U_{ds}$$
(5)
as the new piece of $`W_{ds}`$ in the effective Hamiltonian that corresponds to the basic process, $`sd\nu \overline{\nu }`$.
Upper bounds for $`U_{ds}`$ have been determined by other processes involving K mesons and were summarized recently to be
$`|\mathrm{R}e(U_{ds})|`$ $``$ $`10^5,`$ (6)
$`|U_{ds}|`$ $``$ $`3\times 10^5,`$ (7)
$`|\mathrm{R}e(U_{ds})\mathrm{I}m(U_{ds})|`$ $``$ $`1.3\times 10^9,`$ (8)
$`|\mathrm{I}m(U_{ds})|`$ $``$ $`10^5.`$ (9)
The bound on $`|U_{ds}|`$ arises from the decay $`K^+\pi ^+\nu \overline{\nu }`$, whose width is proportional to $`|W_{ds}|^2`$. It can be improved by using the most recent measurement, of the branching ratio, BR($`K^+\pi ^+\nu \overline{\nu })=1.5^{+3.4}_{1.2}\times 10^{10}`$. This value is consistent with what is expected in the Standard Model and corresponds to $`|W_{ds}|=0.98{}_{0.54}{}^{+0.80}\times 10^3`$. If we were to assume that the total branching ratio were due to new physics arising from $`U_{ds}`$, then the bound on $`|U_{ds}|`$ would be reduced from that in Eq. (7) to $`|U_{ds}|<1.6\times 10^5`$.
* Supersymmetry. A dominant supersymmetric effect arises from penguin diagrams involving charged-Higgs plus top-quark intermediate states or squark and chargino intermediate states. These give additional pieces to the effective Hamiltonian of the form
$$W_{ds}^{NP}=\lambda _{sd}^t\frac{m_H^2}{M_W^2\mathrm{tan}^2\beta }H(x_{tH})+\frac{1}{96}\stackrel{~}{\lambda }_t,$$
(10)
where $`\mathrm{tan}\beta `$ is the ratio of the two Higgs vacuum expectation values and $`x_{tH}=m_{t}^{}{}_{}{}^{2}/M_{H^\pm }^{}{}_{}{}^{2}`$. The quantity $`H(x)`$ is given by
$$H(x)=\frac{x^2}{8}\left[\frac{\mathrm{log}x}{(x1)^2}+\frac{1}{x1}\right],$$
(11)
The parameter $`\stackrel{~}{\lambda }_t`$ can be bounded by similar considerations to those that were used for $`U_{ds}`$. The observed branching ratios for the decays $`K_L\mu ^+\mu ^{}`$ and $`K^+\pi ^+\nu \overline{\nu }`$ have been used to set the limits
$`|\mathrm{R}e\stackrel{~}{\lambda }_t|`$ $``$ $`0.21,`$ (12)
$`|\stackrel{~}{\lambda }_t|`$ $``$ $`0.35.`$ (13)
The most recent branching ratio for $`K^+\pi ^+\nu \overline{\nu }`$ could be used to revise the last limit to $`|\stackrel{~}{\lambda }_t|<0.16`$ .
As we will see shortly, the limitations imposed by experiments on the parameters of both these examples of physics beyond the Standard Model lead to similar restrictions on how large the branching ratios can be for the processes we are studying.
## III $`𝑲\mathbf{}𝝅𝝅𝝂\overline{𝝂}`$ Decays
When the effective four-fermion operator relevant for the decay we are considering is sandwiched between the initial and final states, it factorizes into a product of matrix elements of the hadronic current and the leptonic current. We will use isotopic spin to relate the hadronic matrix elements relevant to $`K_{L,S}\pi \pi \nu \overline{\nu }`$ to those for $`K^+\pi \pi e^+\nu `$, where the corresponding branching ratios (and hence squares of matrix elements) have been measured.
We consider first the process $`K_L\pi ^o\pi ^o\nu \overline{\nu }`$. The $`\pi ^o\pi ^o`$ pair forms a CP-even state, and has total isospin, $`I=0`$. The $`\nu \overline{\nu }`$ pair, created by a virtual $`Z^o`$, is CP even as well. Since there must be one unit of orbital angular momentum to allow the total angular momentum of the final state to be that of the initial $`K_L`$, namely zero, the final state is CP-odd. The overall decay process is then CP-conserving for the major (CP-odd) piece of the $`K_L`$, and the resulting amplitude is proportional to the real part of $`W_{ds}`$. The opposite result holds for the piece of the $`K_L`$ that is proportional to $`ϵ`$ and CP-even; the corresponding decay process is CP-violating and the amplitude is proportional to the imaginary part of $`W_{ds}`$. It is thus suppressed on two counts and is completely negligible.
Using the relationship,
$$\pi ^o\pi ^o(\overline{s}d)_{\mathrm{V}\mathrm{A}}K^o=\pi ^o\pi ^o(\overline{s}u)_{\mathrm{V}\mathrm{A}}K^+,$$
(14)
we find that
$$BR(K_L\pi ^o\pi ^o\nu \overline{\nu })=\frac{3\alpha ^2|\mathrm{Re}W_{ds}|^2}{2\pi ^2\mathrm{sin}^4\theta _W|V_{us}|^2}\frac{\tau _{K_L}}{\tau _{K^+}}BR(K^+\pi ^o\pi ^oe^+\nu ),$$
(15)
where the factor of 3 accounts for the three species of neutrinos.
By relating the desired branching ratio to a measured one, we have avoided having either to do a calculation of the hadronic matrix elements or to perform a detailed analysis in terms of invariant amplitudes, as was done in previous analyses. Of course, the final results for the branching ratio must be consistent, since both approaches agree with the available data on charged-current semileptonic decays, and in particular those for the decay rate for the process $`K^+\pi ^o\pi ^oe^+\nu `$. For the purposes of this paper of discussing the absolute and relative size of the various branching ratios within and beyond the Standard Model, it is considerably easier to formulate the results directly in terms of relationships to branching ratios for measured semileptonic decays.
A similar formula can be obtained for $`BR(K_S\pi ^0\pi ^0\nu \overline{\nu })`$, but with $`\mathrm{R}eW_{ds}`$ replaced by $`\mathrm{I}mW_{ds}`$ and $`\tau _{K_L}`$ replaced by $`\tau _{K_S}`$. Since the $`K_S`$ has a much shorter lifetime and the major part of the $`K_S`$ corresponds to a transition that is CP-violating, this branching ratio is orders of magnitude smaller than that for $`K_L\pi ^o\pi ^o\nu \overline{\nu }`$. Although the part of the $`K_S`$ state proportional to $`ϵ`$ corresponds a CP-conserving transition, the smallness of $`ϵ`$ still gives rise to a net decay amplitude that is much smaller than that for the CP-even part of the $`K_S`$.
The analysis for the decay $`K_L\pi ^+\pi ^{}\nu \overline{\nu }`$ can be carried out analogously. It is convenient to break it up into the cases where the $`\pi ^+\pi ^{}`$ pair in the final state has total isospin zero and one, since there is no interference between them in the decay rate. For the isospin zero case, the argument about the CP properties of the final state is the same as given before, and we simply have a factor of two in the rate for the $`\pi ^+\pi ^{}`$ final state compared to that for the $`\pi ^o\pi ^o`$ final state discussed above:
$$BR(K_L(\pi ^+\pi ^{})_{\mathrm{I}=0}\nu \overline{\nu })=\frac{3\alpha ^2|\mathrm{Re}W_{ds}|^2}{\pi ^2\mathrm{sin}^4\theta _W|V_{us}|^2}\frac{\tau _{K_L}}{\tau _{K^+}}BR(K^+\pi ^o\pi ^oe^+\nu ).$$
(16)
This situation is slightly more complicated for the case where the $`\pi ^+\pi ^{}`$ pair has isospin, $`I=1`$. The $`\pi \pi `$ pair is still CP-even, but it must be in a p-wave. There are two possible ways in which the total angular momentum of the final state can be zero, which correspond to the relative orbital angular momentum of the $`\pi \pi `$ and $`\nu \overline{\nu }`$ pairs being zero or one. As is expected when there is such limited phase space, the latter amplitude is strongly suppressed by centrifugal barrier effects compared to the former. So we are left with a single amplitude where the relative orbital angular momentum is zero. The CP of the final state is even and the transition involving the major part of the $`K_L`$ is CP-violating. Using the relationship,
$$\sqrt{2}(\pi ^+\pi ^{})_{\mathrm{I}=1}(\overline{s}d)_{\mathrm{V}\mathrm{A}}K^o=(\pi ^{}\pi ^o)_{\mathrm{I}=1}(\overline{s}u)_{\mathrm{V}\mathrm{A}}K^o,$$
(17)
we find that
$$BR(K_L(\pi ^+\pi ^{})_{\mathrm{I}=1}\nu \overline{\nu })=\frac{3\alpha ^2|\mathrm{Im}W_{ds}|^2}{4\pi ^2\mathrm{sin}^4\theta _W|V_{us}|^2}BR(K_L\pi ^{}\pi ^oe^\pm \nu ).$$
(18)
The total branching ratio for $`K_L\pi ^+\pi ^{}\nu \overline{\nu }`$ is then found by simply adding the two results above: $`BR(K_L\pi ^+\pi ^{}\nu \overline{\nu })=BR(K_L(\pi ^+\pi ^{})_{\mathrm{I}=0}\nu \overline{\nu })+BR(K_L(\pi ^+\pi ^{})_{\mathrm{I}=1}\nu \overline{\nu })`$.
The corresponding formula for $`K_S\pi ^+\pi ^{}\nu \overline{\nu }`$ can again be obtained by the interchange of $`\mathrm{R}eW_{sd}`$ and $`\mathrm{I}mW_{sd}`$ and multiplication of the right-hand-side by $`\tau _{K_S}/\tau _{K_L}`$.
## IV Numerical Calculation
To obtain numerical predictions we have used a set of parameters taken from the Review of Particle Physics, including the fine-structure constant at the weak scale, $`\alpha =1/129`$; $`M_W=80.3`$ GeV; $`\mathrm{sin}^2\theta _W=0.23`$; and the measured semileptonic branching ratios needed in Eqs. (15), (16) and (18). We correspondingly find that:
$`BR(K_L\pi ^o\pi ^o\nu \overline{\nu })`$ $`=`$ $`[(3.1\pm 0.6)\times 10^7]|\mathrm{Re}W_{ds}|^2,`$ (19)
$`BR(K_L(\pi ^+\pi ^{})_{\mathrm{I}=0}\nu \overline{\nu })`$ $`=`$ $`[(6.2\pm 1.2)\times 10^7]|\mathrm{Re}W_{ds}|^2,`$ (20)
$`BR(K_L(\pi ^+\pi ^{})_{\mathrm{I}=1}\nu \overline{\nu })`$ $`=`$ $`[(0.93\pm 0.05)\times 10^7]|\mathrm{Im}W_{ds}|^2,`$ (21)
where the error bars come from those of the experimental measurements of the relevant semileptonic branching ratios. We have not taken account of radiative corrections or isotopic-spin violating differences in form factors and phase space, as has been done for the case of the decays involving a single pion, given the (larger) uncertainties in other parts of input at this stage of the analysis of these decays. Formulas similar to Eq. (19) hold for the related decays of the $`K_S`$ to the same final states.
For the specific calculation of $`W_{ds}`$ in the Standard Model we need the values of $`X(x_t)`$ and of $`X(x_c)`$ in next-to-leading order, and that of the CKM matrix elements
$`\mathrm{R}eV_{td}`$ $`=`$ $`0.0076\pm 0.0015,`$ (22)
$`\mathrm{I}mV_{td}`$ $`=`$ $`0.0031\pm 0.0008,`$ (23)
and $`V_{ts}=V_{cb}=0.040\pm 0.002`$, aside from the well-known matrix elements connecting the first and second generations. Using this and with $`\overline{m_t}=166\pm 5`$ GeV, we find that
$$W_{ds}^{SM}=[(6.7\pm 1.0)+i(1.9\pm 0.5)]\times 10^4,$$
(24)
and the branching ratios for the various processes shown in Table I.
In $`K_L`$ decays the contribution of the isospin one $`\pi ^+\pi ^{}`$ final state is negligible, since it is already suppressed compared to that with isospin zero from Eq. (19) and the magnitude of the real part of $`W_{ds}^{SM}`$ is considerably greater than that of the imaginary part. Thus the ratio between $`\pi ^+\pi ^{}`$ and $`\pi ^o\pi ^o`$ rates is very close to the factor of two characteristic of isospin zero.
Our results are given in Table I, and both $`K_L`$ branching ratios lie between $`10^{13}`$ and $`10^{12}`$. Our predictions in the Standard Model for $`BR(K_L\pi ^o\pi ^o\nu \overline{\nu })`$ are consistent with the previous calculation of $`13\times 10^{13}`$ and those of $`1.15\times 10^{13}`$ and $`25\times 10^{13}`$ for $`BR(K_L\pi ^+\pi ^{}\nu \overline{\nu })`$, but the allowed range is now considerably restricted. The $`K_S`$ branching ratios are in the $`10^{17}`$ range in the Standard Model. The decay $`K_S\pi ^+\pi ^{}\nu \overline{\nu }`$ gets important contributions from both $`I=0`$ and $`I=1`$ $`\pi \pi `$ final states since the suppression of the $`I=1`$ final state is compensated by the ratio of $`(\mathrm{R}eW_{ds}/\mathrm{I}mW_{ds})^2`$ in the Standard Model.
For the representative examples of physics beyond the Standard Model, we also show in Table I values for the branching ratios that correspond to the maximal values one could obtain consistent with the bounds in Eqs. (6)- (9) and (12) - (13), respectively. These maximal values are similar in both cases and arise when the parameters of the new physics are chosen to maximize $`W_{ds}^{NP}`$ consistent with the constraints coming from known K physics. Among these constraints is the recent branching ratio measurement for $`K^+\pi ^+\nu \overline{\nu }`$, which is equally sensitive to both the real and imaginary parts of the total $`W_{ds}`$ with minimal theoretical assumptions. Hence, similar maximal values of $`W_{ds}^{NP}`$ are obtained in any model of new physics. Note that when taking ratios to the Standard Model branching ratios, a much bigger factor is possible when the new physics enters the imaginary part of $`W_{sd}^{NP}`$, and is CP-violating, since the imaginary part of $`W_{sd}^{SM}`$ is considerably smaller than the real part.
## V Summary
We have used recent information on the CKM matrix to narrow the range of the predicted branching ratios for $`K_{L,S}\pi ^o\pi ^o\nu \overline{\nu }`$ and $`K_{L,S}\pi ^+\pi ^{}\nu \overline{\nu }`$ decays in the Standard Model. These branching ratios are in the neighborhood of $`14\times 10^{13}`$ for the $`K_L`$ decays, which make them possibly observable at the few event level in the next round of experiments that are setting out to see the CP-violating decays with a single $`\pi ^o`$ in the final state. These branching ratios could be larger by up to about an order of magnitude in theories that go beyond the Standard Model.
The branching ratios for the $`K_S`$ decays are in the neighborhood of $`10^{17}`$ to $`10^{16}`$ in the Standard Model and seem unlikely to ever be observed. Here new physics could boost the branching ratios by more than an order of magnitude, although even then the maximum branching ratio of around $`10^{15}`$ is still beyond the limits of observation. For both $`K_S`$ and $`K_L`$ decays, increased experimental accuracy in the measurement of the branching ratio for $`K^+\pi ^+\nu \overline{\nu }`$, assuming it remains consistent with the Standard Model, will put more stringent restrictions on non-Standard-Model physics in the $`sd\nu \overline{\nu }`$ transition and limit the deviations from the Standard Model that can be observed in the decays under discussion here as well.
ACKNOWLEDGMENT
This research work is supported by the Department of Energy under Grant No. DE-FG02-91ER40682. Fred Gilman thanks Y. Wah for discussions that started this investigation.
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# Direct imaging search for planetary companions next to young nearby stars
## 1. Introduction
Despite extensive imaging surveys, only a few sub-stellar companions to normal stars were detected already by direct imaging: Gl 229 B (Nakajima et al. 1995, Oppenheimer et al. 1995), G196-3 B (Rebolo et al. 1998), and Gl 570 D (Burgasser et al. 2000), brown dwarf companions confirmed by spectroscopy and proper motion. Three more candidates were presented: GG Tau Bb (White et al. 1999), HR 7329 B (Lowrance et al. 2000), and CoD$`33^{}7795`$ B (Lowrance et al. 1999), but spectroscopy and/or proper motion were not available so far. Most recently, Neuhäuser et al. (2000b) have shown that CoD$`33^{}7795`$ B has spectral type M8.5 to M9 with an optical and an infrared spectrum, both taken with the VLT, and that it is co-moving with star A, after two years epoch difference with $`5\sigma `$ significance. The extra-solar planet candidate directly detected by the HST near a T Tauri star in Taurus (Terebey et al. 1998) has not yet been confirmed by spectroscopy (Terebey et al. 2000).
Several extra-solar planet candidates have been detected indirectly by radial velocity variations of stars (Latham et al. 1989, Mayor & Queloz 1995, etc, see review by Marcy & Butler 1998), one such candidate is confirmed by a transit event (Charbonneau et al. 2000). There is no direct imaging detection of an extra-solar planet, yet. Direct imaging detection of planets like those in our solar system but orbiting other stars is difficult due to the limited dynamical range: Extra-solar planets are simply too faint and too close to their bright host stars. One can try to avoid the problem of spatial resolution by searching for planetary companions around nearby stars, where the orbit of the outermost solar system planet corresponds to several arcsec, sufficient to resolve a faint object next to a bright star. However, very nearby stars usually are too old, so that their hypothetical planets (like e.g. our Jupiter) are correspondly too faint for direct detection with current technology.
Young planets, on the other hand, are still self-luminous due to on-going accretion and/or contraction (Burrows et al. 1997, Brandner et al. 1997, Malkov, Piskunov & Zinnecker 1998) and, if also nearby, they would be sufficiently bright and resolved for direct detection. Direct imaging of these young planets is optimal at the near-infrared bands H and K, where the brightness difference between young stars and young planets is expected to be the lowest (Burrows et al. 1997) and where also the seeing is better than in the optical. In addition, nearby stars usually have large proper motion, so that one can decide after only a few years whether a companion candidate is co-moving. Finally, there is a crucial advantage in studying companion planets candidates instead of free-floating planet candidates: The mass can be better constrained for companions than for free-floating planets, because age and distance of the primary is usually well-known.
## 2. Our sample: Young nearby stars
We selected young ($`100`$ Myr) nearby ($`75`$ pc) stars, some of them from the literature, others discovered recently among ROSAT sources by ourselves (some of these as yet unpublished). We are confident on the young age of our sample of stars because of Lithium 6708Å absorption lines (and/or H$`\alpha `$ emission, IR excess and/or kinematic membership to a young cluster). We know the distances of most target stars from Hipparcos (and for some of them, from kinematic membership to a cluster with known distance). Our target list includes all members of the TW Hya association (TWA), the T Tauri stars in the nearby and even younger MBM 12 cloud (Hearty et al. 2000a,b), the members of the more recently discovered Tucanae (Zuckerman & Webb 2000) and HorA moving groups (Torres et al. 2000), as well as isolated young stars like GJ 182.
Because we know the distances and ages of the stars in our sample, we can predict the H- and K-band magnitudes of possible substellar companions for different masses, e.g. from 1 to $`80M_{jup}`$, using the non-gray theory by Burrows et al. (1997). In Figure 1, we present the H-band magnitudes of the program stars (either known from infrared photometry, or estimated from their known spectral types and optical magnitudes), their distance distribution, the corresponding angular separation between those stars and possible companions at assumed 50 AU physical separations, and the K-band magnitude distribution of assumed $`10M_{jup}`$ mass companions (estimated from Burrows et al. (1997) for the known ages and distances of the stars). These distributions clearly show that such companions can be detectable and resolvable with current technology.
Hence, we have started an observational program using mainly the SOFI and SHARP infrared cameras at the ESO-3.5m-NTT on La Silla, and the ISAAC imaging camera at the ESO-8.2m-Antu (VLT-UT1), but also other state of the art instruments.
## 3. First results: Case study TWA-7
Using the MPE-build SHARP speckle camera, we have detected a very faint object 2.5 arcsec south-east of the young nearby star T Tauri TWA-7, a member of the TW Hya association (called TWA, see Webb et al. 1999). Four stars of the TWA association have been observed by Hipparcos, their mean distance being 55 pc. By comparison with isochrones, we obtain an age of 1 to 6 Myr (Neuhäuser et al. 2000a).
The faint object TWA-7B detected by SHARP, 2.5 arc sec south-east of TWA-7A has $`H=16.4`$ and $`K=16.3`$ mag, which is more than 9 mag fainter than the primary star at these wavelengths. If TWA-7B were to be a companion to TWA-7A, i.e. if it were to be at the same distance and age, then its apparent H- and K-band magnitudes would correspond to absolute magnitudes (at 55 pc) consistent with an effective temperature T$`{}_{eff}{}^{}1050`$ K and a surface gravity $`3000`$ g/s<sup>2</sup> (see table 5 in Burrows et al. 1997). These values are then consistent with an object with a mass of $`3`$M<sub>jup</sub> and an age of $`10^{6.5}`$ yr (see Figure 9 in Burrows et al. 1997). This derived age is in agreement with the ages of TWA-7A and the other TWA members. The angular separation of 2.5 arc sec (at a distance of 55 pc) corresponds to a physical separation of 138 AU, well within typical T Tauri disk sizes.
TWA-7B has previously been detected by HST Nicmos observations in the F160W, F090M, F165M, and F180M filters. The HST F160W image has been shown in Neuhäuser et al. (2000a), where a coronograph has been used as in the F165M and F180M images. Only in the F090M image with NIC1, no coronograph was used. See Neuhäuser et al. (2000a) for the HST magnitudes, the position angles between TWA-7A and B, and more details on data reduction.
Confirmation of the possibly substellar nature of TWA-7B by spectroscopy is required. To check for a possible spectral signature of a very cool object we took an H-band spectrum of TWA-7B using ISAAC at the ESO-8.2m-Antu (VLT-UT1). In fact, the derived temperature of $`1050`$ K for TWA-7B (if indeed a companion) is similar to those of known old T-dwarfs, so that one might, under some conditions, expect to see methane absorption features in the spectrum of TWA-7B. The technical problem with taking and analysing a spectrum of such a faint object so very close ($``$ 2.5 arcsec) to a much brighter star (contrast $`10^4`$) is again dynamical range. To maximize both the separation between TWA-7B and A and also the fraction of the light coming from TWA-7B, we placed the slit neither along both stars (i.e. along their Position Angle (PA)), nor perpendicular to this PA (with only B in the slit). In the latter case, the signal from B would be on top of the dominant signal from TWA-7A scattered light in the collapsed spectrum (zero separation). In the former case (maximum separation), the light from TWA-7A would have “swamped” the signal from TWA-7B. We used, instead, a slit orientation in between those two extremes.
We modelled the flux of TWA-7A at each wavelength and then subtracted it from TWA-7B. The spectrum of TWA-7B is shown in Figure 2. TWA-7B turned out to be of spectral type K. This means that, given its brightness, TWA-7B cannot be at the same small distance as TWA-7A (55 pc) and it is most likely a background K-type main sequence star. Given that it is a faint source ($`H=16.4`$ mag), this K dwarf should be at a distance between $`2`$ and 4 kpc, in the halo of our galaxy ($`b=21^{}`$). Curiously, the probability of finding a background object as faint as TWA-7B within a 2.5 arsec radius circle in the sky, towards this direction of the Galaxy, is about 1%. This shows how important it is to take spectra of companion candidates.
## 4. Summary
We show that ground-based direct imaging detection of extra-solar planets is in principle possible. As an example, we present the detection of a very faint object ($`H=16.4`$, $`K=16.3`$ mag, i.e. $`9.5`$ mag fainter than a nearby young star), which could have been a few Jupiter-mass giant extra-solar planet, given its brightness and color, and given the fact that it is located very close to a young star. At the age and distance for this star, this faint object could have been the first direct imaging detection of an extra-solar planet. However, our H-band spectrum shows that it is likely to be a background K-dwarf. We demonstrate, nevertheless, that spectroscopic observation of such a faint object that close to a brighter star (contrast $`10^4`$) is possible with current technology from ground-based telescopes.
We are now living in very special times, where we are starting to be able to directly image extra-solar planets and obtain their spectra. We believe that the first direct imaging detection of an extra-solar planet is imminent. For achiving this goal, a large-scale survey is essential.
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# On automorphisms of arithmetic subgroups of unipotent groups in positive characteristic
## 1. Introduction
Roughly speaking, a discrete subgroup $`\mathrm{\Gamma }`$ of a topological group $`G`$ is automorphism rigid if every automorphism of $`\mathrm{\Gamma }`$ extends to a continuous automorphism of $`G`$. However, the formal definition below is slightly more complicated, because it allows for passage to finite-index subgroups.
###### 1.1 Definition.
It is traditional to say that a group $`\mathrm{\Gamma }`$ *virtually* has a property if some finite-index subgroup of $`\mathrm{\Gamma }`$ has the property. It is convenient to extend this terminology to group isomorphisms.
* A *virtual isomorphism* from $`G_1`$ to $`G_2`$ is an isomorphism $`\mathrm{\Lambda }:G_1^{}G_2^{}`$, where $`G_i^{}`$ is a finite-index, open subgroup of $`G_i`$.
* A *virtual automorphism* of $`G`$ is a virtual isomorphism from $`G`$ to $`G`$.
* A virtual isomorphism $`\mathrm{\Lambda }`$ from $`G_1`$ to $`G_2`$ *virtually extends* an isomorphism $`\lambda `$ from $`\mathrm{\Gamma }_1`$ to $`\mathrm{\Gamma }_2`$ if there is a finite-index, open subgroup $`\mathrm{\Gamma }_1^{}`$ of $`\mathrm{\Gamma }_1`$, such that $`\mathrm{\Gamma }_1^{}G_1`$, and $`\mathrm{\Lambda }|_{\mathrm{\Gamma }_1^{}}=\lambda |_{\mathrm{\Gamma }_1^{}}`$.
###### 1.2 Definition.
A discrete subgroup $`\mathrm{\Gamma }`$ of a topological group $`G`$ is *automorphism rigid* in $`G`$ if every virtual automorphism of $`\mathrm{\Gamma }`$ virtually extends to a virtual automorphism of $`G`$.
A classical example is provided by the work of Malcev.
###### 1.3 Definition (\[Rag, Rem. 1.11, p. 21\]).
A discrete subgroup $`\mathrm{\Gamma }`$ of a topological group $`G`$ is a (cocompact) *lattice* if $`G/\mathrm{\Gamma }`$ is compact.
###### 1.4 Theorem (Malcev \[Mal\], \[Rag, Cor. 2.11.1, p. 34\]).
If $`\mathrm{\Gamma }`$ is a lattice in a $`1`$-connected, nilpotent real Lie group $`G`$, then $`\mathrm{\Gamma }`$ is automorphism rigid in $`G`$.
In fact, every virtual automorphism of $`\mathrm{\Gamma }`$ extends to a unique automorphism of $`G`$.
Malcev’s Theorem can be restated in the terminology of algebraic groups (cf. \[Rag, after Thm. 2.12, p. 34\]). Recall that a matrix group $`G`$ is *unipotent* if, for every $`gG`$, there is some $`n`$, such that $`(g\mathrm{Id})^n=0`$. (In other words, $`1`$ is the only eigenvalue of $`g`$.)
###### 1.5 Corollary.
Let $`\mathrm{\Gamma }`$ be an arithmetic subgroup of a unipotent algebraic $``$-group $`𝔾`$. Then $`\mathrm{\Gamma }`$ is an automorphism rigid lattice in $`𝔾()`$.
In this paper, we discuss the analogue of Malcev’s Theorem for unipotent groups over nonarchimedean local fields, instead of $``$. It is well known that if $`𝔾`$ is a unipotent algebraic group over a nonarchimedean local field $`L`$ of characteristic zero, then the group $`𝔾(L)`$ of $`L`$-points of $`𝔾`$ has no nontrivial discrete subgroups. (For example, $``$ is not discrete in the $`p`$-adic field $`_p`$.) Thus the case of characteristic zero is not of interest in this setting; we will consider only local fields of positive characteristic.
For abelian groups, it is easy to prove automorphism rigidity.
###### 1.6 Proposition.
Let $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ be lattices in a totally disconnected, locally compact, abelian group $`G`$. Then every isomorphism $`\lambda :\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$ virtually extends to a virtual automorphism $`\widehat{\lambda }`$ of $`G`$.
###### Proof.
Since $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ are discrete, and $`G`$ is totally disconnected, there exists a compact, open subgroup $`K`$ of $`G`$, such that $`\mathrm{\Gamma }_1K=\mathrm{\Gamma }_2K=e`$. Let $`\widehat{G_1}=\mathrm{\Gamma }_1K`$ and $`\widehat{G_2}=\mathrm{\Gamma }_2K`$, so $`\widehat{G_1}`$ and $`\widehat{G_2}`$ are finite-index, open subgroups of $`G`$, and define $`\widehat{\lambda }:\widehat{G_1}\widehat{G_2}`$ by $`\widehat{\lambda }(\gamma c)=\lambda (\gamma )c`$ for $`\gamma \mathrm{\Gamma }_1`$ and $`cK`$. ∎
For nonabelian groups, automorphism rigidity seems to be surprisingly more difficult to prove, but we provide examples of automorphism rigid lattices. Although we do not have a general theory, and we do not have enough evidence to support a specific conjecture, the examples suggest that there may be mild conditions that imply arithmetic lattices are automorphism rigid.
###### 1.7 Notation.
* Fix a prime $`p`$, and a power $`q`$ of $`p`$.
* $`𝔽_q`$ denotes the finite field of $`q`$ elements.
* $`F`$ denotes the field $`𝔽_q((t))`$ of formal power series over $`𝔽_q`$.
* $`F^{}`$ denotes $`𝔽_q[t^1]`$, the $`𝔽_q`$-subalgebra of $`F`$ generated by $`t^1`$.
Note that $`F`$ is a local field of characteristic $`p`$. (Conversely, any local field of characteristic $`p`$ is isomorphic to $`𝔽_q((t))`$, for some $`q`$ \[Wei, Thm. I.4.8, p. 20\].) The subgroup $`F^{}`$ is a lattice in the additive group $`(F,+)`$.
###### 1.8 Definition.
Let $`G`$ be a closed subgroup of $`\mathrm{GL}(m,F)`$, for some $`m`$.
* Two discrete subgroups $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ of $`G`$ are *commensurable* if $`\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$ is a finite-index subgroup of both $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ \[Mar, p. 8\].
* A subgroup $`\mathrm{\Gamma }`$ of $`G`$ is *arithmetic* if it is commensurable with $`\mathrm{GL}(m,F^{})G`$ (cf. \[Mar, §I.3.1, pp. 60–62\]).
By definition, if $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ are arithmetic subgroups of $`G`$, then $`\mathrm{\Gamma }_1`$ is commensurable with $`\mathrm{\Gamma }_2`$. Thus, $`\mathrm{\Gamma }_1`$ is a lattice in $`G`$ if and only if $`\mathrm{\Gamma }_2`$ is a lattice in $`G`$.
###### 1.9 Definition (cf. \[BS, Ex. 9.2\]).
Fix a power $`r`$ of $`p`$, and let
$$G_2=\left\{\left(\begin{array}{ccc}1& y^r& z\\ 0& 1& y\\ 0& 0& 1\end{array}\right)|y,zF\right\}.$$
So $`G_2`$ is a two-dimensional, unipotent $`F`$-group, and has arithmetic lattices. Note that if $`r>1`$, then $`G_2`$ is nonabelian.
The following theorem describes the virtual automorphisms of any arithmetic lattice in $`G_2`$.
###### 1.10 Definition.
For any continuous field automorphism $`\tau `$ of $`F`$ and any $`aF\{0\}`$, there is a continuous automorphism $`\varphi _{\tau ,a}`$ of $`G_2`$, defined by
$$\varphi _{\tau ,a}\left(\begin{array}{ccc}1& y^r& z\\ 0& 1& y\\ 0& 0& 1\end{array}\right)=\left(\begin{array}{ccc}1& a^r\tau (y)^r& a^{r+1}\tau (z)\\ 0& 1& a\tau (y)\\ 0& 0& 1\end{array}\right).$$
Let us say that $`\varphi _{\tau ,a}`$ is *standard* if
1. there exist $`\sigma \mathrm{Gal}(𝔽_q/𝔽_p)`$, $`\alpha 𝔽_q\{0\}`$, and $`\beta 𝔽_q`$, such that
$$\tau \left(f(t^1)\right)=\sigma \left(f(\alpha t^1+\beta )\right),$$
for all $`f(t^1)F`$, and
2. there exists some nonzero $`bF^{}`$, such that $`abF^{}`$.
Note that if $`\varphi _{\tau ,a}`$ is standard, and $`\mathrm{\Gamma }`$ is an arithmetic lattice in $`G_2`$, then $`\varphi _{\tau ,a}(\mathrm{\Gamma })`$ is commensurable with $`\mathrm{\Gamma }`$.
###### 1.11 Theorem.
Let
* $`\mathrm{\Gamma }`$ be an arithmetic lattice in $`G_2`$; and
* $`\lambda `$ be a virtual automorphism of $`\mathrm{\Gamma }`$.
If $`r>2`$, then there exist
* a standard automorphism $`\varphi _{\tau ,a}`$ of $`G_2`$,
* a finite-index subgroup $`\mathrm{\Gamma }^{}`$ of $`\mathrm{\Gamma }`$, and
* a homomorphism $`\zeta :\mathrm{\Gamma }^{}Z(\mathrm{\Gamma })`$,
such that $`\lambda (\gamma )=\varphi _{\tau ,a}(\gamma )\zeta (\gamma )`$, for all $`\gamma \mathrm{\Gamma }^{}`$.
###### 1.12 Corollary.
If $`r2`$, then any arithmetic lattice in $`G_2`$ is automorphism rigid.
Theorem 1.11 and Corollary 1.12 are proved in Section 2. The authors do not know whether they remain true in the exceptional case $`r=p=2`$.
###### 1.13 Definition.
Assume $`p>2`$, let $`[[,]]:F^{2m}\times F^{2m}F`$ be a symplectic form, and, for notational convenience, let $`Z=F`$. The corresponding *Heisenberg group* is the group $`H=(F^{2m}\times Z,)`$, where
$$(v_1,z_1)(v_2,z_2)=(v_1+v_2,z_1+z_2+[[v_1,v_2]]).$$
We remark that, up to a change of basis, the symplectic form $`[[,]]`$ on $`F^{2m}`$ is unique, so, up to isomorphism, the Heisenberg group $`H`$ is uniquely determined by $`m`$. Note that $`Z`$ is the center of $`H`$.
Because $`H`$ is isomorphic to a subgroup of $`\mathrm{GL}(m+2,F)`$, namely,
$$H\left\{\left(\begin{array}{cccccc}1& x_1& x_2& \mathrm{}& x_m& z\\ & 1& & & & y_1\\ & & 1& \text{0}& & y_2\\ & & & \mathrm{}& & \mathrm{}\\ & \text{0}& & & 1& y_m\\ & & & & & 1\end{array}\right)|\begin{array}{c}x_1,\mathrm{},x_mF,\\ y_1,\mathrm{},y_mF,\\ zF\end{array}\right\},$$
we may speak of arithmetic subgroups of $`H`$.
We assume that $`[[,]]`$ is defined over $`F^{}`$, by which we mean that $`[[F^{},F^{}]]F^{}`$. Then we may assume that the above isomorphism has been chosen so that
a subgroup $`\mathrm{\Gamma }`$ of $`H`$ is arithmetic if and only if it is commensurable with $`(F^{})^{2m}\times F^{}`$.
Thus, $`H`$ has arithmetic lattices.
We remark that one may define Heisenberg groups even if $`p=2`$, but, in this case, they are abelian, so they are not of particular interest.
###### 1.14 Definition.
We say $`TGL(2m,F)`$ is conformally symplectic if there exists some nonzero $`c_TF`$, such that, for all $`v,wV`$, we have
$$[[T(v),T(w)]]=c_T[[v,w]].$$
For every conformally symplectic $`TGL(2m,F)`$, and every continuous field automorphism $`\tau `$ of $`F`$, there is a continuous automorphism $`\varphi _{T,\tau }`$ of $`H`$ defined by
$$\varphi _{T,\tau }(v,z)=(\tau \left(T(v)\right),\tau (c_Tz)).$$
Let us say that $`\varphi _{T,\tau }`$ is standard if
1. there exist $`\sigma \mathrm{Gal}(𝔽_q/𝔽_p)`$, $`\alpha 𝔽_q\{0\}`$, and $`\beta 𝔽_q`$, such that
$$\tau \left(f(t^1)\right)=\sigma \left(f(\alpha t^1+\beta )\right)$$
for all $`f(t^1)F`$; and
2. there exists some nonzero $`bF^{}`$, such that $`bT\mathrm{Mat}(2m,F^{})`$.
Note that if $`\varphi _{T,\tau }`$ is standard, then $`\varphi _{T,\tau }(\mathrm{\Gamma })`$ is commensurable with $`\mathrm{\Gamma }`$ for any arithmetic lattice $`\mathrm{\Gamma }`$ of $`H`$.
###### 1.15 Theorem.
Assume $`p>2`$. Let
* $`\mathrm{\Gamma }`$ be an arithmetic lattice in a Heisenberg group $`H`$; and
* $`\lambda `$ be a virtual automorphism of $`\mathrm{\Gamma }`$.
Then there exist
* a standard automorphism $`\varphi _{T,\tau }`$ of $`H`$;
* a finite index subgroup $`\mathrm{\Gamma }^{}`$ of $`\mathrm{\Gamma }`$; and
* a homomorphism $`\zeta :\mathrm{\Gamma }^{}Z(\mathrm{\Gamma })`$,
such that $`\lambda (\gamma )=\varphi _{T,\tau }(\gamma )\zeta (\gamma )`$, for all $`\gamma \mathrm{\Gamma }^{}`$.
###### 1.16 Corollary.
If $`p>2`$, then any arithmetic lattice in a Heisenberg group $`H`$ is automorphism rigid.
Theorem 1.15 and Corollary 1.16 are proved in Section 3.
###### 1.17 Remark.
Malcev’s Theorem 1.4 does not extend to all lattices in solvable Lie groups. (See the work of A. Starkov \[Sta\] for a thorough discussion.) On the other hand, the Mostow Rigidity Theorem \[Mos\] implies that lattices in most semisimple Lie groups are automorphism rigid.
Superrigidity deals with extending homomorphisms, instead of only isomorphisms. The Margulis Superrigidity Theorem \[Mar, Thm. VII.5.9, p. 230\] implies that lattices in most semisimple Lie groups are superrigid. (Lattices in many non-semisimple Lie groups are also superrigid \[Wit\].) The Superrigidity Theorem also applies to arithmetic subgroups of many semisimple groups defined over nonarchimedean local fields, whether they are of characteristic zero or not \[Mar, Ven\].
###### 1.18 Acknowledgments.
The authors would like to thank the University of Bielefeld (Germany), the Isaac Newton Institute for Mathematical Sciences (Cambridge, U.K.), the University of Virginia, and Oklahoma State University for their hospitality. Most of this research was carried out during productive visits to these institutions. Financial support was provided by the German-Israeli Foundation for Research and Development and the National Science Foundation (DMS-9801136).
## 2. Arithmetic subgroups of the two-dimensional unipotent group $`G_2`$
Recall that $`r`$ and $`G_2`$ are defined in Definition 1.9. (Also recall the definitions of $`p`$, $`q`$, $`F`$, and $`F^{}`$ in Notation 1.7.)
###### Proof of Theorem 1.11.
Let $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ be finite-index subgroups of $`\mathrm{\Gamma }`$, such that $`\lambda `$ is an isomorphism from $`\mathrm{\Gamma }_1`$ to $`\mathrm{\Gamma }_2`$. Then $`\lambda `$ induces isomorphisms
$`\lambda ^{}:\mathrm{\Gamma }_1/Z(\mathrm{\Gamma }_1)\mathrm{\Gamma }_2/Z(\mathrm{\Gamma }_2)`$ and $`\lambda _{}:[\mathrm{\Gamma }_1,\mathrm{\Gamma }_1][\mathrm{\Gamma }_2,\mathrm{\Gamma }_2]`$.
By identifying each of $`G_2/Z(G_2)`$ and $`Z(G_2)`$ with $`F`$ in the natural way (and noting that $`\mathrm{\Gamma }_iZ(G_2)=Z(\mathrm{\Gamma }_i)`$), we may think of $`\mathrm{\Gamma }_i/Z(\mathrm{\Gamma }_i)`$ and $`[\mathrm{\Gamma }_i,\mathrm{\Gamma }_i]`$ as $`𝔽_p`$-subspaces of $`F`$. By replacing $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ with finite-index subgroups, we may assume that these subspaces are contained in $`F^{}`$. Then, because $`\lambda `$ is an isomorphism, we see that the conditions of Notation 2.3 are satisfied, so Theorem 2.4 below implies that there exist
* a standard automorphism $`\varphi _{\tau ,a}`$ of $`G_2`$, and
* a finite-index subgroup $`\mathrm{\Gamma }_1^{}`$ of $`\mathrm{\Gamma }_1`$,
such that $`\lambda (\gamma )\varphi _{\tau ,a}(\gamma )Z(G)`$, for all $`\gamma \mathrm{\Gamma }_1^{}`$.
Because $`\varphi _{\tau ,a}(\mathrm{\Gamma }_1)`$ is an arithmetic lattice, it is commensurable with $`\mathrm{\Gamma }_2`$. Thus, replacing $`\mathrm{\Gamma }_1^{}`$ with a finite-index subgroup, we may assume that $`\varphi _{\tau ,a}(\mathrm{\Gamma }_1^{})\mathrm{\Gamma }_2`$. Then we may define $`\zeta :\mathrm{\Gamma }_1^{}Z(\mathrm{\Gamma }_2)`$ by $`\zeta (\gamma )=\lambda (\gamma )\varphi _{\tau ,a}(\gamma )^1`$. ∎
###### 2.1 Lemma.
Let
* $`\mathrm{\Gamma }`$ be a lattice in a totally disconnected, locally compact group $`G`$,
* $`A`$ be a locally compact, abelian group, and
* $`\zeta :\mathrm{\Gamma }A`$ be a homomorphism.
Assume
1. there is a finite-index subgroup $`\mathrm{\Gamma }^{}`$ of $`\mathrm{\Gamma }`$, such that $`\mathrm{\Gamma }^{}[G,G][\mathrm{\Gamma },\mathrm{\Gamma }]`$, and
2. $`\mathrm{\Gamma }[G,G]`$ is a lattice in $`[G,G]`$.
Then there is a finite-index, open subgroup $`\widehat{G}`$ of $`G`$, such that $`\zeta `$ extends to a continuous homomorphism $`\widehat{\zeta }:\widehat{G}A`$ that is trivial on $`[G,G]`$.
###### Proof.
By assumption, there exists a lattice $`\mathrm{\Gamma }^{}\mathrm{\Gamma }`$ such that $`\mathrm{\Gamma }^{}[G,G][\mathrm{\Gamma },\mathrm{\Gamma }]`$. Since $`\zeta :\mathrm{\Gamma }A`$, and $`A`$ is abelian, we see that $`[\mathrm{\Gamma },\mathrm{\Gamma }]\mathrm{ker}\zeta `$. Therefore $`[\mathrm{\Gamma },\mathrm{\Gamma }]\mathrm{ker}\zeta `$, so, by the choice of $`\mathrm{\Gamma }^{}`$, we have $`\mathrm{\Gamma }^{}[G,G]\mathrm{ker}\zeta `$.
By assumption, $`\mathrm{\Gamma }[G,G]`$ is a lattice in $`[G,G]`$, so $`\mathrm{\Gamma }[G,G]/[G,G]`$ is closed \[Rag, Thm. 1.13, p. 23\], hence discrete. Thus, there is an open compact subgroup $`K/[G,G]G/[G,G]`$, such that $`K(\mathrm{\Gamma }^{}[G,G])=e`$. Let $`\widehat{G}=\mathrm{\Gamma }^{}K[G,G]`$, and extend $`\zeta |_\mathrm{\Gamma }^{}`$ to a homomorphism $`\widehat{\zeta }:\widehat{G}^{}A`$ by defining it to be trivial on $`K[G,G]`$. ∎
###### Proof of Corollary 1.12.
We may assume $`r>2`$. (Otherwise, we must have $`r=1`$, which means $`G_2`$ is abelian, so Proposition 1.6 applies.) From Theorem 1.11, we may assume there exist
* a standard automorphism $`\varphi _{\tau ,a}`$ of $`G_2`$, and
* a homomorphism $`\zeta :\mathrm{\Gamma }_1Z(\mathrm{\Gamma }_2)`$,
such that $`\lambda (\gamma )=\varphi _{\tau ,a}(\gamma )\zeta (\gamma )`$, for all $`\gamma \mathrm{\Gamma }_1`$. From Lemma 2.1, we may assume that there is a finite-index subgroup $`G_2^{}`$ of $`G_2`$, such that $`G_2^{}`$ contains $`[G_2,G_2]`$, and $`\zeta `$ extends to a homomorphism $`\widehat{\zeta }:G_2^{}Z(G_2)`$ that is trivial on $`[G_2,G_2]`$. Let $`G_2^{\prime \prime }=\varphi _{\tau ,a}(G_2^{})`$.
Define $`\widehat{\lambda }:G_2^{}G_2`$ by $`\widehat{\lambda }(g)=\varphi _{\tau ,a}(g)\widehat{\zeta }(g)`$, for $`gG_2^{}`$, so $`\widehat{\lambda }`$ is a continuous homomorphism that extends $`\lambda `$. Because $`\widehat{\zeta }`$ is trivial on $`[G_2,G_2]`$, we know that $`\widehat{\lambda }_{[G_2,G_2]}=\varphi _{\tau ,a}|_{[G_2,G_2]}`$. Also, because $`\widehat{\zeta }(G_2^{})Z(G_2)=[G_2,G_2]`$, we know that $`\widehat{\lambda }(g)\varphi _{\tau ,a}(g)[G_2,G_2]`$ for all $`gG_2^{}`$. Thus, $`\widehat{\lambda }`$ induces an automorphism of $`[G_2,G_2]`$, and an isomorphism $`G_2^{}/[G_2,G_2]G_2^{\prime \prime }/[G_2,G_2]`$, so $`\widehat{\lambda }`$ is an isomorphism. ∎
### 2A. Using linear algebra to prove Theorem 1.11
The remainder of this section is devoted to the statement and proof of Theorem 2.4. This result is a reformulation of Theorem 1.11 in terms of linear algebra. The reformulation is not of intrinsic interest, but it clarifies the essential ideas of the proof, and provides more flexibility, by allowing us to focus on the important aspects of the internal structure of $`\mathrm{\Gamma }`$ that arise from the structure of $`F^{}`$ as a polynomial algebra, without being constrained by the external structure imposed by the group-theoretic embedding of $`\mathrm{\Gamma }`$ in $`G_2`$.
###### 2.2 Notation.
Define an $`𝔽_p`$-bilinear form $`[[,]]:F^{}\times F^{}F^{}`$ by
$$[[a,b]]=a^rbab^r.$$
For any $`V,WF^{}`$, $`[[V,W]]`$ denotes the $`𝔽_p`$-subspace of $`F^{}`$ spanned by $`\{[[v,w]]vV,wW\}`$.
###### 2.3 Notation.
Throughout the remainder of this section, we assume that
* $`r>2`$;
* $`V_1`$ and $`V_2`$ are $`𝔽_p`$-subspaces of finite codimension in $`F^{}`$; and
* $`\lambda ^{}:V_1V_2`$ and $`\lambda _{}:[[V_1,V_1]][[V_2,V_2]]`$ are $`𝔽_p`$-linear bijections,
such that
$$\lambda _{}[[a,b]]=[[\lambda ^{}(a),\lambda ^{}(b)]],$$
for all $`a,bV_1`$.
###### 2.4 Theorem.
There exist
* a subspace $`V_1^{}`$ of finite codimension in $`V_1`$,
* $`ab^1F^{}`$, for some $`bF^{}`$,
* $`\alpha ,\beta 𝔽_q`$, with $`\alpha 0`$, and
* $`\sigma \mathrm{Gal}(𝔽_q/𝔽_p)`$,
such that
$$\lambda ^{}\left(f(t^1)\right)=a\sigma \left(f(\alpha t^1+\beta )\right),$$
for all $`f(t^1)V_1^{}`$.
Let us outline the proof of Theorem 2.4, assuming, for simplicity, that $`V_1=V_2=F^{}`$. For any power $`Q>1`$ of $`r`$, we may define an equivalence relation on $`F^{}\{0\}`$ by $`a_Qb`$ iff $`a/bF^Q`$; let $`[a]`$ denote the equivalence class of $`a`$. For each $`aF^{}`$, the subspace $`[[a,F^{}]]`$ has infinite codimension in $`[[F^{},F^{}]]`$, but Proposition 2.6 shows that $`[[[a],F^{}]]`$ has finite codimension. Because Corollary 2.10 shows that $`\lambda ^{}\left([a]\right)=[\lambda ^{}(a)]`$, this codimension is a useful invariant. Proposition 2.12 shows that it is closely related to the minimum degree of the elements of $`[a]`$. Using this, Corollary 2.22 shows that there is some $`aF^{}`$, a constant $`k`$, and some $`Q`$, such that $`\mathrm{deg}^{}\lambda ^{}(b)=k+\mathrm{deg}^{}b`$ for all $`b_Qa`$. Also, Corollary 2.24 shows that $`\lambda ^{}`$ approximately preserves the degrees of greatest common divisors. Then Proposition 2.25 shows that the restriction of $`\lambda ^{}`$ to the $`𝔽_p`$-rational elements of some equivalence class is of the desired form. Finally, we show that $`\lambda ^{}`$ has the desired form on all of $`F^{}`$.
###### 2.5 Notation.
* We use $`dimW`$ to denote the dimension of a vector space $`W`$ over $`𝔽_p`$.
* Let $`s=dim𝔽_q`$, so $`q=p^s`$.
* For $`a=_{i=0}^n\alpha _it^iF^{}`$, with each $`\alpha _i𝔽_q`$, we let $`\mathrm{deg}^{}a=n`$ if $`\alpha _n0`$.
The following proposition is used in almost all of the following results. Because (2 $``$) requires the assumption that $`e>2`$, it seems that a different approach will be needed for the exceptional case $`p=e=2`$.
###### 2.6 Proposition.
1. The subspace $`[[V_i,V_i]]`$ has finite codimension in $`F^{}`$.
2. Let $`a,bV_i\{0\}`$ and assume $`a/b𝔽_q`$. The subspace $`[[a,V_i]]+[[b,V_i]]`$ has finite codimension in $`[[V_i,V_i]]`$ if and only if $`a/bF^r`$.
###### Proof.
Because $`[[a,V_i]]`$ and $`[[b,V_i]]`$ have finite codimension in $`[[a,F^{}]]`$ and $`[[b,F^{}]]`$, respectively, we see that $`[[a,V_i]]+[[b,V_i]]`$ has finite codimension in $`[[a,F^{}]]+[[b,F^{}]]`$. Thus, in proving (2), we may assume that $`V_i=F^{}`$.
(1) This follows from our proof of (2 $``$) below.
(2 $``$) There are some nonzero $`u,vF^{}`$, such that $`au^r=bv^r`$. Let $`x=a^rub^rv`$.
We claim that $`x0`$. Otherwise, we have
$$a^{r^21}(au^r)=(a^ru)^r=(b^rv)^r=b^{r^21}(bv^r)=b^{r^21}(au^r),$$
so $`a^{r^21}=b^{r^21}`$. This implies $`a/b𝔽_q`$, which is a contradiction. This completes the proof of the claim.
For any $`yF^{}`$, we have
$`[[a,uy]][[b,vy]]`$ $`=`$ $`(a^ruyau^ry^r)(b^rvybv^ry^r)`$
$`=`$ $`(a^ruyb^rvy)(au^ry^rbv^ry^r)`$
$`=`$ $`xy0,`$
so $`[[a,F^{}]]+[[b,F^{}]]`$ contains $`xF^{}`$, which is of finite codimension in $`F^{}`$.
(2 $``$) We may write $`b`$ (uniquely) in the form $`b=x+y^ra`$, with $`x,yF`$, and such that we may write $`x=\alpha _it^i`$ with $`\alpha _i=0`$ whenever $`i\mathrm{deg}^{}(a)(modr)`$. (Note that we do not assume $`x,yF^{}`$.)
For $`u,vF^{}`$, we have
$`[[a,u]][[b,v]]`$ $`=`$ $`(a^ruau^r)(b^rvbv^r)`$
$`=`$ $`(a^rub^rv)\left(au^r(x+y^ra)v^r\right)`$
$`=`$ $`(a^rub^rv)a(uyv)^rxv^r.`$
Whenever either $`\mathrm{deg}^{}(u)`$ or $`\mathrm{deg}^{}(v)`$ is large, it is obvious that $`\mathrm{deg}^{}(a^rub^rv)`$ is much smaller than $`\mathrm{max}\{\mathrm{deg}^{}(uyv)^r,\mathrm{deg}^{}v^r\}`$. Also, we may assume $`x0`$ (otherwise, we have $`b/a=y^rF^r`$, as desired), and, from the definition of $`x`$, we know that $`\mathrm{deg}^{}x\mathrm{deg}^{}a(modr)`$, so
$$\mathrm{deg}^{}\left(a(uyv)^rxv^r\right)=\mathrm{max}\{\mathrm{deg}^{}\left(a(uyv)^r\right),\mathrm{deg}^{}(xv^r)\}.$$
Therefore, we conclude that
$$\mathrm{deg}^{}\left([[a,u]][[b,v]]\right)\{\mathrm{deg}^{}\left(a(uyv)^r\right),\mathrm{deg}^{}(xv^r)\}$$
must be congruent to either $`\mathrm{deg}^{}(a)`$ or $`\mathrm{deg}^{}(x)`$, modulo $`r`$. Thus, because of our assumption that $`r>2`$, we see that $`[[a,F^{}]]+[[b,F^{}]]`$ does not contain elements of all large degrees, so it does not have finite codimension in $`F^{}`$. Then, from (1), we conclude that it does not have finite codimension in $`[[F^{},F^{}]]`$. ∎
###### 2.7 Corollary.
Let $`a_1,a_2V_i\{0\}`$. We have $`a_1/a_2F^r`$ if and only if there is some nonzero $`bV_1`$, such that the subspace $`[[a_j,V_i]]+[[b,V_i]]`$ has finite codimension in $`[[V_i,V_i]]`$, for $`j=1,2`$.
###### Proof.
($``$) Choose $`ba_1F^rV_i(𝔽_qa_1𝔽_qa_2)`$. Then Proposition 2.6(2) implies the desired conclusion.
($``$) From Proposition 2.6(2), we have $`a_1/bF^r`$ and $`a_2/bF^r`$, so $`a_1/a_2F^r`$. ∎
###### 2.8 Lemma.
Let $`a_1,a_2F^{}`$, and let $`Q>1`$ be a power of $`r`$, such that $`\lambda ^{}\left(a_1(F^{})^QV_1\right)=a_2(F^{})^QV_2`$. Define
* subspaces $`W_1`$ and $`W_2`$ of finite codimension in $`F^{}`$ by $`a_i(F^{})^QV_i=a_iW_i^Q`$;
* $`\mu ^{}:W_1W_2`$ by $`\lambda ^{}(a_1w^Q)=a_2\mu ^{}(w)^Q`$; and
* $`\mu _{}:[[W_1,W_1]][[W_2,W_2]]`$ by $`\lambda _{}(a_1^{r+1}w^Q)=a_2^{r+1}\mu _{}(w)^Q`$.
Then $`\mu ^{}`$ and $`\mu _{}`$ are $`𝔽_p`$-linear bijections, and we have
$$\mu _{}[[a,b]]=[[\mu ^{}(a),\mu ^{}(b)]],$$
for all $`a,bW_1`$.
###### 2.9 Definition.
Let $`Q>1`$ be a power of $`p`$. An element of $`F^{}`$ is *$`Q`$-separable* if it is not divisible by a nonconstant $`Q`$th power.
###### 2.10 Corollary.
Let $`aF^{}`$, and let $`Q>1`$ be a power of $`r`$, such that $`a`$ is $`Q`$-separable. Then there is some $`Q`$-separable $`bF^{}`$, such that $`\lambda ^{}\left(a(F^{})^QV_1\right)=b(F^{})^QV_2`$.
###### Proof.
Assume, for the moment, that $`Q=r`$. For $`a_1,a_2F^{}\{0\}`$, define $`a_1a_2`$ iff $`a_1/a_2F^r`$. For nonzero $`a,bV_1`$, we see, from Notation 2.3, that $`[[a,V_1]]+[[b,V_1]]`$ has finite codimension in $`V_1`$ if and only if $`[[\lambda ^{}(a),V_2]]+[[\lambda ^{}(b),V_2]]`$ has finite codimension in $`V_2`$. Therefore, Corollary 2.7 implies that $`ab`$ iff $`\lambda ^{}(a)\lambda ^{}(b)`$. The equivalence classes are precisely the sets of the form $`c(F^{})^rV_i`$, for some $`r`$-separable $`cF^{}`$, so the desired conclusion is immediate.
We may now assume $`Q>r`$. Let $`Q^{}=Q/r`$. There is some $`Q^{}`$-separable $`a^{}F^{}`$, such that $`aa^{}(F^{})^Q^{}`$. By induction on $`Q`$, we know that there is some $`Q^{}`$-separable $`b^{}F^{}`$, such that $`\lambda ^{}\left(a^{}(F^{})^Q^{}V_1\right)=b^{}(F^{})^Q^{}V_2`$.
From the definition of $`a^{}`$, we know there is some $`a_1F^{}`$, such that $`a=a^{}a_1^Q^{}`$. Then, because $`a`$ is $`Q`$-separable, we know that $`a_1`$ is $`r`$-separable.
Define $`W_1`$, $`W_2`$, $`\mu ^{}`$, and $`\mu _{}`$ as in Lemma 2.8 (with $`Q^{}`$, $`a^{}`$, and $`b^{}`$ in the places of $`Q`$, $`a`$, and $`b`$, respectively). Because $`a_1`$ is $`r`$-separable, we know, from the case $`Q=r`$ in the first paragraph of this proof, that there is some $`r`$-separable $`b_1F^{}`$, such that $`\mu ^{}\left(a_1(F^{})^rW_1\right)=b_1(F^{})^rW_2`$. Therefore
$`\lambda ^{}\left(a(F^{})^QV_1\right)`$ $`=`$ $`\lambda ^{}\left[a^{}\left(a_1(F^{})^r\right)^Q^{}V_1\right]`$
$`=`$ $`\lambda ^{}\left[a^{}\left(a_1(F^{})^rW_1\right)^Q^{}\right]`$
$`=`$ $`a^{}\left[\mu ^{}\left(a_1(F^{})^rW_1\right)\right]^Q^{}`$
$`=`$ $`b^{}\left(b_1(F^{})^rW_2\right)^Q^{}`$
$`=`$ $`b^{}\left(b_1(F^{})^r\right)^Q^{}V_2`$
$`=`$ $`b^{}b_1^Q^{}(F^{})^QV_2,`$
as desired. ∎
###### 2.11 Lemma.
Let $`aV_i`$, let $`Q>1`$ be a power of $`r`$, and let $`k`$ be the codimension of $`V_i`$ in $`F^{}`$. Then there is some nonzero $`bF^{}`$ with $`\mathrm{deg}^{}br^2(k+1)`$, such that $`[[a(F^{})^QV_i,V_i]]`$ contains a codimension-$`2k`$ subspace of the ideal $`a^rb^{Q/r}F^{}`$.
###### Proof.
Choose $`cF^{}𝔽_q`$, such that $`ac^QV_i`$ and $`\mathrm{deg}^{}ck+1`$; let $`b=c^{r^2}c`$. For $`yF^{}`$, we have
$`a^rb^{Q/r}y`$ $`=`$ $`a^r(c^{rQ}c^{Q/r})y`$
$`=`$ $`(a^rc^{rQ}yac^Qy^r)(a^rc^{Q/r}yac^Qy^r)`$
$`=`$ $`[[ac^Q,y]][[a,c^{Q/r}y]]`$
$``$ $`[[ac^Q,F^{}]]+[[a,F^{}]],`$
so $`[[ac^Q,F^{}]]+[[a,F^{}]]`$ contains $`a^rb^{Q/r}F^{}`$.
Because $`[[ac^Q,V_i]]`$ and $`[[a,V_i]]`$ contain codimension-$`k`$ subspaces of $`[[ac^Q,F^{}]]`$ and $`[[a,F^{}]]`$, respectively, this implies that $`[[ac^Q,V_i]]+[[a,V_i]]`$ contains a codimension-$`2k`$ subspace of $`a^rb^{Q/r}F^{}`$. Because both $`ac^Q`$ and $`a`$ belong to $`a(F^{})^QV_i`$, the desired conclusion follows. ∎
###### 2.12 Proposition.
Let $`aV_i`$, let $`Q>1`$ be a power of $`r`$, and let $`k`$ be the codimension of $`V_i`$ in $`F^{}`$. Then
$$dim\frac{F^{}}{[[a(F^{})^QV_i,V_i]]}=s(r1)(\mathrm{deg}^{}a)+S+X,$$
where
* $`S=s\mathrm{max}\{\mathrm{deg}^{}cc^r|a\text{}cF^{}\}`$, and
* $`0Xsr(k+1)Q+3k`$.
###### Proof.
Choose $`b`$ as in Lemma 2.11, and let $`I=a^rb^{Q/r}F^{}`$ and $`\overline{F^{}}=F^{}/I`$. It suffices to show
(2.13)
$$dim\overline{F^{}/[[a(F^{})^Q,F^{}]]}s(r1)(\mathrm{deg}^{}a)+S$$
and
(2.14)
$$dim\overline{F^{}/[[a,F^{}]]}S+sr^2(k+1)Q/r+s(r1)\mathrm{deg}^{}a.$$
Let $`u_1,u_2,\mathrm{},u_N`$ be the irreducible factors of $`a^rb^{Q/r}`$. Then we may write
$$a=u_1^{m_1}u_2^{m_2}\mathrm{}u_f^{m_N},b^{Q/r}=u_1^{\epsilon _1}u_2^{\epsilon _2}\mathrm{}u_f^{\epsilon _N},\text{ and }a^rb^{Q/r}=u_1^{n_1}u_2^{n_2}\mathrm{}u_f^{n_N},$$
where $`n_j=rm_j+\epsilon _j`$.
From the Chinese Remainder Theorem, we know that the natural ring homomorphism from $`\overline{F^{}}`$ to
$$\underset{j=1}{\overset{N}{}}\frac{F^{}}{u_j^{n_j}F^{}}$$
is an isomorphism. Thus, we may work in each factor $`F^{}/u_j^{n_j}F^{}`$, and add up the resulting codimensions.
Define $`\varphi _j:F^{}F^{}/(u_j^{rm_j}F^{})`$ by $`\varphi _j(x)=ax^r`$. Then, letting $`m_j^{}=m_jm_j/r`$, we have
$$\mathrm{ker}\varphi _j=\{xF^{}u_j^{m_j^{}}|x\},$$
so
$`dim{\displaystyle \frac{F^{}}{u_j^{rm_j}F^{}+a(F^{})^r}}`$ $`=`$ $`dim{\displaystyle \frac{\mathrm{ker}\varphi _j}{u_j^{rm_j}F^{}}}`$
$`=`$ $`sdim_{𝔽_q}{\displaystyle \frac{\mathrm{ker}\varphi _j}{u_j^{rm_j}F^{}}}`$
$`=`$ $`s(rm_jm_j^{})\mathrm{deg}^{}u_j`$
$`=`$ $`s(r1)\mathrm{deg}^{}u_j^{m_j}+sm_j/r\mathrm{deg}^{}u_j.`$
We have $`a^ru_j^{rm_j}F^{}`$, so
(2.15)
$$[[a(F^{})^Q,F^{}]]a^r(F^{})^{Qr}F^{}+a(F^{})^Q(F^{})^ru_j^{rm_j}F^{}+a(F^{})^r$$
and
(2.16)
$$[[a,F^{}]]+u_j^{rm_j}F^{}=u_j^{rm_j}F^{}+a(F^{})^r.$$
From (2.15), we have
$`dim{\displaystyle \frac{F^{}}{[[a(F^{})^Q,F^{}]]+u_j^{n_j}F^{}}}`$ $``$ $`dim{\displaystyle \frac{F^{}}{[[a(F^{})^Q,F^{}]]+u_j^{rm_j}F^{}}}`$
$``$ $`dim{\displaystyle \frac{F^{}}{u_j^{rm_j}F^{}+a(F^{})^r}}`$
$`=`$ $`s(r1)\mathrm{deg}^{}u_j^{m_j}+sm_j/r\mathrm{deg}^{}u_j,`$
so
$`dim\overline{F^{}/[[a(F^{})^Q,F^{}]]}`$ $``$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}\left(s(r1)\mathrm{deg}^{}u_j^{m_j}+sm_j/r\mathrm{deg}^{}u_j\right)`$
$`=`$ $`s(r1)\mathrm{deg}^{}a+S.`$
This establishes (2.13).
Because $`dim(u_j^{pm_j}F^{}/u_j^{n_j}F^{})=s\epsilon _j\mathrm{deg}^{}u_j`$, and from (2.16), we have
$`dim{\displaystyle \frac{F^{}}{[[a,F^{}]]+u_j^{n_j}F^{}}}`$ $``$ $`dim{\displaystyle \frac{F^{}}{[[a,F^{}]]+u_j^{rm_j}F^{}}}+s\epsilon _j\mathrm{deg}^{}u_j`$
$`=`$ $`dim{\displaystyle \frac{F^{}}{u_j^{rm_j}F^{}+a(F^{})^r}}+s\epsilon _j\mathrm{deg}^{}u_j`$
$`=`$ $`s(r1)\mathrm{deg}^{}u_j^{m_j}+sm_j/r\mathrm{deg}^{}u_j+s\epsilon _j\mathrm{deg}^{}u_j,`$
so
$`dim\overline{F^{}/[[a,F^{}]]}`$ $``$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}\left(s(r1)\mathrm{deg}^{}u_j^{m_j}+sm_j/r\mathrm{deg}^{}u_j+s\epsilon _j\mathrm{deg}^{}u_j\right)`$
$`=`$ $`s(r1)\mathrm{deg}^{}a+S+s\mathrm{deg}^{}b^{Q/r}`$
$``$ $`s(r1)\mathrm{deg}^{}a+S+sr^2(k+1)Q/r.`$
This establishes (2.14). ∎
###### 2.17 Lemma.
For any $`aF^{}`$ and any $`n0`$, we have
$$[[a,F^{}]]+[[1,F^{}]][[a^{r^n},F^{}]]+[[1,F^{}]].$$
###### Proof.
For any $`vF^{}`$, we have
$`[[v,a]]`$ $`=`$ $`v^rava^r`$
$`=`$ $`v^ra(v^{r^2}a^rv^{r^2}a^r)(v^ra^{r^2}v^ra^{r^2})va^r`$
$`=`$ $`[[v^ra,1]]+[[v^r,a^r]]+[[va^r,1]]`$
$``$ $`[[1,F^{}]]+[[a^r,F^{}]].`$
Then the proof is completed by induction on $`n`$. ∎
###### 2.18 Proposition.
There is some $`N`$ (depending only on the codimensions of $`V_1`$ and $`V_2`$, not on the choice of $`V_1`$, $`V_2`$, $`\lambda ^{}`$, or $`\lambda _{}`$), such that $`\mathrm{deg}^{}\lambda ^{}(1)N`$.
###### Proof.
Let $`k`$ be the codimension of $`V_1`$. Choose a power $`Q>1`$ of $`r`$ so large that $`\lambda ^{}(1)`$ is $`Q`$-separable. Then Corollary 2.10 implies $`\lambda ^{}\left((F^{})^QV_1\right)=\lambda ^{}(1)(F^{})^QV_2`$.
Choose $`cF^{}𝔽_q`$, such that $`c^QV_1`$ and $`\mathrm{deg}^{}cr+1`$. We have
$$\begin{array}{ccc}[[(F^{})^QV_1,V_1]]& & [[1,V_1]]+[[c^Q,V_1]]\\ & & [[1,F^{}]]+[[c^Q,F^{}]]\\ & & [[1,F^{}]]+[[c^r,F^{}]]& \text{(see }\text{2.17}\text{)}\\ & & (c^{r^2}c)F^{}& \text{(proof of (}\text{2.11}\text{))}.\end{array}$$
So $`[[(F^{})^QV_1,V_1]]`$ has small codimension in $`[[V_1,V_1]]`$. Therefore $`[[\lambda ^{}(1)(F^{})^QV_2,V_2]]=\lambda _{}[[(F^{})^QV_1,V_1]]`$ must have small codimension in $`[[V_2,V_2]]`$, so $`\mathrm{deg}^{}\lambda ^{}(1)`$ must be small, as desired. ∎
###### 2.19 Corollary.
There is some $`N`$ (depending only on the codimensions of $`V_1`$ and $`V_2`$, not on the choice of $`V_1`$, $`V_2`$, $`\lambda ^{}`$, or $`\lambda _{}`$), such that, for every power $`Q>1`$ of $`r`$ and every $`Q`$-separable element $`a`$ of $`V_1`$, we have $`\mathrm{deg}^{}\lambda ^{}(a)\mathrm{deg}^{}a^{}QN`$, where $`a^{}`$ is the $`Q`$-separable element of $`\lambda ^{}(a)F^Q`$.
###### Proof.
Apply Proposition 2.18 to the map $`\mu ^{}`$ of Lemma 2.8. ∎
###### 2.20 Proposition.
There is a power $`Q>1`$ of $`r`$, and some $`d>0`$, such that, for every $`vV_i`$ with $`\mathrm{deg}^{}v>d`$, there are $`Q`$-separable elements $`v_1,\mathrm{},v_m`$ of $`V_i`$, such that $`v=v_1+\mathrm{}+v_m`$ and $`\mathrm{deg}^{}v_j\mathrm{deg}^{}v`$, for $`j=1,\mathrm{},m`$.
###### Proof.
Let $`k`$ be the codimension of $`V_i`$ in $`F^{}`$, and choose $`Q>k+4`$ so large that, for every $`mQ`$, the subspace $`V_i`$ contains elements of degree $`m`$ whose leading coefficients span $`𝔽_q`$. For any element of $`V_1`$ of degree $`m`$, we show that there is a $`Q`$-separable element of $`V_i`$ of degree $`m`$ with the same leading coefficient.
Let $`\alpha `$ be the leading coefficient of some element of $`V_1`$ of degree $`m`$. Then $`V_i`$ contains exactly $`r^{mk}`$ elements of degree $`m`$ with leading coefficient $`\alpha `$.
On the other hand, if $`a`$ is an element of $`F^{}`$ that is of degree $`m`$ and is not $`Q`$-separable, then $`a`$ must be of the form $`a=x^Qy`$, where $`x`$ is an element of $`F^{}`$ of some degree $`j`$, and $`y`$ is an element of $`F^{}`$ of degree $`mQj`$. Thus, the number of such elements $`a`$ of degree $`m`$ is no more than
$$\underset{j=1}{\overset{\mathrm{}}{}}q^{j+1}q^{mQj+1}=q^{m+2}\underset{j=1}{\overset{\mathrm{}}{}}q^{j(1Q)}=\frac{q^{m+2}}{q^{Q1}1}\frac{q^{m+2}}{q^{Q2}}<Q^{mQ+4}<\frac{Q^m}{r^k}.$$
Therefore, not every element of $`V_i`$ of degree $`m`$ whose leading coefficient is $`\alpha `$ can be such an element $`a`$, so $`V_i`$ has a $`Q`$-separable element of degree $`m`$ with leading term $`\alpha `$, as desired. ∎
###### 2.21 Corollary.
For each $`bF^{}`$, there exists $`N`$, such that, for every $`ab(F^{})^rV_1`$, we have $`|\mathrm{deg}^{}\lambda ^{}(a)\mathrm{deg}^{}a|N`$.
###### Proof.
By symmetry, it suffices to show $`\mathrm{deg}^{}\lambda ^{}(a)\mathrm{deg}^{}a+N`$. We may assume $`b`$ is $`r`$-separable. By combining Proposition 2.20 with Lemma 2.8, we may choose a power $`Q>1`$ of $`r`$, such that each element of $`b(F^{})^r`$ is a sum of $`Q`$-separable elements of $`b(F^{})^r`$ of smaller degree. Thus, we may assume $`a`$ is $`Q`$-separable (and our bound $`N`$ may depend on $`Q`$).
Define $`S`$ as in the statement of Proposition 2.12, and let $`k_i`$ be the codimension of $`V_i`$. Because $`ab(F^{})^r`$ and $`b`$ is $`r`$-separable, we have $`S=s(\mathrm{deg}^{}a\mathrm{deg}^{}b)/r`$, so Proposition 2.12 implies
$$\left|dim\frac{F^{}}{[[a(F^{})^QV_1,V_1]]}s\left(r1+\frac{1}{r}\right)\mathrm{deg}^{}a\right|s\frac{\mathrm{deg}^{}b}{r}+\left(sr(k_1+1)Q+3k_1\right)$$
is bounded. Similarly, letting $`a^{}`$ be the $`Q`$-separable element of $`\lambda ^{}(a)F^Q`$, and $`b^{}`$ be the $`r`$-separable element of $`\lambda ^{}(b)F^r`$, we know that
$$\left|dim\frac{F^{}}{[[a^{}(F^{})^QV_2,V_2]]}s\left(r1+\frac{1}{r}\right)\mathrm{deg}^{}a^{}\right|s\frac{\mathrm{deg}^{}b^{}}{r}+\left(sr(k_2+1)Q+3k_2\right)$$
is bounded. Then, because
$$dim\frac{[[V_1,V_1]]}{[[a(F^{})^QV_1,V_1]]}=dim\frac{[[V_2,V_2]]}{[[a^{}(F^{})^QV_2,V_2]]},$$
we conclude that $`|\mathrm{deg}^{}a^{}\mathrm{deg}^{}a|`$ is bounded. Corollary 2.19 asserts that $`|\mathrm{deg}^{}\lambda ^{}(a)\mathrm{deg}^{}a^{}|`$ is also bounded. ∎
###### 2.22 Corollary.
For each $`bF^{}`$, there is a power $`Q`$ of $`r`$, such that, for every $`a_1,a_2b(F^{})^QV_1`$, we have $`\mathrm{deg}^{}\lambda ^{}(a_1)\mathrm{deg}^{}\lambda ^{}(a_2)=\mathrm{deg}^{}a_1\mathrm{deg}^{}a_2`$.
###### Proof.
Choose $`N`$ as in Corollary 2.21. Now choose $`Q>2N`$. Because
$$\mathrm{deg}^{}\lambda ^{}(a_1)\mathrm{deg}^{}\lambda ^{}(a_2)(modQ)\text{ and }\mathrm{deg}^{}a_1\mathrm{deg}^{}a_2(modQ),$$
we have
$$\mathrm{deg}^{}\lambda ^{}(a_1)\mathrm{deg}^{}a_1\mathrm{deg}^{}\lambda ^{}(a_2)\mathrm{deg}^{}a_2(modQ),$$
so, from the choice of $`N`$ and $`Q`$, we conclude that $`\mathrm{deg}^{}\lambda ^{}(a_1)\mathrm{deg}^{}a_1=\mathrm{deg}^{}\lambda ^{}(a_2)\mathrm{deg}^{}a_2`$. ∎
###### 2.23 Proposition.
There is a constant $`C>0`$, such that, for all $`a_1,a_2V_i`$, and every power $`Q`$ of $`r`$, we have
$`s\mathrm{deg}^{}\mathrm{gcd}(a_1,a_2)C`$ $``$ $`dim{\displaystyle \frac{[[V_i,V_i]]}{[[a_1(F^{})^QV_i,V_i]]+[[a_2(F^{})^QV_i,V_i]]}}`$
$``$ $`C\mathrm{deg}^{}\mathrm{gcd}(a_1,a_2)+C.`$
###### Proof.
Because
$$[[a_1(F^{})^QV_i,V_i]]+[[a_2(F^{})^QV_i,V_i]]\mathrm{gcd}(a_1,a_2)F^{},$$
the left-hand inequality is obvious.
Let $`c=\mathrm{gcd}(a_1,a_2)`$ and let $`k`$ be the codimension of $`V_i`$. Then Lemma 2.11 implies that there exist nonzero $`b_1,b_2F^{}`$ with $`\mathrm{deg}^{}b_ir^2(k+1)`$, such that $`[[a_j(F^{})^QV_i,V_i]]`$ contains a codimension-$`2k`$ subspace of $`a_j^rb_j^{Q/r}F^{}`$ for $`j=1,2`$. Then, letting $`b=b_1b_2`$, we have $`\mathrm{deg}^{}b2r^2(k+1)`$, and $`[[a_1(F^{})^QV_i,V_i]]+[[a_2(F^{})^QV_i,V_i]]`$ contains a codimension-$`4k`$ subspace of the ideal $`I=c^rb^{Q/r}F^{}`$. Thus, it suffices to show that the codimension of $`[[a_1,F^{}]]+[[a_2,F^{}]]+I`$ in $`F^{}`$ is bounded above by $`s(r+2)\mathrm{deg}^{}c+s\mathrm{deg}^{}b`$.
Let $`u_1,\mathrm{},u_N`$ be the irreducible factors of $`c^rb^{Q/r}`$, so we may write $`c=u_1^{m_1}\mathrm{}u_N^{m_N}`$, $`b=u_1^{\epsilon _1}\mathrm{}u_N^{\epsilon _N}`$, and $`c^rb^{Q/r}=u_1^{n_1}\mathrm{}u_N^{n_N}`$, where $`n_j=rm_j+\epsilon _jQ/r`$. From the Chinese Remainder Theorem, we have $`F^{}/I_{j=1}^NF^{}/u^{n_j}F^{}`$, so we may calculate the codimension in each factor, and then add them up.
Fix $`j`$. By interchanging $`a_1`$ and $`a_2`$ if necessary, we may assume that $`u_j^{m_j+1}a_1`$. It suffices to show that
$$dim\frac{F^{}}{[[a_1,F^{}]]+u^{n_j}F^{}}s\left((r+2)m_j+\epsilon _j\right)\mathrm{deg}^{}u_j;$$
thus (because $`m_j+\epsilon _j1`$), we need only show that $`u_j^{(r+1)m_j+1}F^{}[[a_1,F^{}]]+u_j^{n_j}F^{}`$. To show this, let $`M`$ be minimal, such that $`u_j^{M+1}F^{}[[a,F^{}]]+u_j^{n_j}F^{}`$. (Obviously, we have $`M<n_j`$; we wish to show $`M(r+1)m_j`$.) Suppose $`M>(r+1)m_j`$. (This will lead to a contradiction.) We have $`m_j+r(Mrm_j)>M`$, so
$`u_j^MF^{}`$ $`=`$ $`u_j^{rm_i}u_j^{Mrm_j}F^{}`$
$``$ $`a_1^ru_j^{Mrm_j}F^{}+u^{n_j}F^{}`$
$``$ $`[[a_1,u_j^{Mrm_j}F^{}]]+a_1u_j^{r(Mrm_j)}F^{}+u^{n_j}F^{}`$
$``$ $`[[a_1,F^{}]]+u_j^{m_j+r(Mrm_j)}F^{}+u^{n_j}F^{}`$
$``$ $`[[a_1,F^{}]]+u_j^{M+1}F^{}+u^{n_j}F^{}`$
$`=`$ $`[[a_1,F^{}]]+u^{n_j}F^{}.`$
This contradicts the minimality of $`M`$. ∎
###### 2.24 Corollary.
There is a constant $`C>0`$, such that, for all $`a,bV_1`$, we have
$$\frac{\mathrm{deg}^{}\mathrm{gcd}(a,b)}{C}C\mathrm{deg}^{}\mathrm{gcd}(\lambda ^{}(a),\lambda ^{}(b))C\mathrm{deg}^{}\mathrm{gcd}(a,b)+C.$$
###### 2.25 Proposition.
There exist $`bV_1`$, $`b^{}V_2`$, $`\alpha ,\beta 𝔽_q`$, and some $`Q`$ that is a power of both $`r`$ and $`q`$, such that, for all $`bf(t^Q)b\left(𝔽_p[t^1]\right)^QV_1`$, we have $`\lambda ^{}\left(bf(t^Q)\right)=b^{}f(\alpha t^Q+\beta )`$.
###### Proof.
Corollary 2.22 shows that, by replacing $`V_1`$ with some $`(F^{})^QV_1`$ (using Lemma 2.8), we may assume $`\mathrm{deg}^{}\lambda ^{}(a)=\mathrm{deg}^{}a`$, for every $`aV_1`$.
The terms $`C`$ and $`+C`$ in Corollary 2.24 are significant only when $`\mathrm{deg}^{}\mathrm{gcd}(a,b)`$ is small. On the other hand, $`\mathrm{deg}^{}\mathrm{gcd}(a,b)`$ can never be small (and nonzero) if $`a,b(F^{})^Q`$ for some large $`Q`$. Thus, by replacing $`V_1`$ with some $`(F^{})^QV_1`$ (using Lemma 2.8), we may assume
$$\frac{1}{C}\mathrm{deg}^{}\mathrm{gcd}(a,b))\mathrm{deg}^{}\mathrm{gcd}(\lambda ^{}(a),\lambda ^{}(b))C\mathrm{deg}^{}\mathrm{gcd}(a,b),$$
for every $`a,bV_1`$. In particular, $`\mathrm{gcd}(a,b)=1`$ if and only if $`\mathrm{gcd}(\lambda ^{}(a),\lambda ^{}(b))=1`$.
Let $`k`$ be the codimension of $`V_1`$ in $`F^{}`$. Choose some
$$N>4\left(C(C+k)p^{C+k+1}+k+1\right).$$
Choose a power $`Q`$ of $`r`$, such that $`Q>Nk`$. There is some nonzero $`b𝔽_p[t^1]`$, with $`\mathrm{deg}^{}bNk`$, such that
$$b(𝔽_p+t^Q𝔽_p+t^{2Q}𝔽_p+\mathrm{}+t^{NQ}𝔽_p)V_1.$$
Because $`\mathrm{deg}^{}b<Q`$, we know that $`b`$ is $`Q`$-separable, so, by applying Lemma 2.8 to $`b(F^{})^QV_1`$, we may assume
$$𝔽_p+t^1𝔽_p+t^2𝔽_p+\mathrm{}+t^N𝔽_pV_1.$$
By composing $`\lambda ^{}`$ with a map of the form $`f(t^1)\gamma f(\alpha t^1+\beta )`$, for some $`\alpha ,\beta ,\gamma 𝔽_q`$ (with $`\alpha \gamma 0`$), we may assume $`\lambda ^{}(1)=1`$ and $`\lambda ^{}(t^1)=t^1`$, so $`\lambda ^{}|_{𝔽_p+𝔽_pt^1}=\mathrm{Id}`$.
Let $`V_1^{𝔽_p}=V_1𝔽_p[t^1]`$. It suffices to show $`\lambda ^{}(a)=a`$ for every $`aV_1^{𝔽_p}`$.
Suppose $`\lambda ^{}|_{V_1^{𝔽_p}}\mathrm{Id}`$, and let
$$m=\mathrm{min}\left\{\mathrm{deg}^{}a|\lambda ^{}(a)a\text{}aV_1^{𝔽_p}\right\}2.$$
Let $`\mathrm{\Delta }=\lambda ^{}(a)a`$, for any monic $`aV_1^{𝔽_p}`$ with $`\mathrm{deg}^{}a=m`$. (Note that the definition of $`m`$ implies that $`\mathrm{\Delta }`$ is independent of the choice of $`a`$.)
Case 1 . Assume $`mN`$. Let $`u`$ be any irreducible element of $`𝔽_p[t^1]`$ with $`\mathrm{deg}^{}um1`$.
We claim that $`V_1^{𝔽_p}`$ contains a (monic) element $`a`$, such that $`\mathrm{deg}^{}a=m`$ and $`u|a`$. To see this, let $`bV_1^{𝔽_p}`$ with $`\mathrm{deg}^{}b=m`$. There is some $`aF^{}`$, such that $`u|a`$ and $`\mathrm{deg}^{}(ab)<\mathrm{deg}^{}u<\mathrm{deg}^{}b`$. Because $`\mathrm{deg}^{}bN`$, this implies $`abV_1^{𝔽_p}`$, so $`aV_1^{𝔽_p}`$.
Because $`u|a`$ (and $`\lambda ^{}(u)=u`$), we know $`\mathrm{gcd}(u,\lambda ^{}(a))1`$. Because $`u`$ is irreducible, we conclude that $`u|\lambda ^{}(a)`$. We also have $`u|a`$, so this implies $`u|(\lambda ^{}(a)a)=\mathrm{\Delta }`$.
Thus, we see that $`\mathrm{\Delta }`$ is divisible by every irreducible polynomial over $`𝔽_p`$ of degree $`m1`$, so $`\mathrm{\Delta }`$ is divisible by $`t^{p^{m1}}t^1`$. Therefore $`\mathrm{deg}^{}\mathrm{\Delta }p^{m1}`$. However, we also know $`\mathrm{deg}^{}\mathrm{\Delta }\mathrm{deg}^{}a=m`$ (and all nonzero polynomials in $`𝔽_2[t^1]`$ are monic, so $`\mathrm{deg}^{}\mathrm{\Delta }<m`$ if $`p=2`$). This is a contradiction.
Case 2 . Assume $`m>N`$. Choose some monic $`aV_1^{𝔽_p}`$, with $`\mathrm{deg}^{}a=m`$. By subtracting a polynomial of degree $`k`$, we may assume $`t^{(k+1)}|a`$; let $`u=a/t^{(k+1)}`$. There is some nonzero $`x𝔽_p[t^1]`$ with $`\mathrm{deg}^{}xk`$, such that $`uxV_1^{𝔽_p}`$. (Note that $`\mathrm{deg}^{}uxk+\mathrm{deg}^{}u<m`$.)
Let
$$𝒞=\{c𝔽_p[t^1]\{0\}\mathrm{deg}^{}c<C\},$$
and
$$b=\underset{\mathrm{deg}^{}cC+k}{}c,$$
so $`\mathrm{deg}^{}b<(C+k)p^{C+k+1}`$. Now, for each $`c𝒞`$, let
$`u_c=(u+c)x`$ and $`u_c^{}={\displaystyle \frac{u_c}{\mathrm{gcd}(u_c,b)}}`$.
For $`c𝒞`$, we have $`\{cx,ct^{(k+1)}\}V_1^{𝔽_p}`$, so $`u_cV_1^{𝔽_p}`$ and $`a+ct^{(r+1)}V_1^{𝔽_p}`$. Also, because $`a=ut^{(k+1)}`$, we have $`(u+c)|(a+ct^{(k+1)})`$. Then, since $`\lambda ^{}(u+c)=u+c`$, we have $`\mathrm{deg}^{}\mathrm{gcd}(\lambda ^{}(a+ct^{(k+1)}),u+c)\left(\mathrm{deg}^{}(u+c)\right)/C`$, so
$`\mathrm{deg}^{}\mathrm{gcd}(\mathrm{\Delta },u_c^{})`$ $``$ $`\mathrm{deg}^{}\mathrm{gcd}(\mathrm{\Delta },u_c)\mathrm{deg}^{}b`$
$`=`$ $`\mathrm{deg}^{}\mathrm{gcd}(\lambda ^{}(a+ct^{(k+1)})(a+ct^{(k+1)}),u_c)\mathrm{deg}^{}b`$
$``$ $`{\displaystyle \frac{\mathrm{deg}^{}(u+c)}{C}}\mathrm{deg}^{}b`$
$``$ $`{\displaystyle \frac{mk1}{C}}(C+k)p^{C+k+1}`$
$``$ $`{\displaystyle \frac{m}{4C}}.`$
Also, for $`c_1,c_2𝒞`$, we have
$$\mathrm{deg}^{}\mathrm{gcd}(u_{c_1},u_{c_2})\mathrm{deg}^{}(u_{c_1}u_{c_2})=\mathrm{deg}^{}\left((c_1c_2)x\right)C+k,$$
so we see that $`\mathrm{gcd}(u_{c_1}^{},u_{c_2}^{})=1`$ whenever $`c_1c_2`$. Thus, we conclude that
$$\mathrm{deg}^{}\mathrm{\Delta }p^C\frac{m}{4C}>m.$$
This is a contradiction. ∎
###### Proof of Theorem 2.4.
Choose $`b,b^{},\alpha ,\beta ,Q`$ as in Proposition 2.25. By replacing $`\lambda ^{}`$ with $`x(b^{})^1\lambda ^{}(bx)`$ and replacing $`\lambda _{}`$ with $`x(b^{})^{(r+1)}\lambda ^{}(b^{r+1}x)`$, we may assume $`b=b^{}=1`$. Then, by composing $`\lambda ^{}`$ and $`\lambda _{}`$ with $`t^1\alpha ^1(t^1\beta )`$, we may assume $`\alpha =1`$ and $`\beta =0`$. Thus,
(2.26) $`\lambda ^{}(a)=a`$ for all $`a𝔽_p[t^Q]V_1`$.
We wish to show that there is some $`\sigma \mathrm{Gal}(𝔽_q/𝔽_p)`$, such that, for every $`aV_1`$, we have $`\lambda ^{}(a)=\sigma (a)`$.
Step 1 . For each $`aV_1`$, there is some $`\sigma \mathrm{Gal}(𝔽_q/𝔽_p)`$, such that $`\lambda ^{}(a)=\sigma (a)`$. Fix $`aV_1`$. Choose $`C`$ as in Corollary 2.24, let $`k`$ be the codimension of $`V_1`$, and choose $`b𝔽_p[t^Q]V_1`$, such that
$$\frac{\mathrm{deg}^{}b}{C}C>Q\left(s\left(\mathrm{deg}^{}a+\mathrm{deg}^{}\lambda ^{}(a)\right)+k\right).$$
Let
$$c=\underset{\sigma \mathrm{Gal}(𝔽_q/𝔽_p)}{}(b\sigma (a))𝔽_p[t^1],$$
and choose some nonzero $`x𝔽_p[t^1]`$, such that $`(cx)^QV_1`$ and $`\mathrm{deg}^{}xk`$.
We have
$$\begin{array}{ccc}Q\left(\mathrm{deg}^{}\mathrm{gcd}(b\lambda ^{}(a),c)+k\right)& & \mathrm{deg}^{}\mathrm{gcd}(b\lambda ^{}(a),(cx)^Q)\\ & =& \mathrm{deg}^{}\mathrm{gcd}(\lambda ^{}(ba),\lambda ^{}\left((cx)^Q\right))& \text{(see }\text{2.26}\text{)}\\ & & \frac{\mathrm{deg}^{}\mathrm{gcd}(ba,(cx)^Q)}{C}C& \text{(choice of }C\text{)}\\ & =& \frac{(\mathrm{deg}^{}b)}{C}C& \text{(}(ba)|c\text{)}\\ & >& Q\left(s\left(\mathrm{deg}^{}a+\mathrm{deg}^{}\lambda ^{}(a)\right)+k\right)& \text{(choice of }b\text{)}.\end{array}$$
Thus, from the definition of $`c`$, we conclude that there is some $`\sigma \mathrm{Gal}(𝔽_q/𝔽_p)`$, such that
$`\mathrm{deg}^{}\mathrm{gcd}(b\lambda ^{}(a),b\sigma (a))`$ $`>`$ $`\mathrm{deg}^{}a+\mathrm{deg}^{}\lambda ^{}(a)`$
$`=`$ $`\mathrm{deg}^{}\sigma (a)+\mathrm{deg}^{}\lambda ^{}(a)`$
$``$ $`\mathrm{deg}^{}\left(\sigma (a)\lambda ^{}(a)\right)`$
$`=`$ $`\mathrm{deg}^{}\left((b\lambda ^{}(a))(b\sigma (a))\right).`$
Therefore $`\left(b\lambda ^{}(a)\right)\left(b\sigma (a)\right)=0`$, so $`\lambda ^{}(a)=\sigma (a)`$.
Step 2 . There is some $`\sigma \mathrm{Gal}(𝔽_q/𝔽_p)`$, such that $`\lambda ^{}(a)=\sigma (a)`$ for every $`aV_1`$. For $`vF^{}`$, let $`\overline{v}`$ denote the leading coefficient of $`v`$. Choose $`bV_1`$, such that $`\overline{b}`$ generates $`𝔽_q`$, that is, $`𝔽_q=𝔽_p[\overline{b}]`$. From Step 2A, we know there is some $`\sigma \mathrm{Gal}(𝔽_q/𝔽_p)`$, such that $`\lambda ^{}(b)=\sigma (b)`$. We show $`\lambda ^{}(a)=\sigma (a)`$ for every $`aV_1`$.
Given $`aV_1`$, choose some $`cV_1`$, such that $`\overline{c}`$ generates $`𝔽_q`$, and such that $`\mathrm{deg}^{}c>\mathrm{max}\{\mathrm{deg}^{}a,\mathrm{deg}^{}b\}`$. From Step 2A, there exist $`\sigma ^{},\sigma ^{\prime \prime }\mathrm{Gal}(𝔽_q/𝔽_p)`$, such that $`\lambda ^{}(c)=\sigma ^{}(c)`$ and $`\lambda ^{}(a+c)=\sigma ^{\prime \prime }(a+c)`$. Because $`\mathrm{deg}^{}c>\mathrm{deg}^{}a`$, we have $`\overline{c}=\overline{a+c}`$ and $`\overline{\lambda ^{}(a+c)}=\overline{\lambda ^{}(c)}`$. Thus, we have
$$\sigma ^{\prime \prime }(\overline{c})=\sigma ^{\prime \prime }(\overline{a+c})=\overline{\sigma ^{\prime \prime }(a+c)}=\overline{\lambda ^{}(a+c)}=\overline{\lambda ^{}(a)+\lambda ^{}(c)}=\overline{\lambda ^{}(c)}=\sigma ^{}(\overline{c}).$$
Because $`\overline{c}`$ generates $`𝔽_q`$, we conclude that $`\sigma ^{\prime \prime }=\sigma ^{}`$. Therefore
$$\lambda ^{}(a)=\lambda ^{}(a+c)\lambda ^{}(c)=\sigma ^{\prime \prime }(a+c)\sigma ^{}(c)=\sigma ^{}(a+c)\sigma ^{}(c)=\sigma ^{}(a).$$
Similarly, we have $`\lambda ^{}(b)=\sigma ^{}(b)`$. Because we also have $`\lambda ^{}(b)=\sigma (b)`$, and $`\overline{b}`$ generates $`𝔽_q`$, we conclude that $`\sigma ^{}=\sigma `$.
Therefore $`\lambda ^{}(a)=\sigma ^{}(a)=\sigma (a)`$, as desired. ∎
## 3. Arithmetic subgroups of Heisenberg groups
###### Proof of Theorem 1.15.
Let $`\mathrm{\Gamma }_1`$, $`\mathrm{\Gamma }_2`$ be finite-index subgroups of $`\mathrm{\Gamma }`$, such that $`\lambda :\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$ is an isomorphism. Let $`\overline{\mathrm{\Gamma }_i}`$, $`i=1,2`$ be the image of $`\mathrm{\Gamma }_i`$ in $`F^{2m}`$ under the projection $`HF^{2m}`$ with kernel $`Z`$. By passing to a finite-index subgroup, we can assume that $`\overline{\mathrm{\Gamma }_i}(F^{})^{2m}`$. Since $`Z(\mathrm{\Gamma }_i)=\mathrm{\Gamma }_iZ`$, we can identify $`\overline{\mathrm{\Gamma }_i}`$ with $`\mathrm{\Gamma }_i/Z(\mathrm{\Gamma }_i)`$, so $`\lambda `$ induces an isomorphism $`\overline{\lambda }:\overline{\mathrm{\Gamma }_1}\overline{\mathrm{\Gamma }_2}`$.
Step 1 . We can assume $`\overline{\lambda }(av)=a\overline{\lambda }(v)`$ for all $`aF^{}`$ and $`v\overline{\mathrm{\Gamma }_1}`$, such that $`av\overline{\mathrm{\Gamma }_1}`$. For each nonzero $`v\overline{\mathrm{\Gamma }_1}`$, let $`A_v=\{aF^{}av\overline{\mathrm{\Gamma }_1}\}`$. Note that $`A_v`$ is a finite-index subgroup of $`F^{}`$. For $`g,h\mathrm{\Gamma }_i`$, we have $`F\overline{g}=F\overline{h}`$ if and only if $`C_{\mathrm{\Gamma }_i}(g)=C_{\mathrm{\Gamma }_i}(h)`$, so $`\overline{\lambda }(A_vv)=F\overline{\lambda }(v)\overline{\mathrm{\Gamma }_2}`$. Thus, we can define a function $`\tau _v:A_vF`$ by $`\tau _v(a)\overline{\lambda }(v)=\overline{\lambda }(av)`$. Let $`w\overline{\mathrm{\Gamma }_1}`$ be such that $`[[v,w]]0`$, and let $`aA_vA_w`$. Then
$`\tau _v(a)[[\overline{\lambda }(v),\overline{\lambda }(w)]]`$ $`=`$ $`[[\overline{\lambda }(av),\overline{\lambda }(w)]]`$
$`=`$ $`\lambda \left([[av,w]]\right)`$
$`=`$ $`\lambda \left([[v,aw]]\right)`$
$`=`$ $`[[\overline{\lambda }(v),\overline{\lambda }(aw)]]`$
$`=`$ $`\tau _w(a)[[\overline{\lambda }(v),\overline{\lambda }(w)]].`$
Thus
(3.1) $`\tau _v=\tau _w`$ on $`A_vA_w`$ whenever $`[[v,w]]0`$.
For any nonzero $`v,w\overline{\mathrm{\Gamma }_1}`$ and any $`aA_vA_w`$, since $`\overline{\mathrm{\Gamma }_1}a^1\overline{\mathrm{\Gamma }_1}`$ is of finite index in $`\overline{\mathrm{\Gamma }_1}`$, we can find $`u\overline{\mathrm{\Gamma }_1}`$ so that $`aA_u`$, $`[[u,v]]0`$, and $`[[u,w]]0`$. Then it follows from Equation (3.1) that $`\tau _v(a)=\tau _u(a)=\tau _w(a)`$. Since $`aA_vA_w`$ was arbitrary, we conclude that
(3.2) $`\tau _v=\tau _w`$ on $`A_vA_w`$, for all nonzero $`v,w\overline{\mathrm{\Gamma }}_1`$.
For an arbitrary $`aF^{}`$ we can always find $`w\overline{\mathrm{\Gamma }_1}`$ so that $`aA_w`$, thus we can define a function $`\tau :F^{}F`$, by $`\tau (a)=\tau _w(a)`$. Equation (3.2) implies that $`\tau `$ is well defined. Note that $`\tau (1)=1`$. Since
$`\tau (a)\tau (b)[[\overline{\lambda }(u),\overline{\lambda }(v)]]`$ $`=`$ $`[[\overline{\lambda }(au),\overline{\lambda }(bv)]]`$
$`=`$ $`\lambda \left([[au,bv]]\right)`$
$`=`$ $`\lambda \left([[abu,v]]\right)`$
$`=`$ $`[[\overline{\lambda }(abu),\overline{\lambda }(v)]]`$
$`=`$ $`\tau (ab)[[\overline{\lambda }(u),\overline{\lambda }(v)]],`$
we have $`\tau (a)\tau (b)=\tau (ab)`$. Since $`\tau `$ is also an additive homomorphism, and $`\overline{\lambda }`$ is an isomorphism, we conclude that $`\tau `$ is a ring automorphism of $`F^{}`$. Therefore $`\tau \left(f(t^1)\right)=\sigma \left(f(\alpha t^1+\beta )\right)`$ for $`f(t^1)F^{}`$, where $`\sigma \mathrm{Gal}(𝔽_q/𝔽_p)`$, $`\alpha 𝔽_q\{0\}`$, and $`\beta 𝔽_q`$. Hence, by composing with the standard automorphism $`T_{\mathrm{Id},\tau ^1}`$, we obtain the claim.
Step 2 . We may assume that $`\lambda |_{Z(\mathrm{\Gamma }_1)}`$ is the identity map. Let $`v_1,w_1,v_2,w_2\overline{\mathrm{\Gamma }}`$ with $`[[v_i,w_i]]0`$. There is a finite-index subgroup $`A`$ of $`F^{}`$, such that $`av_i\overline{\mathrm{\Gamma }}`$, for every $`aA`$ and $`i=1,2`$. Then, for all $`aA`$, Step 3 implies that
$$\frac{\lambda \left(a[[v_i,w_i]]\right))}{a[[v_i,w_i]]}=\frac{\lambda \left([[v_i,w_i]]\right)}{[[v_i,w_i]]}.$$
Thus, choosing $`a_1,a_2A`$, such that $`a_1[[v_1,w_1]]=a_2[[v_2,w_2]]`$, we have
$$\frac{\lambda \left([[v_1,w_1]]\right)}{[[v_1,w_1]]}=\frac{\lambda \left(a_1[[v_1,w_1]]\right)}{a_1[[v_1,w_1]]}=\frac{\lambda \left(a_2[[v_2,w_2]]\right)}{a_2[[v_2,w_2]]}=\frac{\lambda \left([[v_2,w_2]]\right)}{[[v_2,w_2]]}.$$
We conclude that $`\lambda (z)/z=C`$ is constant, for $`z[[\overline{\mathrm{\Gamma }_1},\overline{\mathrm{\Gamma }_1}]]\{0\}`$.
By composing with a standard automorphism $`\varphi _{T,\mathrm{Id}}`$, such that $`c_T=1/C`$, we may assume that $`C=1`$, so $`\lambda |_{[\mathrm{\Gamma }_1,\mathrm{\Gamma }_1]}=\mathrm{Id}`$. Then, by replacing $`\mathrm{\Gamma }_1`$ with a finite-index subgroup $`\mathrm{\Gamma }_1^{}`$, such that $`\mathrm{\Gamma }_1^{}Z[\mathrm{\Gamma }_1,\mathrm{\Gamma }_1]`$, we may assume $`\lambda |_{Z(\mathrm{\Gamma }_1)}=\mathrm{Id}`$.
Step 3 . $`\overline{\lambda }:\overline{\mathrm{\Gamma }_1}\overline{\mathrm{\Gamma }_1}`$ can be extended to a conformally symplectic map $`\overline{\mathrm{\Lambda }}:F^{2m}F^{2m}`$, with $`c_{\overline{\mathrm{\Lambda }}}=1`$. By Step 3, $`\overline{\lambda }(av)=a\overline{\lambda }(v)`$ for all $`aF^{}`$ and $`v\overline{\mathrm{\Gamma }_1}`$ such that $`av\overline{\mathrm{\Gamma }_1}`$. Because $`\overline{\mathrm{\Gamma }_1}`$ is commensurable with $`(F^{})^{2m}`$, this implies that $`\overline{\lambda }`$ extends (uniquely) to an $`F`$-linear map $`\overline{\mathrm{\Lambda }}:F^{2m}F^{2m}`$. For any $`v,w\overline{\mathrm{\Gamma }_1}`$, we have
$$[[\overline{\mathrm{\Lambda }}(v),\overline{\mathrm{\Lambda }}(w)]]=[[\overline{\lambda }(v),\overline{\lambda }(w)]]=\lambda \left([[v,w]]\right)=[[v,w]],$$
by Step 3. Because $`\overline{\mathrm{\Gamma }_1}`$ spans $`F^{2m}`$, this implies that $`\overline{\mathrm{\Lambda }}`$ is conformally symplectic, with $`c_{\overline{\mathrm{\Lambda }}}=1`$.
Step 4 . Completion of the proof. Define $`\widehat{\mathrm{\Lambda }}:HH`$ by $`\widehat{\mathrm{\Lambda }}(v,z)=(\overline{\mathrm{\Lambda }}(v),z)`$. From Step 3, we see that $`\widehat{\mathrm{\Lambda }}`$ is an automorphism. Denote by $`\zeta :\mathrm{\Gamma }_1Z(H)`$ the map defined by $`\zeta (\gamma )=\widehat{\mathrm{\Lambda }}(\gamma )^1\lambda (\gamma )`$. Then $`\zeta `$ is a homomorphism and $`\lambda (\gamma )=\zeta (\gamma )\widehat{\mathrm{\Lambda }}(\gamma )`$, for $`\gamma \mathrm{\Gamma }_1`$. ∎
###### Proof of Corollary 1.16.
From Theorem 1.15, we may assume there exist
* a standard automorphism $`\varphi _{T,\tau }`$ of $`H`$; and
* a homomorphism $`\zeta :\mathrm{\Gamma }_1Z(H)`$,
such that $`\lambda (\gamma )=\varphi _{T,\tau }(\gamma )\zeta (\gamma )`$ for all $`\gamma \mathrm{\Gamma }_1`$. By Lemma 2.1, there exists a finite-index open subgroup $`\widehat{H}`$ of $`H`$, containing $`[H,H]`$, such that $`\zeta `$ extends to $`\widehat{\zeta }:\widehat{H}Z(H)`$. Let $`H^{}=\varphi _{T,\tau }(\widehat{H})`$.
Define $`\widehat{\mathrm{\Lambda }}:\widehat{H}H`$ by $`\widehat{\mathrm{\Lambda }}(h)=\varphi _{T,\tau }(h)\widehat{\zeta }(h)`$, so that $`\widehat{\mathrm{\Lambda }}`$ is a continuous homomorphism virtually extending $`\lambda `$. Because $`\widehat{\zeta }`$ is trivial on $`[H,H]`$, we have $`\widehat{\mathrm{\Lambda }}|_{[H,H]}=\varphi _{T,\tau }|_{[H,H]}`$, so $`\widehat{\mathrm{\Lambda }}|_{[H,H]}`$ is an automorphism. Because $`\widehat{\zeta }(\widehat{H})Z(H)=[H,H]`$, we see that $`\widehat{\mathrm{\Lambda }}`$ induces an isomorphism $`\widehat{H}/[H,H]H^{}/[H,H]`$. So $`\widehat{\mathrm{\Lambda }}:\widehat{H}H^{}`$ is an isomorphism. ∎
###### 3.3 Definition.
Let
$$H_p=\left\{\left(\begin{array}{cccccc}1& x_1^p& x_2^p& \mathrm{}& x_m^p& z\\ & 1& & & & y_1^p\\ & & 1& \text{0}& & y_2^p\\ & & & \mathrm{}& & \mathrm{}\\ & \text{0}& & & 1& y_m^p\\ & & & & & 1\end{array}\right)|\begin{array}{c}x_1,\mathrm{},x_mF,\\ y_1,\mathrm{},y_mF,\\ zF\end{array}\right\}.$$
###### 3.4 Remark.
$`H_p`$ could also be described as the $`F`$-points of the group obtained from $`H`$ by applying the isogeny of factoring by the Lie algebra of $`Z(H)`$ \[Bor, Prop. V.17.4, p. 215\].
###### 3.5 Corollary.
Any arithmetic lattice in $`H_p`$ is automorphism rigid.
###### Proof.
Let $`\lambda _p:\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$ be an isomorphism, where $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ are arithmetic lattices in $`H_p`$. Define
$$H_p^{}=\left\{\left(\begin{array}{cccccc}1& x_1^p& x_2^p& \mathrm{}& x_m^p& z^p\\ & 1& & & & y_1^p\\ & & 1& \text{0}& & y_2^p\\ & & & \mathrm{}& & \mathrm{}\\ & \text{0}& & & 1& y_m^p\\ & & & & & 1\end{array}\right)|\begin{array}{c}x_1,\mathrm{},x_mF,\\ y_1,\mathrm{},y_mF,\\ zF\end{array}\right\}$$
and
$$A=\left\{\left(\begin{array}{cccccc}1& 0& 0& \mathrm{}& 0& z\\ & 1& & & & 0\\ & & 1& \text{0}& & 0\\ & & & \mathrm{}& & \mathrm{}\\ & \text{0}& & & 1& 0\\ & & & & & 1\end{array}\right)|\begin{array}{c}z=\underset{\begin{array}{c}0in\\ i0\left(modp\right)\end{array}}{}\alpha _it^i,\\ \\ n,\\ \alpha _i𝔽_q\end{array}\right\}$$
Then $`H_p=H_p^{}\times A`$. By passing to a finite-index subgroup we may assume that $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_1^{}\times \mathrm{\Gamma }_{1,A}`$, where $`\mathrm{\Gamma }_1^{}=\mathrm{\Gamma }_1H_p^{}`$ and $`\mathrm{\Gamma }_{1,A}=\mathrm{\Gamma }_1A`$. Let $`\mathrm{\Omega }=\lambda _p(\mathrm{\Gamma }_{1,A})Z(\mathrm{\Gamma }_2)`$ and $`\mathrm{\Gamma }_2^{}=\lambda _p(\mathrm{\Gamma }_1^{})`$. Then, by passing to a finite-index subgroup, we may assume $`\mathrm{\Omega }H_p^{}=e`$ and $`\mathrm{\Gamma }_2^{}A=e`$.
Step 1 . Let $`\pi _A:Z(H_p)A`$ denote the projection with kernel $`H_p^{}`$. Then $`\pi _A\lambda _p:\mathrm{\Gamma }_{1,A}\pi _A(\mathrm{\Omega })`$ virtually extends to a virtual automorphism $`\mathrm{\Psi }`$ of $`A`$. It is easy to see that $`\pi _A(Z(\mathrm{\Gamma }_2))`$ is closed in $`A`$ and hence is a lattice. Because $`Z(\mathrm{\Gamma }_1^{})\times \mathrm{\Gamma }_{1,A}`$ has finite index in $`Z(\mathrm{\Gamma }_1)`$, we know $`\lambda _p(Z(\mathrm{\Gamma }_1^{}))\times \lambda _p(\mathrm{\Gamma }_{1,A})`$ has finite index in $`Z(\mathrm{\Gamma }_2)`$. Then, since $`[\mathrm{\Gamma }_1^{},\mathrm{\Gamma }_1^{}]`$ has finite index in $`Z(\mathrm{\Gamma }_1^{})`$ and
$$\lambda _p([\mathrm{\Gamma }_1^{},\mathrm{\Gamma }_1^{}])[\mathrm{\Gamma }_2^{},\mathrm{\Gamma }_2^{}]H_p^{}=\mathrm{ker}\pi _A$$
we conclude that $`\pi _A(\mathrm{\Omega })=\pi _A(\lambda _p(\mathrm{\Gamma }_{1,A}))`$ has finite index in $`\pi _A(Z(\mathrm{\Gamma }_2))`$. Hence $`\pi _A(\mathrm{\Omega })`$ is a lattice in $`A`$. By Proposition 1.6 $`\pi _A\lambda _p:\mathrm{\Gamma }_{1,A}\pi _A(\mathrm{\Omega })`$ virtually extends to a virtual automorphism $`\mathrm{\Psi }`$ of $`A`$.
Step 2 . Let $`\pi ^{}:H_pH_p^{}`$ be the projection with kernel $`A`$, and let $`\mu _p=\pi ^{}\lambda _p|_{\mathrm{\Gamma }_1^{}}:\mathrm{\Gamma }_1^{}\pi ^{}(\mathrm{\Gamma }_2^{})`$. Then $`\mu _p`$ virtual extends to a virtual automorphism of $`H_p^{}`$.
We claim that $`\pi ^{}(\mathrm{\Gamma }_2^{})`$ is an arithmetic lattice in $`H_p^{}`$. Because $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_1^{}\times \mathrm{\Gamma }_{1,A}`$ and $`\mathrm{\Gamma }_{1,A}Z(\mathrm{\Gamma }_1)`$, we have
$$\mathrm{\Gamma }_2=\mathrm{\Gamma }_2^{}\times \mathrm{\Omega }\mathrm{\Gamma }_2^{}Z(H_p).$$
Then, because $`\mathrm{\Gamma }_2^{}\mathrm{\Gamma }_2`$, we conclude that $`\mathrm{\Gamma }_2^{}Z(H_p)=\mathrm{\Gamma }_2^{}Z(H_p)`$ is a lattice in $`H_p/Z(H_p)H_p^{}/Z(H_p^{})`$. So the image of $`\pi ^{}(\mathrm{\Gamma }_2^{})`$ in $`H_p^{}/Z(H_p^{})`$ is a lattice. Also,
$$\pi ^{}(\mathrm{\Gamma }_2^{})Z(H_p^{})[\mathrm{\Gamma }_2^{},\mathrm{\Gamma }_2^{}]=[\mathrm{\Gamma }_2,\mathrm{\Gamma }_2],$$
so $`\pi ^{}(\mathrm{\Gamma }_2^{})Z(H_p^{})`$ is a lattice in $`[H_p,H_p]=Z(H_p^{})`$. Thus, we conclude that $`\pi ^{}(\mathrm{\Gamma }_2^{})`$ is a lattice in $`H_p^{}`$. Because $`\pi ^{}(\mathrm{\Gamma }_2^{})`$ is contained in the arithmetic lattice $`\pi ^{}(\mathrm{\Gamma }_2)`$, this implies that $`\pi ^{}(\mathrm{\Gamma }_2^{})`$ is arithmetic.
From the preceding paragraph, we know that $`\mu _p`$ is an isomorphism of arithmetic lattices in $`H_p^{}`$. Let $`\mathrm{Fr}:HH_p^{}`$ denote the group isomorphism induced by the Frobenius automorphism $`xx^p`$ of the ground field $`F`$. Then there exist arithmetic lattices $`\widehat{\mathrm{\Gamma }_1},\widehat{\mathrm{\Gamma }_2}`$ in $`H`$, such that $`\mathrm{Fr}(\widehat{\mathrm{\Gamma }_1})=\mathrm{\Gamma }_1^{}`$ and $`\mathrm{Fr}(\widehat{\mathrm{\Gamma }_2})=\pi ^{}(\mathrm{\Gamma }_2^{})`$, and an isomorphism $`\lambda =\mathrm{Fr}^1\mu _p\mathrm{Fr}:\widehat{\mathrm{\Gamma }_1}\widehat{\mathrm{\Gamma }_2}`$. By Corollary 1.16, we can virtually extend $`\lambda `$ to a virtual automorphism $`\mathrm{\Lambda }`$ of $`H`$. Then $`\mathrm{\Lambda }_p^{}=\mathrm{Fr}\mathrm{\Lambda }\mathrm{Fr}^1`$ is a virtual automorphism of $`H_p^{}`$ virtually extending $`\mu _p`$.
Let $`\stackrel{~}{\mathrm{\Lambda }}_p=\mathrm{\Lambda }_p^{}\times \mathrm{\Psi }`$, so $`\stackrel{~}{\mathrm{\Lambda }}_p`$ is a virtual automorphism of $`H_p`$. We can define a map $`\zeta `$ on some finite index subgroup of $`\mathrm{\Gamma }_1`$ by $`\zeta (\gamma )=\lambda _p(\gamma )\stackrel{~}{\mathrm{\Lambda }}_p(\gamma )^1`$. By Lemma 2.1, $`\zeta `$ virtually extends to $`\widehat{\zeta }:H_pZ(H_p)`$. Then $`\mathrm{\Lambda }_p=\stackrel{~}{\mathrm{\Lambda }}_p\zeta `$ is a virtual endomorphism of $`H_p`$. Since $`\mathrm{ker}(\zeta )[H_p,H_p]`$ we conclude (much as in the proof of Corollary 1.12) that $`\mathrm{\Lambda }_p`$ is a virtual automorphism. It is easy to see that it virtually extends $`\lambda _p`$. ∎
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# Spin-wave spectrum of a two-dimensional itinerant electron system: Analytic results for the incommensurate spiral phase in the strong-coupling limit
## I Introduction
Spin fluctuations about incommensurate static magnetic configurations represent an interesting problem, especially in the view of the accumulating experimental evidence on the parent compounds of the high-$`T_c`$ cuprate materials. The dynamic structure factor $`S(𝐪,\omega )`$ at wave vector $`𝐪`$ and frequency $`\omega `$, as measured by inelastic neutron scattering, shows noticeable peaks over the background when $`𝐪=𝐐=𝐐_{AF}+\mathrm{\Delta }𝐐`$ ($`|\mathrm{\Delta }𝐐||𝐐_{AF}|`$) for increasing dopant concentration, at finite albeit small values of $`\omega `$ (with $`\omega J`$, $`J`$ being the exchange coupling). Further experimental studies for larger values of $`|𝐪𝐐_{AF}|`$ have also detected the presence of well-defined spin-wave-type excitations close to the boundary of the antiferromagnetic (AF) Brillouin zone (BZ). The fact that spin-wave excitations with small wavelength (i.e., comparable to lattice spacing) can be detected even in the absence of long-range magnetic order, has actually been well established experimentally since the eighties also in more conventional magnetic materials. This state of affairs has, in turn, prompted a number of theoretical studies on dynamical excitations about incommensurate spin configurations.
From the theoretical point of view, inclusion of incommensurate spin configurations in cluster mean-field calculations would require one to consider cluster sizes at least as big as the spatial extension of the incommensurability. In this way, complex incommensurate patterns could as well be included in the calculation, at the price of considerable numerical effort and without full analytic control on the results. Alternatively, one may set up calculations for an infinite system with necessarily simpler incommensurate spin configurations, with the advantage, however, of obtaining analytic results (at least) in some limits, which in turn may admit a simple physical interpretation. In this respect, the *spiral spin configuration* appears to be the only one for which *analytic* calculations can be performed, in the sense that the constituent equations can be brought to a *closed form* which is manageable for controlled analytic approximations. Indeed previous analytic calculations have considered *static* long-range background spiral spin configurations, on top of which dynamical spin excitations have been considered.
The intrinsic simplicity of the background spiral spin configuration and the neglect of its dynamics have resulted, however, into an instability of the spin-wave spectrum in a limited region of the Brillouin zone. This outcome could, in principle, make the resulting spin-wave spectrum altogether unreliable, as the underlying spin configuration (over and above which spin-wave excitations are constructed) could be much more complex than a spiral one and even possess a dynamics of its own. In addition, a static or dynamic charge modulation can be present as well, like the stripe structures observed in some underdoped cuprates. One expects, however, on physical grounds the *short-wavelength* spin-wave excitations, obtained on top of a spiral configuration, to preserve their dispersion relation even when considering more complex long-range underlying structures.
Within this framework, we have pursued the analytic calculation with an underlying *spiral* spin configuration for an itinerant electronic system described by a single-band Hubbard Hamiltonian. The spin (and charge) dynamics have bee described within the electronic random-phase (RPA) approximation for the dynamical susceptibilities, based on a broken-symmetry Hartree-Fock (HF) mean-field solution, along the lines of Ref. . In particular, we have considered the large $`U/t`$ limit in detail (where $`t`$ is the nearest-neighbor hopping matrix element and $`U`$ is the local on-site repulsion in the Hubbard Hamiltonian), for which the spin-wave dispersion relation (as obtained from the poles of the dynamical susceptibilities) can be expressed in analytic form by systematically expanding in the *small parameter* $`t/U`$. In this way, we have obtained an analytic expression for the spin-wave dispersion relation valid in the limit $`t/U1`$ \[cf. Eq. (106) below\], showing a novel characteristic structure, namely, a hybrid form between the dispersion relations obtained within the nearest-neighbor Heisenberg model for localized spins and the long-range RKKY interaction mediated by the conduction electrons.
This analytic form of the dispersion relation could admittedly not have been guessed *a priori*, by fitting the dispersion relation obtained numerically (from the location of the poles of the dynamical susceptibilities) with a Heisenberg model extending in principle to a large albeit *finite* number of neighbors. Neither, this analytic form can be simply reduced to a nearest-neighbor Heisenberg dispersion relation with a doping-dependent exchange integral. Rather, the characteristic long-range RKKY contribution would require fitting to a Heisenberg model with an infinite number of neighbors. This would contradict the spirit with which the Heisenberg model was introduced to start with, namely, as a fitting model that makes physical sense when the interactions extend to a limited number of neighbors only. Note that this situation contrasts that found at half-filling of the Hubbard band (i.e., in the absence of doping), where the antiferromagnetic (AF) spin-wave spectrum can be nicely fitted by a Heisenberg model extending at most to a few neighbors. Our results also show that the magnitude of the overall exchange integral, which characterizes the spin-wave spectrum, decreases with increasing doping and vanishes when the transition to a ferromagnetic case occurs at the mean-field level.
The plan of the paper is as follows. Section 2 obtains the implicit form of the spin-wave dispersion relation within the HF-RPA approximations, by solving for the dynamical susceptibilities of the itinerant electron system in the presence of an incommensurate spin spiral ground state. Section 3 focuses on the small $`t/U`$ expansion of the results of Section 2, which requires a careful analysis of the doping dependence of the relevant quantities. Section 4 discusses the main results of this paper and Section 5 gives our conclusions. For the sake of completeness, we report in the Appendices details of the analytic calculations, and adapt know results for the Heisenberg and RKKY spin-wave spectra to the present context.
## II Dynamical susceptibilities within the itinerant-electron RPA approximation with an incommensurate spin-spiral ground state
In this Section, we give the derivation of the spin-wave dispersion relation for a two-dimensional itinerant-electron system in the presence of an incommensurate spiral magnetic structure *with a generic characteristic wave vector* $`𝐐`$, within the electronic RPA approximation. Although our results for the dispersion relation coincide with those previously given by Brenig in Ref. , we provide here in addition the expression of the correlation functions which can be relevant for a direct comparison with neutron scattering experiments. Further details of the calculation are reported in Appendix A.
We emphasize that the finding of a closed-form expression for the correlation functions for *any* characteristic wave vector $`𝐐`$ (incommensurate with the lattice spacing) and not just for the (commensurate) antiferromagnetic wave vector $`𝐐_{AF}`$ is altogether a nontrivial result, being intrinsically related to the peculiar pattern of the spiral magnetic solution for the ground state.
We begin by considering the generalized *correlation function* at zero temperature in the broken-symmetry phase:
$$𝒳_{\mu ,\nu }(𝐫t,𝐫^{}t^{})=iT[S_\mu (𝐫,t)S_\nu (𝐫^{},t^{})]+iS_\mu (𝐫,t)S_\nu (𝐫^{},t^{})$$
(1)
where the average $`\mathrm{}`$ is taken over the ground state, $`T`$ stands here for Wick’s time-ordering operator, and $`S_\mu (𝐫)`$ is given by (we set $`\mathrm{}=1`$ throughout)
$$S_\mu (𝐫)=\frac{1}{2}\underset{\alpha ,\beta }{}\psi _\alpha ^{}(𝐫)\sigma _{\alpha ,\beta }^\mu \psi _\beta (𝐫).$$
(2)
In these expressions, $`\mu ,\nu =(0,x,y,z)`$, $`\sigma ^\mu `$ is a Pauli matrix (with $`\sigma ^0`$ equal to the $`2\times 2`$ identity matrix), and $`\alpha ,\beta `$ are spin labels. Note that $`2S_0(𝐫)`$ coincides with the density operator, which couples with the spin operator in the presence of an incommensurate spiral magnetic structure.
For the simple band we are considering, the field operator in Eq.(2) can be represented in the form
$$\psi _\alpha (𝐫)=\underset{i}{}\varphi (𝐫𝐑_i)c_{i\alpha },$$
(3)
where $`\varphi (𝐫)`$ is the atomic (Wannier) orbital associated with the simple band, $`𝐑_i`$ is the lattice vector locating site $`i`$, and $`c_{i\alpha }`$ is a destruction operator. Time evolution in Eq. (1) is governed by the Heisenberg representation:
$$\psi _\alpha (𝐫,t)=\mathrm{e}^{iHt}\psi _\alpha (𝐫)\mathrm{e}^{iHt}$$
(4)
where for $`H`$ we take the simple-band two-dimensional Hubbard Hamiltonian.
In terms of the two-particle correlation function $`L`$, the generalized correlation function (1) takes the form
$$𝒳_{\mu ,\nu }(𝐫t,𝐫^{}t^{})=\frac{i}{4}\underset{\alpha ,\beta }{}\underset{\alpha ^{},\beta ^{}}{}\sigma _{\alpha ,\beta }^\mu \sigma _{\alpha ^{},\beta ^{}}^\nu L(𝐫t\beta ,𝐫^{}t^{}\beta ^{};𝐫t^+\alpha ,𝐫^{}t^+\alpha ^{})$$
(5)
with $`t^+=t+\eta `$ ($`\eta =0^+`$), where $`L`$ satisfies the Bethe-Salpeter equation:
$`L(1,2;1^{},2^{})`$ $`=`$ $`G(1,2^{})G(2,1^{})`$ (6)
$`+`$ $`{\displaystyle d3d4d5d6G(1,3)G(4,1^{})\mathrm{\Xi }(3,5;4,6)L(6,2;5,2^{})}`$ (7)
($`1,2,\mathrm{}`$ signifying the set of space, spin, and time variables). In the above expression,
$$G(1,2)=iT\left[\psi (1)\psi ^{}(2)\right]$$
(8)
is the single-particle Green’s function and the kernel $`\mathrm{\Xi }`$ represents an effective two-particle interaction. In particular, within the RPA approximation we are adopting, the kernel $`\mathrm{\Xi }`$ takes the form:
$`\mathrm{\Xi }(3,5;4,6)`$ $`=`$ $`iv_0\delta (3,4)\delta (5,6^+)\delta (x_3,x_6)\delta (\alpha _3,\alpha _6)`$ (9)
$`+`$ $`iv_0\delta (3,6)\delta (4,5)\delta (x_3^+,x_4)\delta (\alpha _3,\alpha _4)`$ (10)
with the notation $`x(𝐫,t)`$ and where the constant $`v_0`$ can be related to the parameter $`U`$ of the Hubbard Hamiltonian as follows:
$$U=v_0𝑑𝐫|\varphi (𝐫)|^4.$$
(11)
Entering Eqs.(6) and (10) into Eq.(5), we obtain for the generalized correlation function within the RPA approximation:
$`𝒳_{\mu ,\nu }(x,x^{})`$ $`=`$ $`{\displaystyle \frac{i}{4}}{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle \underset{\alpha ^{},\beta ^{}}{}}\sigma _{\alpha ,\beta }^\mu \sigma _{\alpha ^{},\beta ^{}}^\nu G(x\beta ,x^{}\alpha ^{})G(x^{}\beta ^{},x\alpha )`$ (12)
$`+`$ $`{\displaystyle \frac{v_0}{4}}(i)^2{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle \underset{\alpha ^{},\beta ^{}}{}}\sigma _{\alpha ,\beta }^\mu \sigma _{\alpha ^{},\beta ^{}}^\nu {\displaystyle d3G(x\beta ,x_3\alpha _3)G(x_3\alpha _3,x\alpha )}`$ (13)
$`\times `$ $`L(x_3\overline{\alpha }_3,x^{}\beta ^{};x_3^+\overline{\alpha _3},x^+\alpha ^{})`$ (14)
$``$ $`{\displaystyle \frac{v_0}{4}}(i)^2{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle \underset{\alpha ^{},\beta ^{}}{}}\sigma _{\alpha ,\beta }^\mu \sigma _{\alpha ^{},\beta ^{}}^\nu {\displaystyle d3G(x\beta ,x_3\alpha _3)G(x_3\overline{\alpha _3},x\alpha )}`$ (15)
$`\times `$ $`L(x_3\alpha _3,x^{}\beta ^{};x_3^+\overline{\alpha _3},x^+\alpha ^{})`$ (16)
with $`\overline{\alpha }=\alpha `$. This is apparently not a closed-form equation for $`𝒳`$ itself. By the manipulations reported in Appendix A, however, Eq. (16) can be cast in the form of a coupled set of equations for the matrix components of the correlation function, as follows:
$$𝒳_{\mu ,\nu }(x,x^{})=𝒳_{\mu ,\nu }^{(0)}(x,x^{})+2v_0\underset{\mu ^{},\nu ^{}}{}𝑑x^{\prime \prime }𝒳_{\mu ,\mu ^{}}^{(0)}(x,x^{\prime \prime })ϵ_{\mu ^{},\nu ^{}}𝒳_{\nu ^{},\nu }(x^{\prime \prime },x^{})$$
(17)
where we have introduced the tensor
$$ϵ_{\mu ,\nu }=\left(\begin{array}{cccc}\hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1\end{array}\right).$$
(18)
To solve Eq. (17), the explicit form of the *non-interacting counterpart* $`𝒳^{(0)}`$ of $`𝒳`$ is required. To this end, the ground-state average in Eq. (8) is evaluated as shown in Appendix A, in terms of the eigenvalues ($`ϵ_r`$) and eigenvectors ($`W_{\xi ,r}`$) of the mean-field Hubbard Hamiltonian (see also Section 3), yielding for the time Fourier transform of $`𝒳^{(0)}`$ the expression
$`𝒳_{\mu ,\nu }^{(0)}(𝐫,𝐫^{};\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4𝒩^2}}{\displaystyle \underset{i,j}{}}{\displaystyle \underset{𝐤,𝐤^{}}{\overset{BZ}{}}}\mathrm{e}^{i(𝐤𝐤^{})(𝐑_i𝐑_j)}|\varphi (𝐫𝐑_i)|^2|\varphi (𝐫^{}𝐑_j)|^2`$ (19)
$`\times `$ $`{\displaystyle \underset{r,r^{}}{}}{\displaystyle \underset{\mu ^{},\nu ^{}}{}}T_{\mu ,\mu ^{}}(\mathrm{\Omega }_i)T_{\nu ,\nu ^{}}(\mathrm{\Omega }_j)F_{r^{},r}^\mu ^{}(𝐤^{},𝐤)F_{r,r^{}}^\nu ^{}(𝐤,𝐤^{})`$ (20)
$`\times `$ $`_{r,r^{}}(𝐤,𝐤^{},\omega )`$ (21)
where $`BZ`$ stands for the two-dimensional Brillouin zone, $`𝒩`$ is the number of lattice sites, and the quantities $`T`$, $`F`$, and $``$ are defined in Appendix A ($`\mathrm{\Omega }_i`$ standing for the angles defining the local spin quantization axis). We consider further the space Fourier transform
$`𝒳_{ab}(𝐪,𝐪^{};\omega )`$ $`=`$ $`{\displaystyle \frac{1}{𝒩V_0}}{\displaystyle 𝑑𝐫𝑑𝐫^{}\mathrm{e}^{i𝐪𝐫}𝒳_{ab}(𝐫,𝐫^{};\omega )\mathrm{e}^{i𝐪^{}𝐫^{}}}`$ (22)
$``$ $`{\displaystyle \frac{1}{V_0}}S(𝐪)S^{}(𝐪^{})\widehat{𝒳}_{ab}(𝐪,𝐪^{};\omega )`$ (23)
where $`V_0`$ is the volume of the elementary crystal cell and
$$S(𝐪)𝑑𝐫\mathrm{e}^{i\mathrm{𝐪𝐫}}|\varphi (𝐫)|^2$$
(24)
is a form factor (which can be set equal to unity for all practical purposes). Applying a suitable unitary transformation \[cf. Eq. (A31)\] which renders the matrix $`T(\mathrm{\Omega }_i)`$ of Eq. (21) diagonal, one gets for its lattice Fourier transform:
$`\overline{T}_{ab}(𝐤)`$ $`=`$ $`{\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{i}{}}\mathrm{e}^{i𝐤𝐑_i}\overline{T}_{ab}(\mathrm{\Omega }_i)`$ (25)
$`=`$ $`\left(\begin{array}{cccc}\mathrm{\Delta }(𝐤)& 0& 0& 0\\ 0& \mathrm{\Delta }(𝐤𝐐)& 0& 0\\ 0& 0& \mathrm{\Delta }(𝐤)& 0\\ 0& 0& 0& \mathrm{\Delta }(𝐤+𝐐)\end{array}\right)`$ (30)
$`\mathrm{\Delta }(𝐤)`$ being the lattice Kronecker delta function. In the new basis, we thus obtain for the non-interacting correlation function the expression:
$`\widehat{𝒳}_{ab}^{(0)}(𝐪,𝐪^{};\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4𝒩}}{\displaystyle \underset{\mathrm{𝐤𝐤}^{}}{\overset{BZ}{}}}{\displaystyle \underset{a^{},b^{}}{}}{\displaystyle \underset{r,r^{}}{}}\overline{T}_{aa^{}}(𝐤𝐤^{}𝐪)\overline{T}_{b^{}b}(𝐤^{}𝐤+𝐪^{})`$ (31)
$`\times `$ $`\overline{F}_{r^{},r}^a^{}(𝐤^{},𝐤)\overline{F}_{r,r^{}}^b^{}(𝐤,𝐤^{})_{r,r^{}}(𝐤,𝐤^{},\omega )`$ (32)
where the overbar denotes matrices transformed according to the above unitary transformation \[cf. Eqs. (A27) and (A31)\]. In this way, the integral equation (17) reduces to the form:
$`\widehat{𝒳}_{ab}(𝐪,𝐪^{};\omega )`$ $`=`$ $`\widehat{𝒳}_{ab}^{(0)}(𝐪,𝐪^{};\omega )`$ (33)
$`+`$ $`2U{\displaystyle \underset{a^{}b^{}}{}}{\displaystyle \underset{𝐪^{\prime \prime }}{}}\widehat{𝒳}_{aa^{}}^{(0)}(𝐪,𝐪^{\prime \prime };\omega )\overline{ϵ}_{a^{}b^{}}\widehat{𝒳}_{b^{}b}(𝐪^{\prime \prime },𝐪^{};\omega )`$ (34)
with $`\overline{ϵ}`$ given by Eq. (A30). This equation can be solved by the methods of Appendix A, yielding the closed-form expression:
$`\widehat{𝒳}_{ab}(𝐪+\gamma _a𝐐,𝐪^{};\omega )`$ $`=`$ $`{\displaystyle \underset{a^{}}{}}[\mathrm{𝟏}+X(𝐪,\omega )]_{aa^{}}^1`$ (35)
$`\times `$ $`X_{a^{}b}^{(0)}(𝐪+\gamma _a^{}𝐐;\omega |𝐐)\mathrm{\Delta }(𝐪𝐪^{}\gamma _b𝐐)`$ (36)
where $`\mathrm{𝟏}`$ is the $`4\times 4`$ unit matrix, the matrix $`X(𝐪,\omega )`$ is defined by Eq. (A54), and with the notation $`\gamma _a=0`$ for $`a=0,2`$, $`\gamma _a=1`$ for $`a=1`$, and $`\gamma _a=1`$ for $`a=3`$. Although still expressed in the transformed basis, Eq. (36) is the desired expression for the Fourier transform of the generalized correlation function, which holds within the RPA approximation *for any value of* $`𝐐`$.
The spin-wave *dispersion relation* can eventually be obtained by searching for the zeros of the inverse matrix on the right-hand side of Eq. (36), which is equivalent to imposing the condition:
$$det[\mathrm{𝟏}+X(𝐪,\omega )]=\mathrm{\hspace{0.17em}0}.$$
(37)
It can be verified that the condition (37) can be mapped onto the result reported in Ref. , where the dispersion relation has then been obtained numerically for chosen values of $`𝐐`$. In the present paper we proceed instead to deriving the *analytic* form of the dispersion relation for *small values of the parameter* $`t/U`$ of the Hubbard Hamiltonian.
To this end, it is convenient to rewrite first the matrix $`X`$ in Eq. (37) in a more conventional basis identified by the labels ($`0,+,,z`$), with $`\sigma ^\pm =(\sigma ^x\pm i\sigma ^y)/\sqrt{2}`$. The matrix $`\mathrm{𝟏}+X(𝐪,\omega )`$ is then transformed into:
$$M(𝐪,\omega )=\mathbf{\hspace{0.17em}1}+2U\left(\begin{array}{cccc}𝒳_0^{0,0}(𝐪,\omega )& 𝒳_0^{0,}(𝐪,\omega )& 𝒳_0^{0,+}(𝐪,\omega )& 𝒳_0^{0,z}(𝐪,\omega )\\ 𝒳_0^{+,0}(𝐪,\omega )& 𝒳_0^{+,}(𝐪,\omega )& 𝒳_0^{+,+}(𝐪,\omega )& 𝒳_0^{+,z}(𝐪,\omega )\\ 𝒳_0^{,0}(𝐪,\omega )& 𝒳_0^,(𝐪,\omega )& 𝒳_0^{,+}(𝐪,\omega )& 𝒳_0^{,z}(𝐪,\omega )\\ 𝒳_0^{z,0}(𝐪,\omega )& 𝒳_0^{z,}(𝐪,\omega )& 𝒳_0^{z,+}(𝐪,\omega )& 𝒳_0^{z,z}(𝐪,\omega )\end{array}\right)$$
(38)
where now
$$𝒳_0^{\alpha ,\beta }(𝐪,\omega )=\frac{1}{4𝒩}\underset{r,r^{}}{}\underset{𝐤}{\overset{BZ}{}}F_{r^{},r}^\alpha (𝐤𝐪,𝐤)F_{r,r^{}}^\beta (𝐤,𝐤𝐪)_{r,r^{}}(𝐤,𝐤𝐪,\omega )$$
(39)
and
$$F_{r,r^{}}^\mu (𝐤,𝐤^{})=\underset{\xi ,\xi ^{}}{}W_{r,\xi }^{}(𝐤)\sigma _{\xi ,\xi ^{}}^\mu W_{\xi ^{},r^{}}(𝐤^{})$$
(40)
with $`\sigma ^+=\sqrt{2}\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)`$ and $`\sigma ^{}=\sqrt{2}\left(\begin{array}{cc}\hfill 0& \hfill 0\\ \hfill 1& \hfill 0\end{array}\right)`$. Note that two columns in the expression (38) appear interchanged with respect to the original order, owing to the presence in the final basis of the tensor
$$\stackrel{~}{ϵ}=\left(\begin{array}{cccc}\hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1\end{array}\right)$$
(41)
in the place of $`\overline{ϵ}`$ given by Eq. (A30).
We pass now to perform the small $`t/U`$ expansion of the spin-wave dispersion relation obtained from the condition $`detM(𝐪,\omega )=0`$.
## III Expansion in the small parameter $`t/U`$
In this Section, we obtain explicitly the spin wave dispersion relation to *second order* in the small parameter $`t/U`$, from the general condition $`detM(𝐪,\omega )=0`$ obtained in Section 2 within the RPA approximation for the zero-temperature broken-symmetry phase with a generic incommensurate wave vector $`𝐐`$. To this end, we will preliminary expand the self-consistency parameters of the mean-field Hamiltonian, as well as the equations they satisfy, at the relevant order in $`t/U`$; we will then expand the matrix elements of the matrix $`M(𝐪,\omega )`$ defined by Eq. (38) at the relevant order in $`t/U`$.
A. Mean-field equations
The mean-field equations for a single-band Hubbard Hamiltonian in the presence of an incommensurate spiral spin structure have been discussed in Ref.. Introducing a local set of spin quantization axis, with the $`z`$ axis transformed locally into the axis specified by the spherical angles $`\mathrm{\Omega }_i(\theta _i=𝐐𝐑_i,\phi _i=0)`$ at site $`i`$, one transforms the destruction operators $`c_{i\alpha }`$ according to Eq. (A9) and performs the (Hartree-Fock) mean-field decoupling of the Hubbard Hamiltonian, yielding
$$H(𝐐)=\underset{𝐤}{\overset{BZ}{}}\underset{\xi ,\xi ^{}}{}d_{𝐤\xi }^{}_{\xi ,\xi ^{}}(𝐤)d_{𝐤\xi ^{}}U𝒩\left(m_1^2m_2^2\right)$$
(42)
with $`\xi ,\xi ^{}=(+,)`$ and where $``$ is the $`2\times 2`$ matrix
$$(𝐤)=\left(\begin{array}{cc}ϵ_0\mu +tT_e(𝐤)+U(m_1m_2)& itT_o(𝐤)\\ itT_o(𝐤)& ϵ_0\mu +tT_e(𝐤)+U(m_1+m_2)\end{array}\right).$$
(43)
In the above expressions, $`ϵ_0`$ is the site energy, $`\mu `$ the chemical potential, $`m_1`$ and $`m_2`$ represent the occupation number and magnetization along the local quantization axis, respectively, and $`T_{e/o}(𝐤)`$ read
$$T_e(𝐤)=\mathrm{\hspace{0.17em}2}\left[\mathrm{cos}k_x\mathrm{cos}\left(\frac{Q_x}{2}\right)+\mathrm{cos}k_y\mathrm{cos}\left(\frac{Q_y}{2}\right)\right]$$
(44)
$$T_o(𝐤)=\mathrm{\hspace{0.17em}2}\left[\mathrm{sin}k_x\mathrm{sin}\left(\frac{Q_x}{2}\right)+\mathrm{sin}k_y\mathrm{sin}\left(\frac{Q_y}{2}\right)\right].$$
(45)
The eigenvalues and eigenvectors of the matrix (43) are thus given by:
$$ϵ_r(𝐤)=ϵ_0\mu +tT_e(𝐤)+Um_1+(1)^r\sqrt{U^2m_2^2+t^2T_o(𝐤)^2}$$
(46)
($`r=1,2`$) and
$$W_1(𝐤)=\frac{1}{N_1(𝐤)}\left(\begin{array}{c}1\\ i\frac{\frac{t}{U}\frac{1}{m_2}T_o(𝐤)}{1+\sqrt{1+\left(\frac{t}{U}\right)^2\frac{1}{m_2^2}T_o(𝐤)^2}}\end{array}\right)$$
(47)
$$W_2(𝐤)=\frac{1}{N_2(𝐤)}\left(\begin{array}{c}i\frac{\frac{t}{U}\frac{1}{m_2}T_o(𝐤)}{1+\sqrt{1+\left(\frac{t}{U}\right)^2\frac{1}{m_2^2}T_o(𝐤)^2}}\\ 1\end{array}\right)$$
(48)
where $`N_r(𝐤)`$ stands for the normalization factor.
The parameters $`m_1`$, $`m_2`$, and $`𝐐`$ of the mean-field Hamiltonian (42) are obtained, as usual, by minimizing the average value of the Hamiltonian with respect to the parameters themselves. One obtains:
$$m_1=\frac{1}{2𝒩}\underset{𝐤}{\overset{BZ}{}}\underset{r}{}f_F(ϵ_r(𝐤))=\frac{1+\delta }{2},$$
(49)
where $`f_F(ϵ)`$ is the (zero-temperature) Fermi function and $`\delta `$ is the *doping parameter*,
$$m_2=\frac{1}{2𝒩}\underset{𝐤}{\overset{BZ}{}}\underset{\xi }{}\underset{r}{}\xi W_{r,\xi }^{}(𝐤)W_{\xi ,r}(𝐤)f_F(ϵ_r(𝐤)),$$
(50)
and
$$\underset{𝐤}{\overset{BZ}{}}\underset{\xi ,\xi ^{}}{}\stackrel{}{}_𝐐_{\xi ,\xi ^{}}(𝐤)\underset{r}{}W_{r,\xi }^{}(𝐤)W_{\xi ^{},r}(𝐤)f_F(ϵ_r(𝐤))=\mathrm{\hspace{0.17em}0}.$$
(51)
In the following, we shall restrict to the *diagonal solution* $`𝐐=Q(1,1)`$, since it is known to be favored *for sufficiently small values* of $`t/U`$. Accordingly, at the order we are considering of the small parameter $`t/U`$ we expand formally:
$`ϵ_r(𝐤)`$ $`=`$ $`U[ϵ_r^{(0)}(𝐤)\mu ^{(0)}+\left({\displaystyle \frac{t}{U}}\right)(ϵ_r^{(1)}(𝐤)\mu ^{(1)})`$ (52)
$`+`$ $`\left({\displaystyle \frac{t}{U}}\right)^2(ϵ_r^{(2)}(𝐤)\mu ^{(2)})+\mathrm{}]`$ (53)
as well as
$$m_2=m_2^{(0)}+\left(\frac{t}{U}\right)m_2^{(1)}+\left(\frac{t}{U}\right)^2m_2^{(2)}+\mathrm{},$$
(54)
while, at the relevant order we can take
$$W_1(𝐤)=\frac{1}{N_r(𝐤)}\left(\begin{array}{c}1\\ i\frac{t}{U}\frac{T_o(𝐤)}{2m_2^{(0)}}\end{array}\right)$$
(55)
$$W_2(𝐤)=\frac{1}{N_r(𝐤)}\left(\begin{array}{c}i\frac{t}{U}\frac{T_o(𝐤)}{2m_2^{(0)}}\\ 1\end{array}\right)$$
(56)
where
$$\frac{1}{N_r(𝐤)^2}=1\left(\frac{t}{U}\right)^2\left(\frac{T_o(𝐤)}{2m_2^{(0)}}\right)^2+𝒪\left(\left(\frac{t}{U}\right)^3\right)$$
(57)
is independent from $`r`$. Note that it is sufficient to retain the lowest-order term $`m_2^{(0)}`$ in the above equations. The parameter $`m_1`$, on the other hand, is given by Eq. (49) and is thus formally independent from $`t/U`$ (we anticipate, however, that the doping parameter $`\delta `$ will turn out to be at most of the order $`t/U`$ for our expansions to be internally consistent).
The coefficients $`ϵ_r^{(n)}(𝐤)`$ of Eq. (53) with $`n=0,1,2,\mathrm{}`$ can be readily obtained from the expression (46) (where we may set $`ϵ_0=0`$ for simplicity) in terms of the $`m_2^{(n)}`$ (for given $`𝐐`$). The value of $`\mu ^{(0)}`$ can also be readily obtained in terms of $`ϵ_r^{(0)}`$ (cf. Appendix B). The remaining coefficients of the expansions (53) and (54) can further be determined by solving the coupled equations for the self-consistency parameters according to the methods developed in the Appendices B and C. In particular, for $`\delta >0`$ we obtain:
$`m_2^{(0)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1\delta )`$ (58)
$`m_2^{(1)}`$ $`=`$ $`0`$ (59)
$`m_2^{(2)}`$ $`=`$ $`4\mathrm{sin}^2(Q/2)+𝒪(\delta ),`$ (60)
$`\mu ^{(0)}m_2^{(0)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1+\delta )`$ (61)
$`\mu ^{(1)}m_2^{(1)}`$ $`=`$ $`\mathrm{cos}(Q/2)[44\pi \delta +𝒪(\delta ^2)]`$ (62)
$`\mu ^{(2)}m_2^{(2)}`$ $`=`$ $`0+𝒪(\delta ),`$ (63)
and
$`ϵ_r^{(0)}(𝐤)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(1+\delta )+(1)^r(1\delta )\right]`$ (64)
$`ϵ_r^{(1)}(𝐤)`$ $`=`$ $`T_e(𝐤)`$ (65)
$`ϵ_r^{(2)}(𝐤)`$ $`=`$ $`(1)^r\left[T_o^2(𝐤)4\mathrm{sin}^2(Q/2)\right]+𝒪(\delta ),`$ (66)
where for the spiral configuration we are considering $`Q`$ is determined by the condition
$$\mathrm{cos}(Q/2)=\frac{U\delta }{2t}+𝒪\left(t/U\right)$$
(67)
with $`\delta 2(t/U)+𝒪((t/U)^2)`$, as anticipated. (The other allowed solution $`\mathrm{sin}(Q/2)=0`$ describes instead the ferromagnetic case.) As discussed in Appendix B, the above results have been obtained with the further assumption that $`\delta `$ is small enough but not infinitesimal, i.e., $`\delta `$ satisfies the condition $`(t/U)^2<<\delta `$. The case $`\delta =0`$, on the other hand, can be considered separately. Note that Eq. (67) implies that in the spiral phase $`\delta `$ is at most of the order $`t/U`$. This property has to be taken into account to get a consistent expansion up to the desired order in $`t/U`$.
B. Susceptibilities and spin-wave dispersion
Before performing the $`t/U`$ expansion of the matrix elements of the non-interacting susceptibility tensor (39) to get the spin-wave dispersion relation, it is convenient to exploit some symmetry properties that reduce the number of matrix elements to be considered. Specifically, from the property
$$_{r,r^{}}(𝐤,𝐤𝐪,\omega )=_{r^{},r}(𝐤𝐪,𝐤,\omega )$$
(68)
it follows that:
$$𝒳_0^{\alpha ,\beta }(𝐪,\omega )=𝒳_0^{\beta ,\alpha }(𝐪,\omega ).$$
(69)
By direct inspection it can also be verified that:
$`𝒳_0^{0,}(𝐪,\omega )`$ $`=`$ $`𝒳_0^{+,0}(𝐪,\omega ),`$ (70)
$`𝒳_0^{z,}(𝐪,\omega )`$ $`=`$ $`𝒳_0^{+,z}(𝐪,\omega ),`$ (71)
$`𝒳_0^{+,+}(𝐪,\omega )`$ $`=`$ $`𝒳_0^,(𝐪,\omega ).`$ (72)
In this way, the matrix (38) acquires the simplified form:
$$M(𝐪,\omega )=\left(\begin{array}{cccc}1a(𝐪,\omega )& ib(𝐪,\omega )& ib(𝐪,\omega )& c(𝐪,\omega )\\ ib(𝐪,\omega )& 1+d(𝐪,\omega )& e(𝐪,\omega )& if(𝐪,\omega )\\ ib(𝐪,\omega )& e(𝐪,\omega )& 1+d(𝐪,\omega )& if(𝐪,\omega )\\ c(𝐪,\omega )& if(𝐪,\omega )& if(𝐪,\omega )& 1+g(𝐪,\omega )\end{array}\right)$$
(73)
where we have set
$`a(𝐪,\omega )`$ $`=`$ $`2U𝒳_0^{0,0}(𝐪,\omega )`$ (74)
$`b(𝐪,\omega )`$ $`=`$ $`2iU𝒳_0^{0,}(𝐪,\omega )`$ (75)
$`c(𝐪,\omega )`$ $`=`$ $`2U𝒳_0^{0,z}(𝐪,\omega )`$ (76)
$`d(𝐪,\omega )`$ $`=`$ $`2U𝒳_0^{+,}(𝐪,\omega )`$ (77)
$`e(𝐪,\omega )`$ $`=`$ $`2U𝒳_0^{+,+}(𝐪,\omega )`$ (78)
$`f(𝐪,\omega )`$ $`=`$ $`2iU𝒳_0^{+,z}(𝐪,\omega )`$ (79)
$`g(𝐪,\omega )`$ $`=`$ $`2U𝒳_0^{z,z}(𝐪,\omega ).`$ (80)
Entering then the expansions (55)-(57) for the eigenvectors of the mean-field Hamiltonian (with $`m_2^{(0)}`$ given by Eq. (60)) into the definition (40), we obtain the expressions for the relevant matrix elements of the non-interacting susceptibility tensor (39) reported in Appendix B at the order in $`t/U`$ we are considering. Utilizing further the method developed in Appendix C to perform the $`𝐤`$ summation when the doping parameter $`\delta `$ is small, we obtain eventually the following expressions for the matrix elements (77):
$`a(𝐪,\omega )`$ $`=`$ $`a(𝐪)={\displaystyle \frac{1}{2\mathrm{cos}(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2)}}{\displaystyle \frac{\delta U}{t}}+𝒪\left(t/U\right),`$ (81)
$`b(𝐪,\omega )`$ $`=`$ $`b(𝐪)={\displaystyle \frac{\mathrm{sin}(Q/2)(\mathrm{sin}q_x+\mathrm{sin}q_y)}{\sqrt{2}\mathrm{cos}(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2)}}\delta +𝒪\left(\left(t/U\right)^2\right),`$ (82)
$`c(𝐪,\omega )`$ $`=`$ $`c(𝐪)=a(𝐪,\omega )+𝒪\left(t/U\right),`$ (83)
$`e(𝐪,\omega )`$ $`=`$ $`e(𝐪)=4\left({\displaystyle \frac{t}{U}}\right)^2\mathrm{sin}^2(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y)+𝒪\left(\left(t/U\right)^3\right),`$ (84)
$`f(𝐪,\omega )`$ $`=`$ $`f(𝐪)=b(𝐪,\omega )+𝒪\left(\left(t/U\right)^2\right),`$ (85)
$`g(𝐪,\omega )`$ $`=`$ $`g(𝐪)=a(𝐪,\omega )+𝒪\left(t/U\right),`$ (86)
as well as
$`d(𝐪,\omega )`$ $``$ $`1\stackrel{~}{\omega }+\alpha (𝐪)+𝒪\left(\left(t/U\right)^3\right)`$ (87)
$`=`$ $`1\stackrel{~}{\omega }+\left({\displaystyle \frac{t}{U}}\right)^2\mathrm{\hspace{0.17em}4}\mathrm{cos}^2(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2)+8\left({\displaystyle \frac{t}{U}}\right)^2\mathrm{sin}^2(Q/2)`$ (88)
$`+`$ $`2\left({\displaystyle \frac{t}{U}}\right)\delta \mathrm{cos}(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2)`$ (89)
$``$ $`2\left({\displaystyle \frac{t}{U}}\right)\delta {\displaystyle \frac{\mathrm{sin}^2(Q/2)(\mathrm{sin}q_x+\mathrm{sin}q_y)^2}{\mathrm{cos}(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2)}}+𝒪\left(\left(t/U\right)^3\right)`$ (90)
where $`\stackrel{~}{\omega }\omega /U`$ will turn out to be of order $`t^2/U^2`$ at the spin-wave poles.
To obtain the above expressions, we have considered only the real part of the functions (A24). This is definitely possible for every pairs of bands when we restrict to values of $`𝐪`$ such that $`|𝐪|k_F`$, where $`k_F`$ (by our definition) is the maximum value of the function $`k(\varphi )`$ introduced in Appendix C, which coincides with the Fermi momentum for small $`\delta `$. Since we have shown in the same Appendix that $`k(\varphi )\sqrt{\delta }\sqrt{t/U}`$, taking the $`𝐪0`$ limit implies letting $`t/U`$ to vanish before $`𝐪`$.
Note also that, although the expressions (86) and (90) have been calculated at different orders in $`t/U`$, the ensuing expression for the determinant of the matrix (73) is obtained consistently at the fourth order in $`t/U`$, as required for the frequency of the spin-wave mode to be of the order of $`t^2/U`$. In fact, by writing the determinant explicitly we obtain:
$`\left[\left(1a(𝐪)\right)\left(1+a(𝐪)\right)+a^2(𝐪)+𝒪\left(t/U\right)\right]`$ (91)
$`\times \left[\left(\stackrel{~}{\omega }+\alpha (𝐪)\right)\left(\stackrel{~}{\omega }+\alpha (𝐪)\right)e^2(𝐪)+𝒪\left(\left(t/U\right)^5\right)\right]`$ (92)
$`+2\left[1a(𝐪)+𝒪\left(t/U\right)\right]\left[b^2(𝐪)+𝒪\left(\left(t/U\right)^3\right)\right]\left[e(𝐪)\alpha (𝐪)+𝒪\left(\left(t/U\right)^3\right)\right]`$ (93)
$`2\left[1+a(𝐪)+𝒪\left(t/U\right)\right]\left[b^2(𝐪)+𝒪\left(\left(t/U\right)^3\right)\right]\left[e(𝐪)\alpha (𝐪)+𝒪\left(\left(t/U\right)^3\right)\right]`$ (94)
$`+\mathrm{\hspace{0.17em}4}\left[a(𝐪)+𝒪\left(t/U\right)\right]\left[b^2(𝐪)+𝒪\left(\left(t/U\right)^3\right)\right]\left[e(𝐪)\alpha (𝐪)+𝒪\left(\left(t/U\right)^3\right)\right]=\mathrm{\hspace{0.17em}0}.`$ (95)
(96)
Note that the last three terms on the left-hand side add up to zero at the fourth order in $`t/U`$ we are considering. We are thus left with the expression
$$\stackrel{~}{\omega }^2\alpha ^2(𝐪)+e^2(𝐪)+𝒪\left(\left(t/U\right)^5\right)=\mathrm{\hspace{0.17em}0},$$
(97)
yielding
$$\omega ^2(𝐪)=U^2\left(\alpha ^2(𝐪)e^2(𝐪)\right)+𝒪\left(t^5/U^3\right).$$
(98)
Taking into account the expressions for $`\alpha (𝐪)`$ and $`e(𝐪)`$ given above, we get eventually the desired spin-wave dispersion relation, in the form:
$`\omega ^2(𝐪)`$ $`=`$ $`16{\displaystyle \frac{t^4}{U^2}}\mathrm{cos}Q(\mathrm{cos}q_x+\mathrm{cos}q_y2)\left[\mathrm{cos}q_x+\mathrm{cos}q_y2\mathrm{cos}Q\right]`$ (99)
$`+`$ $`16{\displaystyle \frac{t^3}{U}}\delta \left[\mathrm{cos}(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2){\displaystyle \frac{\mathrm{sin}^2(Q/2)(\mathrm{sin}q_x+\mathrm{sin}q_y)^2}{\mathrm{cos}(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2)}}\right]`$ (100)
$`\times `$ $`\left[\mathrm{cos}^2(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2)+2\mathrm{sin}^2(Q/2)\right]`$ (101)
$`+`$ $`4\delta ^2t^2\left[\mathrm{cos}(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2){\displaystyle \frac{\mathrm{sin}^2(Q/2)(\mathrm{sin}q_x+\mathrm{sin}q_y)^2}{\mathrm{cos}(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2)}}\right]^2.`$ (102)
To obtain the physical dispersion relation, however, there remains to enter in the above expression the relation among $`Q`$, $`\delta `$, and $`t/U`$ as given by the mean-field condition (67), as discussed in the next Section.
## IV Results and discussion
In this Section, we discuss the physical consequences of the dispersion relation (LABEL:final-dispersion) for spin-wave excitations over an incommensurate (diagonal) spiral magnetic configuration of a two-dimensional Hubbard Hamiltonian.
We have already remarked that the dispersion relation (LABEL:final-dispersion) is not yet in its final form, since the connection among $`Q`$, $`\delta `$, and $`t/U`$ needs still to be specified. Before considering the general case, however, it is interesting to recover from Eq. (LABEL:final-dispersion) the spin-wave dispersion relations corresponding to the limiting cases of an antiferromagnet and of a ferromagnet.
When $`\delta =0`$, Eq. (B32) yields $`\mathrm{cos}(Q/2)=0`$, that is $`𝐐=𝐐_{AF}=(\pi ,\pi )`$. In this case we obtain from Eq. (LABEL:final-dispersion) (by setting $`\delta =0`$ identically therein):
$$\omega (𝐪)=J_{\mathrm{e}ff}^{(AF)}\left[4\left(\mathrm{cos}q_x+\mathrm{cos}q_y\right)^2\right]^{\frac{1}{2}}$$
(104)
with $`J_{\mathrm{e}ff}^{(AF)}=4t^2/U`$. This result coincides with the spin-wave dispersion relation of a two-dimensional Heisenberg *antiferromagnet* at leading order in $`t/U`$.
When $`𝐐=(2\pi ,2\pi )`$ and $`\delta `$ arbitrary, we obtain instead from Eq. (LABEL:final-dispersion)
$$\omega (𝐪)=\frac{4t^2}{U}\left|\left(1\frac{\delta U}{2t}\right)\right|(2\mathrm{cos}q_x\mathrm{cos}q_y),$$
(105)
which coincides with the spin-wave dispersion relation of a nearest-neighbor Heisenberg *ferromagnet* with $`J_{\mathrm{e}ff}^{(F)}=4t^2/U\left[1\delta U/(2t)\right]`$. Note that $`J_{\mathrm{e}ff}^{(F)}`$ is negative since the (mean-field) ferromagnetic solution is actually stable when $`\delta 2t/U`$.
In the general case of a (diagonal) spiral spin configuration, the relation among $`Q`$, $`\delta `$, and $`t/U`$ is given by Eq. (B32), that is, $`\delta =(2t/U)\mathrm{cos}Q/2+𝒪((t/U)^2)`$ at the leading order in $`t/U`$. Eliminating $`Q`$ in favor of $`\delta `$ and $`t/U`$ via this relation in Eq. (LABEL:final-dispersion), we obtain eventually the following expression:
$$\omega (𝐪)=J_{\mathrm{e}ff}\left\{\left[2\frac{(\mathrm{sin}q_x+\mathrm{sin}q_y)^2}{2\mathrm{cos}q_x\mathrm{cos}q_y}\right]^2(\mathrm{cos}q_x+\mathrm{cos}q_y)^2\right\}^{\frac{1}{2}}$$
(106)
where now the effective exchange integral is given by
$$J_{\mathrm{e}ff}=\frac{4t^2}{U}\mathrm{sin}^2(Q/2)=\frac{4t^2}{U}\left(1\left(\frac{U}{2t}\delta \right)^2\right)=\frac{4t^2}{U}\left(1\left(\frac{\delta }{\delta _c}\right)^2\right)$$
(107)
with $`\delta _c2t/U`$. Equation (106) constitutes the main result of this paper. Note that when $`\delta `$ reaches the *critical* value $`\delta _c`$, $`J_{\mathrm{e}ff}`$ vanishes. Past this value, the spiral solution evolves into the ferromagnetic solution, which becomes the stable solution (at the mean-field level).
The real and imaginary parts of $`\omega (𝐪)`$ (in units of $`J_{\mathrm{e}ff}`$), as obtained from the analytic expression (106), are plotted in Figs. 1(a) and 1(b), respectively, over the two-dimensional Brillouin zone.
It is clear from the expression (106) that $`\omega (𝐪)`$ is *either* real *or* purely imaginary, so that, when its real part is nonvanishing, its imaginary part vanishes identically, and viceversa. Actually, this is true in the region of the Brillouin zone where our expansion holds, namely, for $`|𝐪|k_F`$ \[cf. the discussion below Eq. (90)\]. For $`|𝐪|<k_F`$, the spin-wave dispersion acquires a damping due to the mixing with the particle-hole continuum. An exception is represented by the line $`q_x=q_y`$, along which the real and imaginary parts vanish simultaneously. The softening of the dispersion relation along the whole line $`q_x=q_y`$ and not only when $`𝐪=0`$ and $`𝐪=𝐐`$ (as one would expect on general grounds in the presence of a spiral spin configuration ) corresponds to the fact that $`\omega (𝐪)`$ given by Eq. (106) can be cast in the form $`\omega (𝐪)=f(𝐐)g_{\widehat{Q}}(𝐪)`$; as $`\omega (𝐪=𝐐)=0`$ invariably, it follows that $`g_{\widehat{Q}}(𝐪=𝐐)=0`$ since $`g_{\widehat{Q}}(𝐪)`$ does not depend on $`|𝐐|`$. This implies that $`g_{\widehat{Q}}(q_x=q_y)0`$ identically.
One expects this result to be modified, however, at higher order in $`t/U`$. In this respect, it is interesting to compare with the results obtained by Brenig by solving numerically the condition (37) directly, without performing the expansion in the small parameter $`t/U`$. One sees, in particular, from Fig. 6 of Brenig’s paper (obtained for $`\delta =0.075`$) that $`\omega (𝐪)`$ remains indeed finite along $`q_x=q_y`$ already when $`t/U=0.1`$.
Figure 2 shows the region of the Brillouin zone (shaded area) where $`\omega (𝐪)`$ is *overdamped* (purely imaginary) (with the region about $`𝐪=0`$ excluded according to the argument given in Section 3). A finite region of the Brillouin zone where the dispersion relation is overdamped is also reported in Ref. , even though a direct comparison with our results is not possible owing to the different ranges of the parameter $`t/U`$ explored. This overdamping signals an instability of the system (due to the merging into the particle-hole continuum) toward a different ground state, reflecting possibly a more complicated long-range spin (and charge) structure than the spin-spiral one considered in the present paper. Nonetheless, we expect on physical grounds that close to the boundary of the Brillouin zone (where overdamping of spin waves does not occur in our calculation) the spin-wave spectrum with small wavelength obtained by our approach would survive the inclusion of more complicated long-range spin structures.
It is also interesting to compare the form of the dispersion relation (106) with the spin-wave dispersion relation obtained with the Heisenberg model including second and third nearest-neighbor couplings, for the same value of the incommensurability wave vector $`Q`$ (cf. Appendix D). This comparison is shown in Fig. 3 along the symmetry lines of the Brillouin zone and evidences marked differences between the two dispersion relations.
Returning to the dispersion relation (106) for $`\delta 0`$, we emphasize that its functional form could not be obtained from the dispersion relation (104) valid when $`\delta =0`$, by simply modifying the numerical values of the exchange integral therein. Nor, it would be sufficient to include a *finite number* of exchange integrals in the Heisenberg model to account for the finite doping, owing to the presence of an RKKY-type term in Eq. (106) which contains trigonometric functions in the denominator rather than in the numerator only (cf. Appendix E). Recall, in fact, that for practical purposes the Heisenberg model can be regarded as a *fitting model* which could, in principle, reproduce *any* type of spin-wave dispersion relation, provided a sufficiently large number of terms associated with distant neighbors were included. In addition, it appears fair to say that it would have been certainly difficult to guess *a priori* the functional form of the dispersion relation (106), by fitting the numerical results for the itinerant model onto the Heisenberg model extended to a large number of neighbors.
It is also interesting to comment on the effective exchange integral (107) retaining the $`t^2/U`$ dependence of the nearest-neighbor Heisenberg model, even when couplings between far apart neighbors are considered. From a perturbative point of view, when $`U>>t`$ the magnetic interaction between a given lattice site and far apart neighbors is provided by the *mobility* of the holes in the doped configurations. In this respect, one should consider all possible configurations with empty sites distributed at random over the lattice sites, the mobility of the holes then resulting by diagonalizing the Hamiltonian in this basis. As a consequence, the magnetic exchange coupling turns out to be proportional to $`t^2/U`$ at leading order, even for coupling between sites at arbitrary distances.
As anticipated in the Introduction, the form (106) is somewhat hybrid between the one obtained with a nearest-neighbor Heisenberg model (cf. Appendix D) and the long-range RKKY interaction mediated by the conduction electrons (cf. Appendix E). These two contributions to Eq. (106) cannot be separated in a clear and unambiguous way. However, it is possible to trace their origin by considering the expressions (B35) and (LABEL:plus-plus) for the relevant matrix elements of the susceptibility.
If one could set “by hand” $`_{22}=0`$ in those expressions, thus keeping only the interband terms, one would in fact obtain for the dispersion relation:
$$\omega (𝐪)=J_{\mathrm{e}ff}\left[4(\mathrm{cos}q_x+\mathrm{cos}q_y)^2\right]^{\frac{1}{2}},$$
(108)
which corresponds to a nearest-neighbor Heisenberg antiferromagnet with exchange coupling given by Eq. (107). Similarly, if one could set $`_{21}=_{12}=0`$ “by hand”, thus keeping only the intraband terms, one would instead obtain an expression of the RKKY-type:
$$\omega (𝐪)=J_{\mathrm{e}ff}\frac{(\mathrm{sin}q_x+\mathrm{sin}q_y)^2}{2\mathrm{cos}q_x\mathrm{cos}q_y}$$
(109)
with $`J_{\mathrm{e}ff}`$ still given by Eq. (107). Note that in both cases the spectrum would be real and no overdamping would occur. Our general result (106) can then be cast in the following appealing form:
$$\omega (𝐪)=\left\{\left[J(𝐪=0)\omega (𝐪)_{RKKY}\right]^2J(𝐪)^2\right\}^{1/2}$$
(110)
where $`J(𝐪)=J_{\mathrm{e}ff}(\mathrm{cos}q_x+\mathrm{cos}q_y)`$ and $`\omega (𝐪)_{RKKY}`$ is given by the expression (109). Note, in particular, that if one could set $`\omega (𝐪)_{RKKY}=0`$, the Heisenberg form (108) would result from Eq. (110).
A final comment on the possible comparison between the spin-wave spectrum we have obtained and the available experimental data is in order. That the classical spin-wave theory (possibly with quantum corrections ) can accurately describe the spin-wave spectrum over the whole Brillouin zone for the parent compounds of high-temperature superconductors (zero doping) is well established at this point. That when carriers are added (finite doping) incommensurate spin fluctuations occur with rather short coherence length has also been well established in several materials. It is interesting to point out, in addition, that the observation of well-defined spin waves only close to the boundary of the Brillouin zone has been reported, consistently with what we have obtained by our calculation. However, detailed comparison with our form (106) of the spin-wave dispersion relation with the experimental data may not be possible, because our result is valid in the asymptotic limit $`U>>W`$, where $`W8t`$ is the bandwidth, which may not be realized for real materials.
## V Concluding remarks
In this paper, we have studied the spin-wave spectrum for a two-dimensional Hubbard Hamiltonian in the presence of an incommensurate spin-spiral phase within the electronic RPA approximation in the broken-symmetry phase. We have, in particular, obtained the analytic form of this spectrum at the leading order in the small parameter $`t/U`$ of the Hubbard Hamiltonian. Specifically, it has been possible to obtain the spectrum in a *closed form* even in the presence of an *incommensurate* structure, owing to the peculiar symmetry which is intrinsic to the spiral phase. In this respect, starting from an alternative mean-field configuration different from the spiral one (such as, for instance, a “stripe” structure) would have not enabled us to solve for the spin wave spectrum in a *closed form*. By our approach, we have thus pushed the analytic results for the spin wave spectrum of an incommensurate structure as far as possible (apart, obviously, from including higher orders in the $`t/U`$ expansion which could still be done by our approach).
Alternatively, the spin-wave spectrum could have been obtained numerically for any value of $`t/U`$ without performing the small $`t/U`$ expansion, as already reported in Ref. . In this way, however, it would have been rather difficult (if not impossible) to arrive at the functional form (106) for the spin-wave dispersion relation in the strong-coupling limit, which has an hybrid form between the dispersion relation for a nearest-neighbor Heisenberg model and that obtained within the (long-range) RKKY interaction in the presence of a finite doping. Owing to the itinerant character of the system, it is, in fact, the presence of a band of metallic character (crossed by the Fermi level) which generates the RKKY magnetic interaction between the localized spins associated with the (filled) lower band. This novel feature constitutes the main result of the present paper. We have concluded accordingly that, even for small doping, the itinerant model we have considered results in a dispersion relation $`\omega (𝐪)`$ that cannot be effectively represented by the Heisenberg model, for which the couplings extend to a few neighbors only. This occurs because the free carriers, associated with the itinerant character of our starting Hamiltonian, necessarily introduce long-range RKKY-type magnetic interactions among the localized spins.
A serious concern, which is related to the spin-wave spectrum we have obtained, regards the instability occurring about the center of the Brillouin zone, where the spin-wave spectrum becomes purely imaginary. To overcome this point, one should possibly start from a more complicated incommensurate mean-field solution other than the spiral configuration, with a lower ground-state energy. In this way, however, one would unavoidably not obtain a closed-form equation for the spin-wave spectrum, since the incommensurability could in general prevent it. Nonetheless, we expect on physical grounds the spin-wave spectrum we have obtained with the spiral pattern to survive inclusion of more realistic spin structures, if one considers only the region close to the boundary of the Brillouin zone, for which knowledge of the detailed form of the underlying (long-range) spin pattern appears to be less crucial.
## Acknowledgments
We are indebted to P. Salvi for help during the initial stage of this work, especially in developing part of the material contained in Appendix C. We are also indebted to N. Majlis, W. Brenig, and G. Aeppli for discussions.
## A Solution of the RPA equations for the dynamical susceptibilities with an incommensurate spin-spiral ground state
In this Appendix, we provide the details of the derivation of the explicit form of the generalized correlation function within RPA and of the associated spin-wave dispersion relation, which were reported only schematically in Section 2.
We begin by showing how Eq. (16) of the text can be reduced to a closed-form equation for the correlation function $`𝒳`$ itself. To this end, we note from Eq. (5) that the following relations hold for $`\mu =0`$ and $`\mu =z`$, in the order:
$`𝒳_{0,\nu }(x,x^{})`$ $`=`$ $`{\displaystyle \frac{i}{4}}{\displaystyle \underset{\alpha ^{},\beta ^{}}{}}\sigma _{\alpha ^{},\beta ^{}}^\nu L(x+,x^{}\beta ^{};x^++,x^+\alpha ^{})`$ (A1)
$``$ $`{\displaystyle \frac{i}{4}}{\displaystyle \underset{\alpha ^{},\beta ^{}}{}}\sigma _{\alpha ^{},\beta ^{}}^\nu L(x,x^{}\beta ^{};x^+,x^+\alpha ^{})`$ (A2)
and
$`𝒳_{z,\nu }(x,x^{})`$ $`=`$ $`{\displaystyle \frac{i}{4}}{\displaystyle \underset{\alpha ^{},\beta ^{}}{}}\sigma _{\alpha ^{},\beta ^{}}^\nu L(x+,x^{}\beta ^{};x^++,x^+\alpha ^{})`$ (A3)
$`+`$ $`{\displaystyle \frac{i}{4}}{\displaystyle \underset{\alpha ^{},\beta ^{}}{}}\sigma _{\alpha ^{},\beta ^{}}^\nu L(x,x^{}\beta ^{};x^+,x^+\alpha ^{}).`$ (A4)
By adding and subtracting both sides of the above equations, we obtain:
$$\{\begin{array}{cc}i_{\alpha ^{},\beta ^{}}\sigma _{\alpha ^{},\beta ^{}}^\nu L(x+,x^{}\beta ^{};x^++,x^+\alpha ^{})=\hfill & 2\left[𝒳_{0,\nu }(x,x^{})+𝒳_{z,\nu }(x,x^{})\right]\hfill \\ i_{\alpha ^{},\beta ^{}}\sigma _{\alpha ^{},\beta ^{}}^\nu L(x,x^{}\beta ^{};x^+,x^+\alpha ^{})=\hfill & 2\left[𝒳_{0,\nu }(x,x^{})𝒳_{z,\nu }(x,x^{})\right]\hfill \end{array}.$$
(A5)
For $`\mu =x`$ and $`\mu =y`$ we obtain instead:
$$\{\begin{array}{cc}i_{\alpha ^{},\beta ^{}}\sigma _{\alpha ^{},\beta ^{}}^\nu L(x,x^{}\beta ^{};x^++,x^+\alpha ^{})=\hfill & 2\left[𝒳_{x,\nu }(x,x^{})+i𝒳_{y,\nu }(x,x^{})\right]\hfill \\ i_{\alpha ^{},\beta ^{}}\sigma _{\alpha ^{},\beta ^{}}^\nu L(x+,x^{}\beta ^{};x^+,x^+\alpha ^{})=\hfill & 2\left[𝒳_{x,\nu }(x,x^{})i𝒳_{y,\nu }(x,x^{})\right]\hfill \end{array},$$
(A6)
with similar relations holding for the non-interacting counterparts of $`L`$ and $`𝒳`$. Entering these relations into Eq. (16), the closed-form expression (17) results eventually after suitable manipulation.
Next, we specify the form of the non-interacting counterpart $`𝒳^{(0)}`$ of $`𝒳`$, as defined by the first term on the right-hand side of Eq. (16), namely:
$`𝒳_{\mu ,\nu }^{(0)}(x,x^{})`$ $``$ $`{\displaystyle \frac{i}{4}}{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle \underset{\alpha ^{},\beta ^{}}{}}\sigma _{\alpha ,\beta }^\mu \sigma _{\alpha ^{},\beta ^{}}^\nu G(x\beta ,x^{}\alpha ^{})G(x^{}\beta ^{},x\alpha )`$ (A7)
$`=`$ $`{\displaystyle \frac{i}{4}}{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle \underset{\alpha ^{},\beta ^{}}{}}\sigma _{\alpha ,\beta }^\mu \sigma _{\alpha ^{},\beta ^{}}^\nu T\left[\psi _\beta (x)\psi _\alpha ^{}^{}(x^{})\right]T\left[\psi _\beta ^{}(x^{})\psi _\alpha ^{}(x)\right].`$ (A8)
To this end, we follow Ref. and introduce the set of destruction operators $`d_{i\xi }`$ along the local spin-quantization axes, which identify the spiral pattern of the magnetic ground state within the mean-field approximation. The field operator (3) acquires then the form:
$`\psi _\alpha (𝐫)`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \underset{\xi }{}}\varphi (𝐫𝐑_i)(\mathrm{\Omega }_i)_{\alpha \xi }d_{i\xi }`$ (A9)
$`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \underset{\xi }{}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\varphi (𝐫𝐑_i)(\mathrm{\Omega }_i)_{\alpha \xi }{\displaystyle \frac{\mathrm{e}^{i𝐤𝐑_i}}{\sqrt{𝒩}}}d_{𝐤\xi }`$ (A10)
where $`(\mathrm{\Omega }_i)`$ is the spin rotation operator associated with the angles $`\mathrm{\Omega }_i`$ and where the Bloch transformation has been introduced as in Section 3, which brings the mean-field Hubbard Hamiltonian into block form for each wave vector $`𝐤`$ belonging to the Brillouin zone ($`BZ`$).
Let $`W_{\xi ,r}(𝐤)`$ (with $`\xi =+,`$) be the matrix which diagonalizes the mean-field Hubbard Hamiltonian at given $`𝐤`$, such that \[cf. Eq. (43)\]
$$\underset{𝐤}{\overset{BZ}{}}\underset{\xi ,\xi ^{}}{}d_{𝐤\xi }^{}_{\xi ,\xi ^{}}(𝐤)d_{𝐤\xi ^{}}=\underset{𝐤}{\overset{BZ}{}}\underset{r}{}\gamma _{𝐤,r}^{}ϵ_r(𝐤)\gamma _{𝐤,r}$$
(A11)
with
$$\gamma _{𝐤,r}=\underset{\xi }{}W_{r,\xi }^{}(𝐤)d_{𝐤\xi }$$
(A12)
and $`ϵ_r(𝐤)`$ given by Eq. (46). Upon averaging over the the broken-symmetry ground state, one then obtains:
$$T[d_{𝐤\xi }(t)d_{𝐤^{}\xi ^{}}^{}(t^{})]=\underset{r,r^{}}{}W_{\xi ,r}(𝐤)W_{r^{},\xi ^{}}^{}(𝐤^{})T[\gamma _{𝐤r}(t)\gamma _{𝐤^{}r^{}}^{}(t^{})]$$
(A13)
where
$`T[\gamma _{𝐤r}(t)\gamma _{𝐤^{}r^{}}^{}(t^{})]`$ (A14)
$`=\delta _{𝐤,𝐤^{}}\delta _{r,r^{}}e^{iϵ_r(𝐤)(tt^{})}\left\{\mathrm{\Theta }(tt^{})\left[1f_F(ϵ_r(𝐤))\right]\mathrm{\Theta }(t^{}t)f_F(ϵ_r(𝐤))\right\}`$ (A15)
$`\mathrm{\Theta }(t)`$ being the unit step function. Introducing further the tensor $`T(\mathrm{\Omega }_i)`$ via the relation
$$\underset{\alpha ,\beta }{}^{}(\mathrm{\Omega }_i)_{\xi \alpha }\sigma _{\alpha ,\beta }^\mu (\mathrm{\Omega }_i)_{\beta ,\xi ^{}}=\underset{\nu }{}T_{\mu \nu }(\mathrm{\Omega }_i)\sigma _{\xi ,\xi ^{}}^\nu ,$$
(A16)
such that
$$T(\mathrm{\Omega }_i)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& \mathrm{cos}\theta _i& 0& \mathrm{sin}\theta _i\\ 0& 0& 1& 0\\ 0& \mathrm{sin}\theta _i& 0& \mathrm{cos}\theta _i\end{array}\right)$$
(A17)
where $`\theta _i=𝐐𝐑_i`$ within the spiral-spin pattern we are considering, and approximating
$$\varphi ^{}(𝐫𝐑_i)\varphi (𝐫𝐑_j)|\varphi (𝐫𝐑_i)|^2\delta _{i,j}$$
(A18)
for localized atomic (Wannier) orbitals, we obtain the following expression for the frequency Fourier transform of the non-interacting correlation function:
$`𝒳_{\mu ,\nu }^{(0)}(𝐫,𝐫^{};\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4𝒩^2}}{\displaystyle \underset{i,j}{}}{\displaystyle \underset{𝐤,𝐤^{}}{\overset{BZ}{}}}e^{i(𝐤𝐤^{})(𝐑_i𝐑_j)}|\varphi (𝐫𝐑_i)|^2|\varphi (𝐫^{}𝐑_j)|^2`$ (A19)
$`\times `$ $`{\displaystyle \underset{r,r^{}}{}}{\displaystyle \underset{\mu ^{},\nu ^{}}{}}T_{\mu ,\mu ^{}}(\mathrm{\Omega }_i)T_{\nu ,\nu ^{}}(\mathrm{\Omega }_j)F_{r^{},r}^\mu ^{}(𝐤^{},𝐤)F_{r,r^{}}^\nu ^{}(𝐤,𝐤^{})`$ (A20)
$`\times `$ $`_{r,r^{}}(𝐤,𝐤^{},\omega ).`$ (A21)
Here we have introduced the notation:
$$F_{r,r^{}}^\mu (𝐤,𝐤^{})\underset{\xi ,\xi ^{}}{}W_{r,\xi }^{}(𝐤)\sigma _{\xi ,\xi ^{}}^\mu W_{\xi ^{},r^{}}(𝐤^{})$$
(A22)
as well as
$`_{r,r^{}}(𝐤,𝐤^{},\omega )`$ $`=`$ $`{\displaystyle \frac{\left[1f_F(ϵ_r(𝐤))\right]f_F(ϵ_r^{}(𝐤^{}))}{\omega ϵ_r(𝐤)+ϵ_r^{}(𝐤^{})+i\eta }}`$ (A23)
$``$ $`{\displaystyle \frac{\left[1f_F(ϵ_r^{}(𝐤^{}))\right]f_F(ϵ_r(𝐤))}{\omega ϵ_r(𝐤)+ϵ_r^{}(𝐤^{})i\eta }}`$ (A24)
where $`f_F(ϵ)`$ is the Fermi function. Equation (21) of the text is thus recovered.
We have mentioned in Section 2 that the solution to the integral equation (17) could be considerably simplified, provided the matrix $`T(\mathrm{\Omega }_i)`$ given by Eq. (A17) were preliminary brought to diagonal form by a suitable unitary transformation. This transformation reads:
$$\overline{\sigma }^a=\underset{\mu =(0,x,y,z)}{}B_{a\mu }\sigma ^\mu $$
(A25)
with $`a=(0,1,2,3)`$ and where
$$B=\left(\begin{array}{cccc}\hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill \frac{1}{\sqrt{2}}& \hfill 0& \hfill \frac{i}{\sqrt{2}}\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill \frac{1}{\sqrt{2}}& \hfill 0& \hfill \frac{i}{\sqrt{2}}\end{array}\right).$$
(A26)
In this way, Eq.(A22) is replaced by:
$$\overline{F}_{r,r^{}}^a(𝐤,𝐤^{})=\underset{\mu =(0,x,y,z)}{}B_{a\mu }F_{r,r^{}}^\mu (𝐤,𝐤^{}),$$
(A27)
and the correlation functions transform according to the rule:
$$\overline{𝒳}_{ab}=\underset{\mu ,\nu }{}B_{a\mu }𝒳_{\mu \nu }\left(B^T\right)_{\nu b}.$$
(A28)
The second term on the right-hand side of the integral equation (17) then transforms as follows:
$$B𝒳^{(0)}ϵ𝒳B^T=B𝒳^{(0)}B^T(B^T)^1ϵB^1B𝒳B^T=\overline{𝒳}^{(0)}\overline{ϵ}\overline{𝒳}$$
(A29)
in matrix notation, where
$$\overline{ϵ}=\left(B^T\right)^1ϵB^1=\left(\begin{array}{cccc}\hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0\end{array}\right);$$
(A30)
at the same time
$$\overline{T}(\mathrm{\Omega }_i)=BT(\mathrm{\Omega }_i)B^1=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& e^{i𝐐𝐑_i}& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& e^{i𝐐𝐑_i}\end{array}\right)$$
(A31)
becomes diagonal, as anticipated. Equation (30) of the text then follows.
There remains to solve Eq. (34) of the text explicitly. To this end, we introduce the compact notation
$$\gamma _a=\{\begin{array}{ccc}\hfill 0& & (a=0)\\ \hfill 1& & (a=1)\\ \hfill 0& & (a=2)\\ \hfill +1& & (a=3)\end{array},$$
(A32)
in such a way that (the lattice Fourier transform of) Eq. (A31) reads:
$$\overline{T}_{aa^{}}(𝐤)=\delta _{a,a^{}}\mathrm{\Delta }(𝐤+\gamma _a𝐐)$$
(A33)
$`\mathrm{\Delta }(𝐤)`$ being the lattice Kronecker delta function. In the transformed basis (cf. Eq. (A28)), the non-interacting part (32) of the correlation function then takes the form:
$$\widehat{𝒳}_{ab}^{(0)}(𝐪,𝐪^{};\omega )=\mathrm{\Delta }(𝐪𝐪^{}(\gamma _a+\gamma _b)𝐐)X_{ab}^{(0)}(𝐪;\omega |𝐐)$$
(A34)
where we have set
$`X_{ab}^{(0)}(𝐪;\omega |𝐐)`$ $`=`$ $`{\displaystyle \frac{1}{4𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}{\displaystyle \underset{r,r^{}}{}}\overline{F}_{r^{},r}^a(𝐤𝐪+\gamma _a𝐐,𝐤)\overline{F}_{r,r^{}}^b(𝐤,𝐤𝐪+\gamma _a𝐐)`$ (A35)
$`\times `$ $`_{rr^{}}(𝐤,𝐤𝐪+\gamma _a𝐐).`$ (A36)
Owing to the wave vector conserving Kronecker delta function in Eq. (A34), the integral equation (34) becomes:
$`\widehat{𝒳}_{ab}(𝐪,𝐪^{};\omega )`$ $`=`$ $`\mathrm{\Delta }(𝐪𝐪^{}(\gamma _a+\gamma _b)𝐐)X_{ab}^{(0)}(𝐪;\omega |𝐐)`$ (A37)
$`+`$ $`2U[X_{a0}^{(0)}(𝐪;\omega |𝐐)\widehat{𝒳}_{0b}(𝐪\gamma _a𝐐,𝐪^{};\omega )`$ (A38)
$``$ $`X_{a1}^{(0)}(𝐪;\omega |𝐐)\widehat{𝒳}_{3b}(𝐪(\gamma _a1)𝐐,𝐪^{};\omega )`$ (A39)
$``$ $`X_{a2}^{(0)}(𝐪;\omega |𝐐)\widehat{𝒳}_{2b}(𝐪\gamma _a𝐐,𝐪^{};\omega )`$ (A40)
$``$ $`X_{a3}^{(0)}(𝐪;\omega |𝐐)\widehat{𝒳}_{1b}(𝐪(\gamma _a+1)𝐐,𝐪^{};\omega )].`$ (A41)
Note that the wave vector arguments on the right-hand side of the above expression depend on the index $`a`$ of the matrix element. To avoid this feature, we let $`𝐪𝐪+\gamma _a𝐐`$ everywhere in the above expression and obtain:
$`\widehat{𝒳}_{ab}(𝐪+\gamma _a𝐐,𝐪^{};\omega )`$ $`=`$ $`\mathrm{\Delta }(𝐪𝐪^{}\gamma _b𝐐)X_{ab}^{(0)}(𝐪+\gamma _a𝐐;\omega |𝐐)`$ (A42)
$`+`$ $`2U[X_{a0}^{(0)}(𝐪+\gamma _a𝐐;\omega |𝐐)\widehat{𝒳}_{0b}(𝐪,𝐪^{};\omega )`$ (A43)
$``$ $`X_{a1}^{(0)}(𝐪+\gamma _a𝐐;\omega |𝐐)\widehat{𝒳}_{3b}(𝐪+𝐐,𝐪^{};\omega )`$ (A44)
$``$ $`X_{a2}^{(0)}(𝐪+\gamma _a𝐐;\omega |𝐐)\widehat{𝒳}_{2b}(𝐪,𝐪^{};\omega )`$ (A45)
$``$ $`X_{a3}^{(0)}(𝐪+\gamma _a𝐐;\omega |𝐐)\widehat{𝒳}_{1b}(𝐪𝐐,𝐐^{};\omega )].`$ (A46)
For any given value of the index $`b`$, this relation can then be cast in the equivalent form:
$$\left[\mathrm{𝟏}+X(𝐪,\omega )\right]\left(\begin{array}{c}\widehat{𝒳}_{0b}(𝐪,𝐪^{};\omega )\\ \widehat{𝒳}_{1b}(𝐪𝐐,𝐪^{};\omega )\\ \widehat{𝒳}_{2b}(𝐪,𝐪^{};\omega )\\ \widehat{𝒳}_{3b}(𝐪+𝐐,𝐪^{};\omega )\end{array}\right)=\left(\begin{array}{c}X_{0b}^{(0)}(𝐪;\omega |𝐐)\\ X_{1b}^{(0)}(𝐪𝐐;\omega |𝐐)\\ X_{2b}^{(0)}(𝐪;\omega |𝐐)\\ X_{3b}^{(0)}(𝐪+𝐐;\omega |𝐐)\end{array}\right)\mathrm{\Delta }(𝐪𝐪^{}\gamma _b𝐐)$$
(A47)
where we have set
$`X(𝐪,\omega )=\mathrm{\hspace{0.17em}2}U`$ (A48)
$`\times \left(\begin{array}{cccc}X_{00}^{(0)}(𝐪;\omega |𝐐)& X_{03}^{(0)}(𝐪;\omega |𝐐)& X_{02}^{(0)}(𝐪;\omega |𝐐)& X_{01}^{(0)}(𝐪;\omega |𝐐)\\ X_{10}^{(0)}(𝐪𝐐;\omega |𝐐)& X_{13}^{(0)}(𝐪𝐐;\omega |𝐐)& X_{12}^{(0)}(𝐪𝐐;\omega |𝐐)& X_{11}^{(0)}(𝐪𝐐;\omega |𝐐)\\ X_{20}^{(0)}(𝐪;\omega |𝐐)& X_{23}^{(0)}(𝐐;\omega |𝐐)& X_{22}^{(0)}(𝐪;\omega |𝐐)& X_{21}^{(0)}(𝐐;\omega |𝐐)\\ X_{30}^{(0)}(𝐪+𝐐;\omega |𝐐)& X_{33}^{(0)}(𝐪+𝐐;\omega |𝐐)& X_{32}^{(0)}(𝐪+𝐐;\omega |𝐐)& X_{31}^{(0)}(𝐐+𝐐;\omega |𝐐)\end{array}\right).`$ (A53)
(A54)
Solving for $`\widehat{𝒳}_{ab}`$ in Eq. (A47), we obtain eventually:
$`\widehat{𝒳}_{ab}(𝐪+\gamma _a𝐐,𝐪^{};\omega )`$ $`=`$ $`{\displaystyle \underset{a^{}}{}}[\mathrm{𝟏}+X(𝐪,\omega )]_{aa^{}}^1`$ (A55)
$`\times `$ $`X_{a^{}b}^{(0)}(𝐪+\gamma _a^{}𝐐;\omega |𝐐)\mathrm{\Delta }(𝐪𝐪^{}\gamma _b𝐐)`$ (A56)
which coincides with Eq. (36) of the text.
## B Details of the $`t/U`$ expansion
In this Appendix, we provide details of the $`t/U`$ expansion of the mean-field parameters and of the matrix elements of the correlation function which are necessary to obtain the spin-wave dispersion relation.
We begin by considering the leading ($`t=0`$) term of the expansion (53) for the mean-field (band) eigenvalues, corresponding to completely flat bands, which is given by \[cf. Eq. (46)\]
$$ϵ_r(𝐤)=U\left(ϵ_r^{(0)}\mu ^{(0)}\right)=U(m_1+(1)^rm_2)\mu _0$$
(B1)
with $`r=1,2`$ and where $`\mu _0=U\mu ^{(0)}`$ is the chemical potential at the lowest order in $`t/U`$. Since at this order the band eigenvalues do not depend on the wave vector, there is no value of the chemical potential $`\mu ^{(0)}`$ satisfying equation (49) at zero temperature and noninteger doping $`\delta `$. More precisely, the doping jumps from $`\delta =1`$ when $`\mu ^{(0)}<ϵ_1^{(0)}`$, to $`\delta =0`$ when $`ϵ_1^{(0)}<\mu ^{(0)}<ϵ_2^{(0)}`$, and to $`\delta =+1`$ when $`\mu ^{(0)}>ϵ_2^{(0)}`$. In order to consider a continuously doped system, it is therefore necessary to include at the outset the next-to-leading order in the expansion (53) of the band eigenvalues, which introduces a band dispersion.
For definiteness, we shall consider the case $`\delta >0`$ from now on. (The case $`\delta <0`$ can be recovered by exploiting particle-hole symmetry.) In the case $`\delta >0`$, the lowest band is always filled, and we will consistently use (the zero-temperature value) $`f_F(ϵ_1(𝐤))=1`$ throughout. On the other hand, the zeroth order chemical potential for $`0<\delta <1`$ must be chosen as $`\mu ^{(0)}=ϵ_2^{(0)}`$, and the Fermi function in the upper band can be expanded as
$`f_F(ϵ_2(𝐤))=\mathrm{\hspace{0.17em}1}\mathrm{\Theta }\left(ϵ_2^{(1)}(𝐤)\mu ^{(1)}\right)`$ (B2)
$``$ $`\delta \left(ϵ_2^{(1)}(𝐤)\mu ^{(1)}\right)\left[\left({\displaystyle \frac{t}{U}}\right)(ϵ_2^{(2)}(𝐤)\mu ^{(2)})+\left({\displaystyle \frac{t}{U}}\right)^2(ϵ_2^{(3)}(𝐤)\mu ^{(3)})\right]`$ (B3)
$``$ $`{\displaystyle \frac{1}{2}}\delta ^{}\left(ϵ_2^{(1)}(𝐤)\mu ^{(1)}\right)\left({\displaystyle \frac{t}{U}}\right)^2(ϵ_2^{(2)}(𝐤)\mu ^{(2)})^2+\mathrm{},`$ (B4)
where we have considered the zero-temperature limit of the Fermi function and of its derivatives, and introduced the step function $`\mathrm{\Theta }`$ as well as the Dirac $`\delta (x)`$ function (not to be confused with doping). In this way, Eq. (49) can be expressed in powers of $`t/U`$, yielding
$`{\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\mathrm{\Theta }\left(ϵ_2^{(1)}(𝐤)\mu ^{(1)}\right)=\mathrm{\hspace{0.17em}1}\delta ,`$ (B5)
$`{\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\delta \left(ϵ_2^{(1)}(𝐤)\mu ^{(1)}\right)(ϵ_2^{(2)}(𝐤)\mu ^{(2)})=\mathrm{\hspace{0.17em}0},`$ (B6)
$`{\displaystyle \frac{1}{2𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\left[\delta ^{}\left(ϵ_2^{(1)}(𝐤)\mu ^{(1)}\right)(ϵ_2^{(2)}(𝐤)\mu ^{(2)})^2+2\delta \left(ϵ_2^{(1)}(𝐤)\mu ^{(1)}\right)(ϵ_2^{(3)}(𝐤)\mu ^{(3)})\right]=\mathrm{\hspace{0.17em}0}.`$ (B7)
(B8)
Similarly, the small $`t/U`$ expansion (54) of Eq. (50) yields \[cf. also Eqs. (55)-(57)\]:
$`m_2^{(0)}+\left({\displaystyle \frac{t}{U}}\right)m_2^{(1)}+\left({\displaystyle \frac{t}{U}}\right)^2m_2^{(2)}+\mathrm{}`$ (B9)
$`={\displaystyle \frac{1}{2𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\left[12\left({\displaystyle \frac{t}{U}}\right)^2\left({\displaystyle \frac{T_o(𝐤)}{2m_2^{(0)}}}\right)^2\right]\left(1f_F(ϵ_2(𝐤))\right)+𝒪\left(\left({\displaystyle \frac{U}{t}}\right)^3\right),`$ (B10)
where the Fermi function should be again expanded in the form (B2). One obtains eventually:
$`m_2^{(0)}`$ $`=`$ $`{\displaystyle \frac{1}{2𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\mathrm{\Theta }\left(ϵ_2^{(1)}(𝐤)\mu ^{(1)}\right),`$ (B11)
$`m_2^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\delta \left(ϵ_2^{(1)}(𝐤)\mu ^{(1)}\right)(ϵ_2^{(2)}(𝐤)\mu ^{(2)}),`$ (B12)
$`m_2^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\left({\displaystyle \frac{T_o(𝐤)}{2m_2^{(0)}}}\right)^2\mathrm{\Theta }\left(ϵ_2^{(1)}(𝐤)\mu ^{(1)}\right)`$ (B13)
$`+`$ $`{\displaystyle \frac{1}{4𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\delta ^{}\left(ϵ_2^{(1)}(𝐤)\mu ^{(1)}\right)(ϵ_2^{(2)}(𝐤)\mu ^{(2)})^2`$ (B14)
$`+`$ $`{\displaystyle \frac{1}{2𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\delta \left(ϵ_2^{(1)}(𝐤)\mu ^{(1)}\right)(ϵ_2^{(3)}(𝐤)\mu ^{(3)})`$ (B15)
$`=`$ $`{\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\left({\displaystyle \frac{T_o(𝐤)}{2m_2^{(0)}}}\right)^2\mathrm{\Theta }\left(ϵ_2^{(1)}(𝐤)\mu ^{(1)}\right),`$ (B16)
where in the last step we have used the last of Eqs. (B8).
There remains to expand in powers of $`t/U`$ the last the self-consistency equation (51). For the diagonal spin-spiral solution we obtain:
$`0=\mathrm{sin}(Q/2){\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}(\mathrm{cos}k_x+\mathrm{cos}k_y)\left(1+f_F(ϵ_2(𝐤))\right)`$ (B18)
$`+\left({\displaystyle \frac{t}{U}}\right)\mathrm{cos}(Q/2){\displaystyle \frac{1}{m_2^{(0)}}}{\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}T_o(𝐤)(\mathrm{sin}k_x+\mathrm{sin}k_y)\left(1f_F(ϵ_2(𝐤))\right)`$ (B19)
where the expansion (B2) has still to be inserted.
From Eq. (46), the argument of the $`\mathrm{\Theta }`$ function is $`T_e(𝐤)+m_2^{(1)}\mu ^{(1)}`$ (with $`T_e(𝐤)=2\mathrm{cos}(Q/2)(\mathrm{cos}k_x+\mathrm{cos}k_y)`$ for the diagonal spiral solution), which can take positive as well as negative values.
At the various orders in $`t/U`$ we then obtain from Eqs. (B8) and (LABEL:B-a2):
(i) order $`(t/U)^0`$
$`{\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\left[1\mathrm{\Theta }\left(T_e(𝐤)+m_2^{(1)}\mu ^{(1)}\right)\right]=\delta ,`$ (B20)
$`m_2^{(0)}={\displaystyle \frac{1}{2𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\mathrm{\Theta }\left(T_e(𝐤)+m_2^{(1)}\mu ^{(1)}\right);`$ (B21)
(ii) order $`(t/U)^1`$
$`{\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\delta \left(T_e(𝐤)+m_2^{(1)}\mu ^{(1)}\right)(ϵ_2^{(2)}(𝐤)\mu ^{(2)})=0,`$ (B22)
$`m_2^{(1)}={\displaystyle \frac{1}{2𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\delta \left(T_e(𝐤)+m_2^{(1)}\mu ^{(1)}\right)(ϵ_2^{(2)}(𝐤)\mu ^{(2)});`$ (B23)
(iii) order $`(t/U)^2`$
$$m_2^{(2)}=\frac{1}{𝒩}\underset{𝐤}{\overset{BZ}{}}\left(\frac{T_o(𝐤)}{2m_2^{(0)}}\right)^2\mathrm{\Theta }\left(T_e(𝐤)+m_2^{(1)}\mu ^{(1)}\right),$$
(B24)
and we don’t need to evaluate $`(ϵ_2^{(3)}(𝐤)\mu ^{(3)})`$ at the order we are considering. It is convenient to solve these equations order-by-order for the two variables $`\mu ^{(i)}m_2^{(i)}`$, and $`m_2^{(i)}`$ ($`i=0,1,2`$).
Using at this point the method developed in Appendix C for performing the $`𝐤`$ summation over the relevant portions of the Brillouin zone, we obtain eventually the following results for the mean-field parameters at the leading orders in $`\delta `$:
$$\{\begin{array}{ccc}\hfill m_2^{(0)}& =& (1\delta )/2\hfill \\ \hfill m_2^{(1)}& =& 0\hfill \\ \hfill m_2^{(2)}& =& 4\mathrm{sin}^2(Q/2)+𝒪(\delta ),\hfill \end{array}$$
(B25)
and
$$\{\begin{array}{ccc}\hfill \mu ^{(0)}m_2^{(0)}& =& m_1=(1+\delta )/2\hfill \\ \hfill \mu ^{(1)}m_2^{(1)}& =& 4\mathrm{cos}(Q/2)(1\pi \delta +𝒪(\delta ^2))\hfill \\ \hfill \mu ^{(2)}m_2^{(2)}& =& 0+𝒪(\delta ).\hfill \end{array}$$
(B26)
There remains to determine the magnitude $`Q`$ of the characteristic wave vector from the self-consistency condition (B18), where for the Fermi function we use the expansion (B2). Equation (B18) then becomes:
$`\left({\displaystyle \frac{t}{U}}\right)\mathrm{sin}(Q/2){\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}(\mathrm{cos}k_x+\mathrm{cos}k_y)[2\mathrm{\Theta }(T_e(𝐤)+m_2^{(1)}\mu ^{(1)})`$ (B27)
$`\left({\displaystyle \frac{t}{U}}\right)\delta (T_e(𝐤)+m_2^{(1)}\mu ^{(1)})(ϵ_2^{(2)}(𝐤)\mu ^{(2)})]+2({\displaystyle \frac{t}{U}})^2\mathrm{cos}(Q/2)`$ (B28)
$`\times {\displaystyle \frac{1}{2m_2^{(0)}}}{\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}(\mathrm{sin}k_x+\mathrm{sin}k_y)T_o(𝐤)\mathrm{\Theta }\left(T_e(𝐤)+m_2^{(1)}\mu ^{(1)}\right)`$ (B29)
$`=0.`$ (B30)
In particular, at the lowest order in $`t/U`$ we obtain:
$$\mathrm{sin}(Q/2)(2\delta +𝒪(\delta ^2))=\mathrm{\hspace{0.17em}0}.$$
(B31)
For fixed $`\delta `$ and large $`U/t`$, one thus has only the *ferromagnetic* solution $`Q=0`$. However, if one allows $`\delta `$ to be of the order $`t/U`$, Eq.(B31) has to be considered together with the next term in $`t/U`$, and we obtain instead:
$$\mathrm{sin}(Q/2)(2\delta +𝒪(\delta ^2))+4\mathrm{sin}(Q/2)\mathrm{cos}(Q/2)(1+𝒪(\delta ))=\mathrm{\hspace{0.17em}0}$$
(B32)
which yields the two solutions
$`\begin{array}{cc}\mathrm{sin}(Q/2)=\mathrm{\hspace{0.17em}0}& \text{ (ferromagnet) }\\ \mathrm{cos}(Q/2)=\delta U/2t& \text{ (diagonal spiral) },\end{array}`$
with the spiral solution being energetically favored over the ferromagnetic solution whenever it exists. From the above equation, it is clear that the spiral solution exists for any $`\delta 2t/U`$, which is consistent with our assumption that the doping parameter $`\delta `$ is at most of order $`t/U`$ in the spiral phase, thus justifying the expansion of Appendix C. The transition to the ferromagnetic state is second order, as the incommensurability vector $`Q`$ decreases continuously from $`Q=\pi `$ to $`Q=2\pi `$ with increasing $`\delta `$.
We pass finally to consider the $`t/U`$ expansion of the matrix elements of the correlation function. To begin with, we use the expressions (55)-(57) for the eigenvectors of the mean-field Hamiltonian and obtain from the definitions (39) and (40):
$`𝒳_0^{+,}(𝐪,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{2𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\{[1\left({\displaystyle \frac{t}{U}}\right)^2(T_o^2(𝐤)+T_o^2(𝐤𝐪))]_{2,1}(𝐤,𝐤𝐪,\omega )`$ (B33)
$`+`$ $`\left({\displaystyle \frac{t}{U}}\right)^2[T_o^2(𝐤)_{1,1}(𝐤,𝐤𝐪,\omega )+T_o^2(𝐤𝐪)_{2,2}(𝐤,𝐤𝐪,\omega )]\}`$ (B34)
$`+`$ $`𝒪\left(\left(t/U\right)^3\right)`$ (B35)
$`𝒳_0^{+,+}(𝐪,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{2𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\{\left({\displaystyle \frac{t}{U}}\right)^2T_o(𝐤)T_o(𝐤𝐪)`$ (B36)
$`+`$ $`[_{1,2}(𝐤,𝐤𝐪,\omega )+_{2,1}(𝐤,𝐤𝐪,\omega )`$ (B37)
$``$ $`_{1,1}(𝐤,𝐤𝐪,\omega )_{2,2}(𝐤,𝐤𝐪,\omega )]\}+𝒪((t/U)^3)`$ (B38)
for the matrix elements entering the final expression (98) of the spin-wave dispersion relation. For the remaining matrix elements we obtain instead:
$`𝒳_0^{0,0}(𝐪,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\{\left({\displaystyle \frac{t}{U}}\right)^2(T_o(𝐤)T_o(𝐤𝐪))^2`$ (B40)
$`\times `$ $`\left[_{1,2}(𝐤,𝐤𝐪,\omega )+_{2,1}(𝐤,𝐤𝐪,\omega )\right]`$ (B41)
$`+`$ $`\left[1\left({\displaystyle \frac{t}{U}}\right)^2\left(T_o(𝐤)T_o(𝐤𝐪)\right)^2\right]`$ (B42)
$`\times `$ $`[_{1,1}(𝐤,𝐤𝐪,\omega )+_{2,2}(𝐤,𝐤𝐪,\omega )]\}+𝒪((t/U)^3),`$ (B43)
$`𝒳_0^{0,}(𝐪,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}i\sqrt{2}\{\left({\displaystyle \frac{t}{U}}\right){\displaystyle \frac{1}{1\delta }}(T_o(𝐤𝐪)T_o(𝐤))_{2,1}(𝐤,𝐤𝐪,\omega )`$ (B45)
$`+`$ $`\left({\displaystyle \frac{t}{U}}\right){\displaystyle \frac{1}{1\delta }}[T_o(𝐤)_{1,1}(𝐤,𝐤𝐪,\omega )T_o(𝐤𝐪)_{2,2}(𝐤,𝐤𝐪,\omega )]\}`$ (B46)
$`+`$ $`𝒪\left(\left(t/U\right)^3\right),`$ (B47)
$`𝒳_0^{0,z}(𝐪,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\{\left({\displaystyle \frac{t}{U}}\right)^2(T_o^2(𝐤)T_o^2(𝐤𝐪))`$ (B48)
$`\times `$ $`\left[_{2,1}(𝐤,𝐤𝐪,\omega )_{1,2}(𝐤,𝐤𝐪,\omega )\right]`$ (B49)
$`+`$ $`\left[1\left({\displaystyle \frac{t}{U}}\right)^2\left(T_o^2(𝐤)+T_o^2(𝐤𝐪)\right)\right]`$ (B50)
$`\times `$ $`[_{1,1}(𝐤,𝐤𝐪,\omega )_{2,2}(𝐤,𝐤𝐪,\omega )]\}+𝒪((t/U)^3),`$ (B51)
$`𝒳_0^{+,z}(𝐪,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}i\sqrt{2}\{\left({\displaystyle \frac{t}{U}}\right){\displaystyle \frac{1}{1\delta }}(T_o(𝐤𝐪)+T_o(𝐤))_{2,1}(𝐤,𝐤𝐪,\omega )`$ (B53)
$`+`$ $`\left({\displaystyle \frac{t}{U}}\right){\displaystyle \frac{1}{1\delta }}[T_o(𝐤)_{1,1}(𝐤,𝐤𝐪,\omega )T_o(𝐤𝐪)_{2,2}(𝐤,𝐤𝐪,\omega )]\}`$ (B54)
$`+`$ $`𝒪\left(\left(t/U\right)^3\right),`$ (B55)
$`𝒳_0^{z,z}(𝐪,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}\{\left({\displaystyle \frac{t}{U}}\right)^2(T_o(𝐤)+T_o(𝐤𝐪))^2`$ (B56)
$`\times `$ $`\left[_{1,2}(𝐤,𝐤𝐪,\omega )+_{2,1}(𝐤,𝐤𝐪,\omega )\right]`$ (B57)
$`+`$ $`\left[1\left({\displaystyle \frac{t}{U}}\right)^2\left(T_o(𝐤)+T_o(𝐤𝐪)\right)^2\right]`$ (B58)
$`\times `$ $`[_{1,1}(𝐤,𝐤𝐪,\omega )+_{2,2}(𝐤,\mathrm{𝐤𝐤}𝐪,\omega )]\}+𝒪((t/U)^3).`$ (B59)
Note that in the above expressions $`_{1,1}(𝐤,𝐤𝐪,\omega )=0`$ when the doping parameter $`\delta 0`$. In addition, $`_{1,2}`$ can be obtained from $`_{2,1}`$ through the symmetry condition (68). Making use of the expansion (53) for the eigenvalues of the mean-field Hamiltonian, we write further
$`_{2,2}(𝐤,𝐤q)`$ $`=`$ $`{\displaystyle \frac{1}{t}}\left[f_F(ϵ_2(𝐤q))f_F(ϵ_2(𝐤))\right]{\displaystyle \frac{1}{T_e(𝐤q)T_e(𝐤)}}`$ (B60)
$`+`$ $`𝒪\left(1/U\right)`$ (B61)
as well as
$`_{2,1}(𝐤,𝐤q)`$ $`=`$ $`[f_F(ϵ_2(𝐤))1]{\displaystyle \frac{1}{U}}\{{\displaystyle \frac{1}{(1\delta \stackrel{~}{\omega })}}`$ (B62)
$``$ $`\left({\displaystyle \frac{t}{U}}\right){\displaystyle \frac{1}{(1\delta \stackrel{~}{\omega })^2}}[T_e(𝐤)T_e(𝐤q)]+\left({\displaystyle \frac{t}{U}}\right)^2[{\displaystyle \frac{1}{(1\delta \stackrel{~}{\omega })^2}}`$ (B63)
$`\times `$ $`\left(T_o^2(𝐤)T_o^2(𝐤q)2m_2^{(2)}\right)`$ (B64)
$`+`$ $`{\displaystyle \frac{1}{(1\delta \stackrel{~}{\omega })^3}}(T_e(𝐤)T_e(𝐤q))^2]\}`$ (B65)
$`+`$ $`𝒪\left(t^3/U^4\right)`$ (B66)
with $`\stackrel{~}{\omega }=\omega /U`$, where the expansion (B2) of the Fermi function has still to be inserted.
## C Sums over the two-dimensional BZ in the small $`\delta `$ limit
In this Appendix, we develop a method suitable for performing the $`𝐤`$ summation over the relevant portions of the Brillouin zone in powers of the doping parameter $`\delta `$. This is, in turn, justified by the fact that $`\delta `$ in the spiral phase is of order $`t/U`$, as shown in Appendix B.
The typical integral to be evaluated is of the form:
$$(\gamma )=\frac{1}{𝒩}\underset{𝐤}{\overset{BZ}{}}g(𝐤)\mathrm{\Theta }\left(T_e(𝐤)+\gamma \right)$$
(C1)
where $`T_e(𝐤)=2\mathrm{cos}(Q/2)(\mathrm{cos}k_x+\mathrm{cos}k_y)`$ for the diagonal spiral solution we are considering and $`\gamma `$ is a parameter (which depends on the chemical potential $`\mu ^{(1)}`$ at the order $`𝒪\left[(t/U)^0\right])`$ that in the following we shall simply call $`\mu `$ for simplicity. Introducing the notation $`f_\mu `$ such that $`\gamma =4\mathrm{cos}(Q/2)(1f_\mu /4)`$ for any given $`Q`$ value, recalling the definition of the $`\mathrm{\Theta }`$ function, and considering that $`\mathrm{cos}(Q/2)<0`$ in our solution, we set further:
$$(\gamma )=\frac{1}{𝒩}\underset{\mathrm{cos}k_x+\mathrm{cos}k_y2f_\mu /2}{}g(𝐤).$$
(C2)
Quite generally, we can introduce a polar representation for the two-dimensional wave vector $`𝐤`$ and determine for each value of the polar angle $`\varphi `$ the magnitude $`k(\varphi )`$, such that the equality
$$\mathrm{cos}\left(k(\varphi )\mathrm{cos}\varphi \right)+\mathrm{cos}\left(k(\varphi )\mathrm{sin}\varphi \right)=\mathrm{\hspace{0.17em}2}\frac{f_\mu }{2}$$
(C3)
is satisfied. In this way we rewrite the integral (C2) as follows:
$$(\gamma )=\frac{1}{4\pi ^2}_{BZ}𝑑k_x𝑑k_yg(k_x,k_y)\frac{1}{4\pi ^2}_0^{2\pi }𝑑\varphi _0^{k(\varphi )}𝑑kkg(k,\varphi ),$$
(C4)
with a slight (albeit harmless) abuse of notation for the function $`g`$.
In practice, when $`f_\mu `$ is small compared to unity (as it is the case for small $`\delta `$) it is convenient to determine $`k(\varphi )`$ in powers of $`f_\mu `$ by expanding $`\mathrm{cos}\left(k(\varphi )\mathrm{cos}\varphi \right)+\mathrm{cos}\left(k(\varphi )\mathrm{sin}\varphi \right)`$ in a Taylor series about $`k=0`$. For instance, at the lowest order we obtain from Eq. (C3):
$$k(\varphi )^2+𝒪\left(k(\varphi )^4\right)=f_\mu $$
(C5)
which gives $`k(\varphi )^2=f_\mu +𝒪\left(f_\mu ^2\right)`$. At the next significant order we obtain instead:
$$k(\varphi )^2\frac{(\mathrm{sin}^4\varphi +\mathrm{cos}^4\varphi )}{12}k(\varphi )^4+𝒪\left(k(\varphi )^6\right)=f_\mu $$
(C6)
which gives
$$k(\varphi )^2=f_\mu +\frac{(\mathrm{sin}^4\varphi +\mathrm{cos}^4\varphi )}{12}f_\mu ^2+𝒪\left(f_\mu ^3\right).$$
(C7)
Consider, for instance, the first of Eqs. (B21), which is of the form (C1) with $`g(𝐤)=1`$ and $`\gamma =m_2^{(1)}\mu ^{(1)}`$. Equation (C4 ) now becomes:
$$\delta =\frac{1}{8\pi ^2}_0^{2\pi }𝑑\varphi k(\varphi )^2.$$
(C8)
Introducing the notation $`f_\mu `$ as above, for small values of $`\delta (>0)`$ we obtain:
$$\delta =\frac{1}{4\pi }\left[f_\mu +\frac{1}{16}f_\mu ^2+𝒪\left(f_\mu ^3\right)\right]$$
(C9)
where use has been made of Eq. (C7). Inverting this relation we obtain eventually:
$$f_\mu =\mathrm{\hspace{0.17em}4}\pi \delta \pi ^2\delta ^2+𝒪\left(\delta ^3\right),$$
(C10)
which confirms the fact that $`f_\mu =𝒪(\delta )`$.
In general, the case of a $`𝐤`$-dependent $`g`$ can be treated by expanding $`g(𝐤)`$ about $`|𝐤|=0`$, since the last integral on the rignt-hand side of Eq. (C4) is restricted to $`|𝐤|<k(\varphi )=𝒪(\delta ^{1/2})`$, and by retaining the relevant orders in $`\delta `$ in the final expression. The integral over the whole BZ (i.e., the first term on the rignt-hand side of Eq. (C4)), on the other hand, can usually be performed analytically.
The Brillouin zone sums involving the Dirac delta functions or its derivatives can be further evaluated according to the following device. Writing
$`\delta \left(T_e(𝐤)+\gamma \right)`$ $`=`$ $`{\displaystyle \frac{}{\gamma }}\mathrm{\Theta }\left(T_e(𝐤)+\gamma \right)`$ (C11)
$`=`$ $`{\displaystyle \frac{1}{\mathrm{cos}(Q/2)}}{\displaystyle \frac{}{f_\mu }}\mathrm{\Theta }\left(T_e(𝐤)+\gamma \right),`$ (C12)
we obtain for the Brillouin zone sum
$`{\displaystyle \frac{1}{𝒩}}`$ $`{\displaystyle \underset{𝐤}{\overset{BZ}{}}}g(𝐤)\delta ^{(n)}\left(T_e(𝐤)+m_2^{(1)}\mu ^{(1)}\right)`$ (C13)
$`=`$ $`{\displaystyle \frac{1}{(\mathrm{cos}(Q/2))^{(n+1)}}}{\displaystyle \frac{^{(n+1)}}{(f_\mu ^{(1)})^{(n+1)}}}{\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}g(𝐤)\mathrm{\Theta }\left(T_e(𝐤)+m_2^{(1)}\mu ^{(1)}\right)`$ (C14)
$`=`$ $`{\displaystyle \frac{1}{(\mathrm{cos}(Q/2))^{(n+1)}}}{\displaystyle \frac{^{(n+1)}}{(f_\mu ^{(1)})^{(n+1)}}}{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _O^{k(\varphi )}}𝑑kkg(k,\varphi )`$ (C15)
since the first term on the right-hand side of Eq. (C4) does not depend on $`f_\mu `$.
As an example, we evaluate the expression
$`{\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}T_o^2(𝐤)\delta \left(T_e(𝐤)+m_2^{(1)}\mu ^{(1)}\right)`$
which appears in the first of Eqs. (B23). According to Eq. (C12), we first evaluate the expression (cf. also Eq. (C4)):
$`{\displaystyle \frac{1}{𝒩}}`$ $`{\displaystyle \underset{𝐤}{\overset{BZ}{}}}T_o^2(𝐤)\mathrm{\Theta }\left(T_e(𝐤)+m_2^{(1)}\mu ^{(1)}\right)`$ (C16)
$`=`$ $`4\mathrm{sin}^2(Q/2){\displaystyle \frac{1}{\pi ^2}}\mathrm{sin}^2(Q/2){\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _0^{k(\varphi )}}𝑑kk\left(\mathrm{sin}(k\mathrm{cos}\varphi )+\mathrm{sin}(k\mathrm{sin}\varphi )\right)^2.`$ (C17)
Expanding the factor within parentheses in the last integral in powers of $`k`$, as explained below Eq. (C10), we obtain:
$`{\displaystyle \frac{1}{𝒩}}`$ $`{\displaystyle \underset{𝐤}{\overset{BZ}{}}}T_o^2(𝐤)\mathrm{\Theta }\left(T_e(𝐤)+m_2^{(1)}\mu ^{(1)}\right)`$ (C18)
$`=`$ $`4\mathrm{sin}^2(Q/2){\displaystyle \frac{1}{\pi ^2}}\mathrm{sin}^2(Q/2){\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _0^{k(\varphi )}}𝑑kk\left(k^2(\mathrm{cos}\varphi +\mathrm{sin}\varphi )^2+𝒪\left(k^4\right)\right)`$ (C19)
$`=`$ $`\mathrm{sin}^2(Q/2)\left(48\pi \delta ^2+𝒪(\delta ^3)\right).`$ (C20)
Using Eq. (C12) we obtain eventually:
$`{\displaystyle \frac{1}{𝒩}}`$ $`{\displaystyle \underset{𝐤}{\overset{BZ}{}}}T_o^2(𝐤)\delta \left(T_e(𝐤)+m_2^{(1)}\mu ^{(1)}\right)`$ (C21)
$`=`$ $`{\displaystyle \frac{1}{\mathrm{cos}(Q/2)}}{\displaystyle \frac{}{f_\mu }}{\displaystyle \frac{1}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}T_o^2(𝐤)\mathrm{\Theta }\left(T_e(𝐤)+m_2^{(1)}\mu ^{(1)}\right)`$ (C22)
$`=`$ $`{\displaystyle \frac{\mathrm{sin}^2(Q/2)}{\mathrm{cos}(Q/2)}}{\displaystyle \frac{1}{\pi }}f_\mu +𝒪((f_\mu ^2)=4\delta {\displaystyle \frac{\mathrm{sin}^2(Q/2)}{\mathrm{cos}(Q/2)}}+𝒪(\delta ^2).`$ (C23)
With these prescriptions, the results (B25), (B26), and (B30) are readily obtained.
It was shown in Section 3 that, as far as the matrix elements of the correlation function are concerned, only the explicit expressions of $`d(𝐪,\omega )`$ and $`e(𝐪,\omega )`$ are required at the order in $`t/U`$ we are considering in this paper. Writing
$`d(𝐪,\omega )`$ $`=`$ $`1\stackrel{~}{\omega }+\alpha (𝐪)+𝒪\left(\left(t/U\right)^3\right)`$ (C24)
$`=`$ $`I(𝐪,\omega )J(𝐪,\omega )K(𝐪,\omega )+M(𝐪,\omega )`$ (C25)
and
$$e(𝐪,\omega )=L(𝐪,\omega )+L(𝐪,\omega )+N(𝐪,\omega )$$
(C26)
where the quantities $`I,J,\mathrm{}`$ are specified below, we obtain for the sums over the wave vector $`𝐤`$ using the method described in this Appendix:
$`I(𝐪,\omega )`$ $``$ $`{\displaystyle \frac{U}{𝒩}}{\displaystyle \underset{𝐤}{\overset{BZ}{}}}_{21}(𝐤,𝐤q,\omega )={\displaystyle \frac{1+\delta }{1\stackrel{~}{\omega }\delta }}`$ (C27)
$`+`$ $`\left({\displaystyle \frac{t}{U}}\right)\mathrm{\hspace{0.17em}2}\delta \mathrm{cos}(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y)`$ (C28)
$`+`$ $`\left({\displaystyle \frac{t}{U}}\right)^2\mathrm{\hspace{0.17em}4}\mathrm{cos}^2(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y)+𝒪\left(\left(t/U\right)^3\right),`$ (C29)
$`J(𝐪,\omega )`$ $``$ $`{\displaystyle \frac{U}{𝒩}}\left({\displaystyle \frac{t}{U}}\right)^2{\displaystyle \underset{𝐤}{\overset{BZ}{}}}T_o^2(𝐤)_{21}(𝐤,𝐤q,\omega )`$ (C30)
$`=`$ $`4\left({\displaystyle \frac{t}{U}}\right)^2\mathrm{sin}^2(Q/2)+𝒪\left(\left(t/U\right)^3\right),`$ (C31)
$`K(𝐪,\omega )`$ $``$ $`{\displaystyle \frac{U}{𝒩}}\left({\displaystyle \frac{t}{U}}\right)^2{\displaystyle \underset{𝐤}{\overset{BZ}{}}}T_o^2(𝐤𝐪)_{21}(𝐤,𝐤q,\omega )`$ (C32)
$`=`$ $`4\left({\displaystyle \frac{t}{U}}\right)^2\mathrm{sin}^2(Q/2)+𝒪\left(\left(t/U\right)^3\right),`$ (C33)
$`L(𝐪,\omega )`$ $``$ $`{\displaystyle \frac{U}{𝒩}}\left({\displaystyle \frac{t}{U}}\right)^2{\displaystyle \underset{𝐤}{\overset{BZ}{}}}T_o(𝐤)T_o(𝐤𝐪)_{21}(𝐤,𝐤q,\omega )`$ (C34)
$`=`$ $`2\left({\displaystyle \frac{t}{U}}\right)^2\mathrm{sin}^2(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y)+𝒪\left(\left(t/U\right)^3\right),`$ (C35)
$`M(𝐪,\omega )`$ $``$ $`{\displaystyle \frac{U}{𝒩}}\left({\displaystyle \frac{t}{U}}\right)^2{\displaystyle \underset{𝐤}{\overset{BZ}{}}}T_o^2(𝐤𝐪)_{22}(𝐤,𝐤q,\omega )`$ (C36)
$`=`$ $`2\delta \left({\displaystyle \frac{t}{U}}\right){\displaystyle \frac{\mathrm{sin}^2(Q/2)(\mathrm{sin}q_x+\mathrm{sin}q_y)^2}{\mathrm{cos}(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2)}}+𝒪\left(\left(t/U\right)^3\right),`$ (C37)
$`N(𝐪,\omega )`$ $``$ $`{\displaystyle \frac{U}{𝒩}}\left({\displaystyle \frac{t}{U}}\right)^2{\displaystyle \underset{𝐤}{\overset{BZ}{}}}T_o(𝐤)T_o(𝐤𝐪)_{22}(𝐤,𝐤q,\omega )`$ (C38)
$`=`$ $`\left({\displaystyle \frac{t}{U}}\right)^2𝒪(\delta )=𝒪\left(\left(t/U\right)^3\right),`$ (C39)
where these results hold for $`|𝐪|k_F`$, as discussed in Section 3-B.
For the relevant matrix elements (C25) and (C26) we obtain eventually :
$`d(𝐪,\omega )`$ $`=`$ $`1\stackrel{~}{\omega }+\left({\displaystyle \frac{t}{U}}\right)^2\mathrm{\hspace{0.17em}4}\mathrm{cos}^2(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2)+8\left({\displaystyle \frac{t}{U}}\right)^2\mathrm{sin}^2(Q/2)`$ (C40)
$`+`$ $`2\left({\displaystyle \frac{t}{U}}\right)\delta \mathrm{cos}(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2)`$ (C41)
$``$ $`2\left({\displaystyle \frac{t}{U}}\right)\delta {\displaystyle \frac{\mathrm{sin}^2(Q/2)(\mathrm{sin}q_x+\mathrm{sin}q_y)^2}{\mathrm{cos}(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y2)}}+𝒪\left(\left(t/U\right)^3\right)`$ (C42)
and
$$e(𝐪,\omega )=4\left(\frac{t}{U}\right)^2\mathrm{sin}^2(Q/2)(\mathrm{cos}q_x+\mathrm{cos}q_y)+𝒪\left(\left(t/U\right)^3\right).$$
(C43)
## D Heisenberg spin waves for a spiral configuration
In this Appendix, we give a simplified derivation of the spin-wave dispersion relation for a two-dimensional Heisenberg antiferromagnet in the presence of a spiral incommensurate magnetic ground state. Although the form of this spectrum is well known, it is worthed to give here a compact derivation in terms of the set of local spin quantization axes utilized in Section 3 and in Appendix A.
We consider the Heisenberg Hamiltonian
$$H=\frac{1}{2}\underset{i,j}{}J_{i,j}𝐒_i𝐒_j$$
(D1)
for spin $`1/2`$, where the sum over the lattice sites $`i`$ and $`j`$ extends in principle over distant neighbors. We transform next the spin locally as follows:
$$S_i^\alpha \underset{\beta }{}T_{\alpha \beta }^{(i)}S_i^\beta $$
(D2)
where $`T^{(i)}`$ represents the $`3\times 3`$ matrix given by the spin part of (A17). The Hamiltonian (D1) then becomes:
$$H=\frac{1}{2}\underset{i,j}{}J_{i,j}\underset{\alpha ,\beta }{}S_i^\alpha T_{\alpha \beta }^{(ji)}S_j^\beta $$
(D3)
where $`T^{(ji)}`$ contains the difference $`\theta _j\theta _i`$ in the place of $`\theta _i`$ of Eq. (A17).
Writing further
$$S_i^x=\frac{1}{2}\left(S_i^++S_i^{}\right),S_i^y=\frac{1}{2i}\left(S_i^+S_i^{}\right),$$
(D4)
and carrying out the standard Holstein-Primakoff transformation
$`S_i^\pm `$ $``$ $`\sqrt{2S}a_i^\pm `$ (D5)
$`S_i^z`$ $`=`$ $`a_i^{}a_iS=n_iS`$ (D6)
at the leading order away from perfect alignment (and for large values of $`S`$), the Hamiltonian (D3) becomes:
$`H`$ $``$ $`{\displaystyle \frac{S^2}{2}}{\displaystyle \underset{i,j}{}}J_{i,j}\mathrm{cos}(\theta _j\theta _i)S{\displaystyle \underset{i,j}{}}n_iJ_{i,j}\mathrm{cos}(\theta _j\theta _i)`$ (D7)
$`+`$ $`{\displaystyle \frac{S}{4}}{\displaystyle \underset{i,j}{}}J_{i,j}\{(a_i^{}a_j+a_j^{}a_i)[\mathrm{cos}(\theta _j\theta _i)+1]`$ (D8)
$`+`$ $`(a_i^{}a_j^{}+a_ia_j)[\mathrm{cos}(\theta _j\theta _i)1]\}.`$ (D9)
We introduce at this point the lattice Fourier transform, specify the choice $`\theta _i=𝐐𝐑_i`$, and write:
$$\underset{i,j}{}a_ib_j\mathrm{cos}𝐐(𝐑_j𝐑_i)J_{i,j}=\frac{1}{2}\underset{𝐤}{\overset{BZ}{}}a(𝐤)b(𝐤)\left(J(𝐤𝐐)+J(𝐤+𝐐)\right),$$
(D10)
where $`a_i`$ and $`b_j`$ is an arbitrary pair of operators and where
$$J(𝐤)=\underset{𝐑}{}\mathrm{exp}(i𝐤𝐑)J(𝐑).$$
(D11)
With the notation
$$\alpha (𝐤)SJ(𝐐)+\frac{S}{2}\left(J(𝐤)+\frac{J(𝐤+𝐐)+J(𝐤𝐐)}{2}\right)$$
(D12)
and
$$\beta (𝐤)\frac{S}{2}\left(J(𝐤)\frac{J(𝐤+𝐐)+J(𝐤𝐐)}{2}\right)$$
(D13)
for any given value of $`𝐐`$, the Hamiltonian (D9) reduces to the form:
$`H`$ $``$ $`{\displaystyle \frac{𝒩S^2}{2}}J(𝐐)+{\displaystyle \underset{𝐤}{\overset{BZ}{}}}[\alpha (𝐤)a^{}(𝐤)a(𝐤)`$ (D14)
$`+`$ $`{\displaystyle \frac{\beta (𝐤)}{2}}(a^{}(𝐤)a^{}(𝐤)+a(𝐤)a(𝐤))].`$ (D15)
The standard Bogoliubov transformation can be used at this point to diagonalize the Hamiltonian (D15). Writing
$$a(𝐤)=u(𝐤)b(𝐤)+v(𝐤)b^{}(𝐤)$$
(D16)
where $`u(𝐤)=\mathrm{cosh}\psi (𝐤)`$ and $`v(𝐤)=\mathrm{sinh}\psi (𝐤)`$ with
$$\mathrm{tanh}2\psi (𝐤)=\frac{\beta (𝐤)}{\alpha (𝐤)}\gamma (𝐤),$$
(D17)
the Hamiltonian (D15) becomes eventually:
$$H=E_0+E_1+\underset{𝐤}{\overset{BZ}{}}\epsilon (𝐤)b^{}(𝐤)b(𝐤).$$
(D18)
In this expression:
$$E_0=\frac{𝒩S^2}{2}J(𝐐),$$
(D19)
$$E_1=\underset{𝐤}{}\left(\alpha (𝐤)v^2(𝐤)+u(𝐤)v(𝐤)\beta (𝐤)\right),$$
(D20)
and
$`\epsilon (𝐤)`$ $`=`$ $`\sqrt{\alpha ^2(𝐤)\beta ^2(𝐤)}`$ (D21)
$`=`$ $`S\sqrt{\left(J(𝐐)J(𝐤)\right)\left(J(𝐐){\displaystyle \frac{J(𝐤+𝐐)+J(𝐤𝐐)}{2}}\right)}.`$ (D22)
The spectrum (D22) coincides (apart for a different normalization of $`J`$) with that given in Ref. (cf. in particular Eq. (67) therein). Note from Eq. (D22) the presence of a Goldstone mode for $`𝐤=0`$ and $`𝐤=\pm 𝐐`$ (owing to the symmetry $`J(𝐤)=J(𝐤)`$).
In the classical limit (that is, for large values of $`S`$), minimization of the ground-state energy is equivalent to finding the minimum of $`J(𝐐)`$. In this way, the appropriate value of $`𝐐`$ is determined. For instance, when second ($`J_2`$) and third ($`J_3`$) nearest-neighbor couplings are included besides the nearest-neighbor antiferromagnetic coupling ($`J_1`$), $`J(𝐐)`$ reads:
$$J(𝐐)=J_1(\mathrm{cos}Q_x+\mathrm{cos}Q_y)+\mathrm{\hspace{0.17em}2}J_2\mathrm{cos}Q_x\mathrm{cos}Q_y+J_3(\mathrm{cos}2Q_x+\mathrm{cos}2Q_y)$$
(D23)
(with the wave vectors measured in units of the inverse of the lattice spacing). In this case, the diagonal spiral configuration is stable for $`J_3>J_1/4J_2/2`$ and $`J_3>J_2/2`$, implying that $`J_3`$ must be nonvanishing for the spiral phase to be stable. The corresponding spin-wave spectrum (D22) for $`J_3=J_2=J_1/5`$ and $`S=1/2`$ is shown in Fig. 4 (full line), where the antiferromagnetic spectrum with $`𝐐=(\pi ,\pi )`$ and $`J_3=J_2=0`$ (broken line) is also shown for comparison.
## E RKKY spin waves for a spiral configuration
In this Appendix, we show how the part of the dispersion relation (106) of the text containing the trigonometric functions in the denominator originates from the Ruderman-Kittel-Kasuya-Yosida (RKKY) interaction between two localized spins and mediated by the conduction electrons. The long-range nature of this interaction is readily evidenced by expanding formally the denominator in Eq. (106) as a a power series in $`\mathrm{cos}q_x`$ and $`\mathrm{cos}q_y`$, so that an infinite number of Heisenberg-like terms appears.
We follow here the treatment by Mattis, and consider the interaction Hamiltonian between the localized spin operators $`𝐒_i`$ (associated with the valence (filled) band) and the itinerant spin operators $`𝐬_c(𝐑_i)`$ (associated with the conduction band), where the suffix $`i`$ specifies the lattice site. We write accordingly:
$$H_{\mathrm{e}xc}=I\underset{i}{}𝐒_i𝐬_c(𝐑_i)$$
(E1)
where $`I`$ is an exchange integral.
For a generic spin-$`1/2`$ operator at site $`i`$ associated with the conduction electrons we write:
$$s_i^\mu =\frac{1}{2𝒩}\underset{𝐤,𝐪}{}\underset{\xi ,\xi ^{}}{}\mathrm{e}^{i𝐪𝐑_i}d_{𝐤+𝐪\xi }^{}\sigma _{\xi ,\xi ^{}}^\mu d_{𝐤\xi ^{}}$$
(E2)
where $`\mu =(+,,z)`$ and $`d_{i\xi }`$ are the destruction operators along the local spin-quantization axes. Combining these operators with the eigenvectors of the mean-field Hubbard Hamiltonian \[cf. Eq. (A12)\] and recalling the notation (A22), we cast Eq. (E2) in the form:
$$s_i^\mu =\frac{1}{2𝒩}\underset{𝐤,𝐪}{}\underset{r,r^{}}{}\mathrm{e}^{i𝐪𝐑_i}\gamma _{𝐤+𝐪,r}^{}F_{r,r^{}}^\mu (𝐤+𝐪,𝐤)\gamma _{𝐤,r^{}}.$$
(E3)
The restriction to the conduction band implies that $`r=r^{}=2`$ in Eq. (E3) when $`\delta >0`$. In this case the Hamiltonian (E1) becomes:
$$H_{\mathrm{e}xc}=\frac{I}{2𝒩}\underset{i}{}\underset{\mu }{}\left(S_i^\mu \right)^{}\left(\underset{𝐤,𝐪}{}e^{i𝐪𝐑_i}\gamma _{𝐤+𝐪,2}^{}F_{22}^\mu (𝐤+𝐪,𝐤)\gamma _{𝐤,2}\right).$$
(E4)
Applying at this point the standard procedure to obtain the energy shift of second-order in the Hamiltonian (E4), the following RKKY-type effective Hamiltonian results:
$$H^{RKKY}=\underset{ij}{}\underset{\mu \nu }{}\left(S_i^\mu \right)^{}J_{ij}^{\mu \nu }S_j^\nu $$
(E5)
with the notation
$$J_{ij}^{\mu \nu }=\frac{I^2}{4𝒩^2}\underset{k<k_F}{}\underset{|𝐤+𝐪|>k_F}{}\mathrm{e}^{i𝐪(𝐑_i𝐑_j)}\frac{F_{22}^\mu (𝐤+𝐪,𝐤)F_{22}^\nu (𝐤+𝐪,𝐤)}{ϵ_2(𝐤)ϵ_2(𝐤+𝐪)}.$$
(E6)
Note that the expression (E6) does not vanish even when sites $`i`$ and $`j`$ are far apart. Note also that the energy denominator on the right-hand side never vanish by construction.
There remains to obtain the spin-wave dispersion associated with the spin Hamiltonian (E5). To this end, we introduce in Eq. (E5) the usual Holstein-Primakoff transformation for the spin operators and truncate the expansion in the bosonic operators $`a`$ and $`a^{}`$ to quadratic order \[cf. Eq. (D6)\], to obtain:
$`H^{RKKY}`$ $`=`$ $`E_0+S{\displaystyle \underset{|𝐪|k_F}{}}[(J^{++}(𝐪)+J^{}(𝐪))a^{}(𝐪)a(𝐪)`$ (E7)
$`+`$ $`J^+(𝐪)a^{}(𝐪)a^{}(𝐪)+J^+(𝐪)a(𝐪)a(𝐪)]`$ (E8)
where now
$$E_0=S^2\underset{|𝐪|k_F}{}J^{zz}(𝐪)+S\underset{|𝐪|k_F}{}J^{}(𝐪)$$
(E9)
with the restriction $`|𝐪|k_F`$ replacing the weaker condition $`|𝐤+𝐪|>k_F`$ (cf. Section 3-B). Note that in Eq. (E8) we have introduced the notation
$$J^{\mu \nu }(𝐪)=\frac{I^2}{4𝒩}\underset{|𝐤|<k_F}{}\frac{F_{22}^\mu (𝐤+𝐪,𝐤)F_{22}^\nu (𝐤+𝐪,𝐤)}{ϵ_2(𝐤)ϵ_2(𝐤+𝐪)}$$
(E10)
and used the identity $`_iJ_{ij}^{\mu \nu }=0`$ for any given $`j`$.
The quadratic Hamiltonian (E8) is then diagonalized by the standard Bogoliubov transformation (cf. Appendix D). The resulting spin-wave spectrum has the form:
$$\omega ^{RKKY}(𝐪)=S\left[\left(J^{++}(𝐪)+J^{}(𝐪)\right)^2\mathrm{\hspace{0.17em}4}J^+(𝐪)J^+(𝐪)\right]^{1/2}.$$
(E11)
To proceed further, we need to specify the quantities $`F_{22}^+`$, $`F_{22}^{}`$, and $`ϵ_2`$ entering the definition (E10), at the relevant order in the small parameter $`t/U`$. Using Eqs. (53)-(56) we obtain for $`\delta >0`$:
$`F_{22}^+(𝐤+𝐪,𝐤)`$ $`=`$ $`i\sqrt{2}\left({\displaystyle \frac{t}{U}}\right)T_o(𝐤+𝐪)+𝒪\left((t/U)^2\right)`$ (E12)
$`F_{22}^{}(𝐤+𝐪,𝐤)`$ $`=`$ $`i\sqrt{2}\left({\displaystyle \frac{t}{U}}\right)T_o(𝐤)+𝒪\left((t/U)^2\right)`$ (E13)
as well as
$$ϵ_2(𝐤)ϵ_2(𝐤+𝐪)=t\left(T_e(𝐤)T_o(𝐤+𝐪)\right)+𝒪\left((t/U)^2\right),$$
(E14)
where we have approximated $`2m_2^{(0)}=1\delta 1`$.
The method developed in Appendix C can be used at this point to perform the sum over $`𝐤`$ in Eq. (E10). The result is:
$$J^{++}(𝐪)I^2\frac{t}{U^2}\delta \frac{\mathrm{sin}^2(Q/2)}{\mathrm{cos}(Q/2)}\frac{(\mathrm{sin}q_x+\mathrm{sin}q_y)^2}{\mathrm{cos}q_x+\mathrm{cos}q_y2},$$
(E15)
$$J^{}(𝐪)I^2\frac{t}{U^2}𝒪(\delta ^2),$$
(E16)
$$J^+(𝐪)=J^+(𝐪)I^2\frac{t}{U^2}𝒪(\delta ^2).$$
(E17)
At the lowest order, only $`J^{++}`$ contributes to Eq. (E11), which then reduces to:
$$\omega ^{RKKY}(𝐪)SI^2\frac{t}{U^2}\delta \frac{\mathrm{sin}^2(Q/2)(\mathrm{sin}q_x+\mathrm{sin}q_y)^2}{|\mathrm{cos}(Q/2)|(2\mathrm{cos}q_x\mathrm{cos}q_y)}.$$
(E18)
The self-consistency equation (67) relating $`Q`$ to $`\delta `$ and $`t/U`$ for the diagonal spiral configuration can eventually be used, to yield (for $`S=1/2)`$:
$$\omega ^{RKKY}(𝐪)J_{\mathrm{e}ff}\frac{(\mathrm{sin}q_x+\mathrm{sin}q_y)^2}{2\mathrm{cos}q_x\mathrm{cos}q_y}$$
(E19)
with $`J_{\mathrm{e}ff}`$ given by Eq. (107) of the text and where we have set $`I=2U`$. Note that Eq. (E19) coincides with Eq. (109) of the text, which was obtained by setting “by hand” $`_{21}=_{12}=0`$ in the full calculation. Note also that the expression (E19) involves transverse spin components only, akin the general expression (98) of the text.
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# Chiral Transition and Baryon-number Susceptibility aafootnote aThe title is slightly changed from the original one “Baryon-number Susceptibility and Related Problems”.
## 1 Introduction
When exploring a phase transition in any physical system, the study of fluctuations of physical quantities, especially ones related to the order parameter is as important as that of the phase diagram for the system in equilibrium. The fluctuations of observables are related with dynamical phenomena including the transport properties of the system. The chirally restored and deconfined phase is expected to be created dynamically in the intermediate stage of the ultra-relativistic heavy ion collisions and in the early universe. Therefore the study of the fluctuations has a great relevance to phenomenology.
In the present report, we discuss the baryon-number susceptibility $`\chi _B`$ $`^{\mathrm{?},\mathrm{?}}`$ and related topics which include the density fluctuations around the critical point of the chiral transition at finite temperature $`T`$ and baryon density $`\rho _B`$.
## 2 QCD phase diagram in $`(T,\rho _B)`$-plane and the vector coupling
The lattice simulations of QCD $`^{\mathrm{?},\mathrm{?}}`$ suggest that the order and even the existence of the phase transition(s) at finite temperature $`T`$ are largely dependent on the number of the active flavors when the physical current quark masses are used: For $`m_um_d10\mathrm{M}\mathrm{e}\mathrm{V}<100\mathrm{M}\mathrm{e}\mathrm{V}\stackrel{<}{}m_s`$, the phase transition may be weak 1st order or 2nd order or not exist. The lattice QCD is, unfortunately, still not matured enough to predict a definite thing about the phase transition at finite baryon density $`\rho _B`$(or chemical potential $`\mu `$).
Low-energy effective models$`^{\mathrm{?},\mathrm{?}}`$ and the chiral random-matrix theory$`^\mathrm{?}`$ have given suggestive pictures of the phase diagram of QCD in the $`T`$-$`\mu `$ (or $`T`$-$`\rho _B`$) plane. For example, the NJL model$`^\mathrm{?}`$ well describes the gross features of the $`T`$ dependence of the quark condensates of the lightest three quarks as given by the lattice QCD, and predict that the chiral transition at $`\mu 0`$ is of rather strong first order at low temperatures $`T`$, say, lower than 50 MeV, provided that the vector coupling between the quarks as given by $`g__V/2(\overline{q}\gamma _\mu q)^2`$ is absent$`^{\mathrm{?},\mathrm{?}}`$. <sup>b</sup><sup>b</sup>b The first order chiral transition in the density direction at low temperatures is also obtained in the chiral random-matrix theory.$`^\mathrm{?}`$ As a matter of fact, the strength and even the existence of the 1st order transition are strongly dependent on the strength of the vector coupling $`g__V`$ $`^{\mathrm{?},\mathrm{?}}`$; the vector term prevents a high-density state.
The reason why the vector coupling weakens the phase transition and postpones the chiral restoration is understood as follows. Thermodynamics tells us that when two phases I and II are coexistent, their temperatures $`T_{\mathrm{I},\mathrm{II}}`$, pressures $`P_{\mathrm{I},\mathrm{II}}`$ and the chemical potentials $`\mu _{\mathrm{I},\mathrm{II}}`$ are the same;
$`T_\mathrm{I}=T_{\mathrm{II}},P_\mathrm{I}=P_{\mathrm{II}},\mu _\mathrm{I}=\mu _{\mathrm{II}}.`$ (2.1)
If the phase I (II) is the chirally broken (chirally restored) phase, the last equality further tells us something because $`\mu _{\mathrm{I},\mathrm{II}}`$ at vanishing temperature are given by $`\mu _\mathrm{i}=\sqrt{M^2+p_{F_i}^2}`$, ($`i=`$ I, II), where $`p_{F_i}`$ is the Fermi momentum of the $`i`$-th phase, and $`M`$ and $`m`$ are the constituent quark mass and the current quark mass that vanishes in the chiral limit. One readily sees that $`p_{F_I}<p_{F_{II}}`$, accordingly, $`\rho _{B_I}<\rho _{B_{II}}`$, i.e., the chirally restored phase is in higher density than the coexistent broken phase. The vector coupling above gives rise to a repulsion proportional to the density squared i.e., $`g_V\rho _B^2/2`$ which is bigger in the restored phase than in the broken phase. Thus the vector coupling weakens and/or postpone the phase transition of the chiral restoration at low temperatures.
Since the vector coupling contribute to the energy repulsively, it also suppress the density fluctuations, or the baryon-number susceptibility, which is the main subject of the present report.
## 3 Baryon-number susceptibility
The baryon-number susceptibility $`\chi _B`$ is the measure of the response of the baryon number density $`\rho _B=_{ı=1N_f}\rho _i`$ to infinitesimal changes in the quark chemical potentials $`\mu _i`$ $`^{\mathrm{?},\mathrm{?}}`$:
$`\chi _B(T,\mu )=\left[{\displaystyle \underset{i=1}{\overset{N_f}{}}}{\displaystyle \frac{}{\mu _i}}\right]({\displaystyle \underset{i=1}{\overset{N_f}{}}}\rho _i)=N_B^2/VT,`$ (3.1)
where $`N_B`$ is the baryon-number operator given by $`N_B_{i=1}^{N_f}N_i,`$ with
$`\rho _i=\mathrm{Tr}N_i\mathrm{exp}[\beta (H{\displaystyle \underset{i=u,d}{}}\mu _iN_i)]/VN_i/V`$ (3.2)
the $`i`$-th quark-number density, $`V`$ the volume of the system and $`\beta =1/T`$.
In the following, we shall confine ourselves to the $`SU_f(N_f)`$-symmetric case; $`\mu _u=\mu _d=\mu _s=\mathrm{}\mu .`$
The baryon-number susceptibility at $`\rho _B0`$ is related with the (iso-thermal) compressibility of the system $`^\mathrm{?}`$ $`\kappa __TN_B^1(V/\mu )_{T,N_B}=\frac{\chi _B}{\rho ^2}`$, which tells that if $`\chi _B`$ is large and so is the density fluctuation, the system is easy to compress.
Another physical meaning of $`\chi _B`$ is that it is the density-density correlation which is nothing but the 0-0 component of the vector-vector correlations or fluctuations$`^\mathrm{?}`$;
$`\chi _B(T,\mu _q)=\beta {\displaystyle 𝑑𝐱S_{00}(0,𝐱)},`$ (3.3)
where $`S_{\mu \nu }(t,𝐱)=j_\mu (t,𝐱)j_\nu (0,\mathrm{𝟎}),`$ with $`j_\mu (t,𝐱)=\overline{q}(t,𝐱)\gamma _\mu q(t,𝐱)`$ being the current operator. Using the fluctuation-dissipation theorem, one has
$`\chi _B(T,\mu _q)=\underset{k0}{lim}L(0,𝐤),`$ (3.4)
where $`L(\omega ,𝐤)`$ is the longitudinal component of the retarded Green’s function or the response function in the vector channel;
$`R_{\mu \nu }(\omega ,k)=\mathrm{F}.\mathrm{T}.(i\theta (t)[j_\mu (t,𝐱),j_\nu (0,\mathrm{𝟎})]_{})`$.
The lattice simulations of QCD $`^\mathrm{?}`$ show that $`\chi _B`$ at $`\mu =0`$ is suppressed in the low temperature phase, increases with $`T`$ sharply around the critical point of the chiral transition and takes almost the free quark gas value then saturates. This behavior may be understood intuitively and roughly as follows$`^\mathrm{?}`$: In the confined phase at low $`T`$, the density fluctuation picks up the Boltzmann factor $`\mathrm{e}^{M_N/T}`$ with $`M_N`$ being the nucleon mass, which is much smaller than the factor $`\mathrm{e}^{M_q/T}`$ with $`M_q`$ being the current quark mass (constituent quark mass), which factor will be picked up in the deconfined and chirally restored (chirally broken) phase.
Our point here is however that the nature of the chiral transition and also presence (or absence) of the vector coupling affect the baryon-number susceptibility, especially when $`\rho _B0`$, for which the lattice data is not available so far$`^\mathrm{?}`$.
### 3.1 Free quark gas
It is instructive to examine $`\chi _B(T,\mu )`$ in the simple free-quark gas model:
$`\rho _B=2N_fN_c{\displaystyle \frac{d𝐩}{(2\pi )^3}(n(T,\mu )\overline{n}(T,\mu ))}`$ (3.5)
where $`n(T,\mu )=1/[\mathrm{exp}\beta (E_p\mu )+\mathrm{\hspace{0.17em}1}]`$ and $`\overline{n}(T,\mu )=1/[\mathrm{exp}\beta (E_p+\mu )+\mathrm{\hspace{0.17em}1}]`$ with $`E_p=\sqrt{M^2+p^2}`$, and $`N_c=3`$ is the number of the colors. Then one readily obtains
$`\chi _B(T,\mu )=2N_fN_c\beta {\displaystyle \frac{d𝐩}{(2\pi )^3}\left\{n(1n)+\overline{n}(1\overline{n})\right\}}\chi _B^{(0)}(T,\mu ),`$ (3.6)
which is reduced to
$`\chi _B^{(0)}(T,0)\chi _B^{(0)}(T)=4N_fN_c\beta {\displaystyle \frac{d𝐩}{(2\pi )^3}\frac{\mathrm{exp}(E_p/T)}{[\mathrm{exp}(E_p/T)+1]^2}},`$ (3.7)
at $`\mu =0`$.
If $`M(T)`$ is decreased as in the chiral restoration, $`\chi _B`$ increases and reaches $`N_fT^2`$ at $`M(T)=0`$: The enhancement is,however, found to be modest and not so large as obtained in the lattice simulations.
### 3.2 Model calculation
To demonstrate the relevance of the nature of the chiral transition and the presence of the vector coupling to $`\chi _B`$, we perform a calculation with an effective model, which is given by adding the vector-coupling terms $`^\mathrm{?}`$ to the Nambu$``$Jona-Lasinio model$`^\mathrm{?}`$:
$``$ $`=`$ $`\overline{q}(i\gamma 𝐦)q+{\displaystyle \underset{a=0}{\overset{N_f^21}{}}}{\displaystyle \frac{g__S}{2}}[(\overline{q}\lambda _aq)^2+(\overline{q}i\lambda _a\gamma _5q)^2]`$ (3.8)
$`{\displaystyle \frac{g__V}{2}}{\displaystyle \underset{a=0}{\overset{N_f^21}{}}}[(\overline{q}\lambda _a\gamma _\mu q)^2+(\overline{q}\lambda _a\gamma _\mu \gamma _5q)^2]`$
The realistic value of $`g__V`$ used in the literature $`^\mathrm{?}`$ is roughly in the range of $`g__V\mathrm{\Lambda }^2=59`$.
In the self-consistent mean field approximation, the constituent mass (dynamical mass) $`M_i(T,\mu _q)`$, the quark condensates $`\overline{q}_iq_i`$ and the quark density $`\rho _i`$ are all coupled with the vector coupling $`g__V`$ and determined by the following equations;
$`M_i`$ $`=`$ $`m_i2g__S\overline{q}_iq_i,`$ (3.9)
$`\overline{q}_iq_i`$ $`=`$ $`2N_c{\displaystyle \frac{d𝐩}{(2\pi )^3}\{1n_i(T,\stackrel{~}{\mu }_i)\overline{n}_i(T,\stackrel{~}{\mu }_i)\}},`$ (3.10)
$`\rho _i`$ $`=`$ $`2N_c{\displaystyle \frac{d𝐩}{(2\pi )^3}(n_i(T,\stackrel{~}{\mu })\overline{n}_i(T,\stackrel{~}{\mu }))},`$ (3.11)
where it is to be noted that the shifted chemical potential $`\stackrel{~}{\mu }=\mu 2g__V\rho _i`$ enters the distribution functions instead of the naive one, $`\mu `$.
Simply differentiating these equations with respect to $`\mu `$, one obtains $`\chi _B(T,\mu )`$. It is noteworthy that when $`\mu 0`$ there arises a coupling between $`\chi _B`$ and the scalar-density susceptibility $`\chi _s`$ owing to the non-vanishing “vector-scalar susceptibility” $`\chi _{_{VS}}`$. They are defined by
$`\chi _s={\displaystyle \frac{d\overline{q}q}{dm}}=\beta {\displaystyle 𝑑𝐱\overline{q}(0,𝐱)q(0,𝐱)\overline{q}(0,\mathrm{𝟎})q(0,\mathrm{𝟎})},`$ (3.12)
$`\chi _{_{VS}}={\displaystyle \frac{\overline{q}q}{\mu _B}}=\beta {\displaystyle 𝑑𝐱\overline{q}(0,𝐱)\gamma _0q(0,𝐱)\overline{q}(0,\mathrm{𝟎})q(0,\mathrm{𝟎})},`$ (3.13)
respectively. $`\chi _s`$ represents the fluctuation of the order parameter of the chiral transition, and is related with the sigma meson propagator. The differentiation leads to the coupled equation
$`\left(\begin{array}{cc}1+2g__V\chi _B^{(0)}& \chi _{_{VS}}^{(0)}\\ 4g_sg__V\chi _{_{VS}}^{(0)}& 12g_s\chi _s^{(0)}\end{array}\right)\left(\begin{array}{c}\chi _B\\ 2g_s\chi _{_{VS}}\end{array}\right)=\left(\begin{array}{c}\chi _B^{(0)}\\ 2g_s\chi _{_{VS}}^{(0)}\end{array}\right)`$ (3.20)
where $`\chi _B^{(0)}(T,\mu _q)`$ is the zero-th order baryon-number susceptibility given before,
$`\chi _{VS}^{(0)}=2N_c\beta {\displaystyle \underset{i=1}{\overset{N_f}{}}}{\displaystyle \frac{d𝐩}{(2\pi )^3}\left\{n_i(1n_i)\overline{n}_i(1\overline{n}_i)\right\}}`$ (3.21)
the zero-th order vector-scalar one and
$`\chi _s^{(0)}`$ $`=`$ $`2N_c{\displaystyle \underset{i=1}{\overset{N_f}{}}}{\displaystyle }{\displaystyle \frac{d𝐩}{(2\pi )^3}}[{\displaystyle \frac{p^2}{E_i^3}}(1n_i\overline{n}_i)`$
$`+\beta {\displaystyle \frac{M_i^2}{E_i^2}}\{n_i(1n_i)+\overline{n}_i(1\overline{n}_i)\}]`$ (3.22)
the scalar-density susceptibility in the zero-th order. One should note that when $`\mu =0`$, then $`\chi _{VS}^{(0)}`$ vanishes. One can readily obtain $`\chi _q`$ and $`\chi _s`$ in terms of $`\chi ^{(0)},\chi _{VS}^{(0)}`$ and $`\chi _s^{(0)}`$, so we do not write them down to save the space.
### 3.3 (A) Finite density case $`g_V=0`$
To see how the nature of the chiral transition can affect the behavior of $`\chi _B`$, let us first take a simple case where $`g_V`$ is negligible.
The essential point lies in the fact that the distribution function $`n(T,\mu )=[\mathrm{e}^{\beta (E_p\mu )}+1]^1`$ depends on $`\mu `$ not only explicitly but also implicitly through the dynamical mass $`M(T,\mu )`$ in $`E_p=\sqrt{M^2+p^2}`$; hence
$`T{\displaystyle \frac{n}{\mu }}=n(1n){\displaystyle \frac{M}{E_p}}{\displaystyle \frac{M}{\mu }}n(1n).`$ (3.23)
Thus
$`\chi _B(T,\mu )`$ $`=`$ $`2N_c{\displaystyle \underset{i=1}{\overset{N_f}{}}}\beta {\displaystyle }{\displaystyle \frac{d𝐩}{(2\pi )^3}}[\{n_i(1n_i)+\overline{n}_i(1\overline{n}_i)\}`$
$`{\displaystyle \frac{M_i}{E_{ip}}}{\displaystyle \frac{M_i}{\mu }}\{n_i(1n_i)\overline{n}_i(1\overline{n}_i)\}].`$ (3.24)
The notable point is the presence of the derivative $`\frac{M}{\mu }`$. The constituent mass $`M`$ varies with the quark condensate $`\overline{q}q`$ by Eq. (3.9), the order parameter of the chiral transition. We saw in §2 that the chiral transition at low temperatures is likely to be of first order in the chemical potential direction. It means that the derivative $`\frac{M}{\mu }`$ diverges at the critical point at low temperatures, hence so does the susceptibility $`\chi _B`$.
### 3.4 (B) Zero-density case with $`g_V0`$
Putting $`\mu =0`$ into the expressions one gets$`^\mathrm{?}`$
$`\chi _B={\displaystyle \frac{\chi _B^{(0)}(T)}{1+2g__V\chi _B^{(0)}(T)}},`$ (3.25)
where $`\chi _B^{(0)}(T)`$ is the susceptibility for the free-quark gas. The denominator of $`\chi _B`$ is essentially the inverse of the propagator of the vector meson in the ring approximation at the vanishing four momenta. The above expression shows that $`\chi _B`$ is suppressed by the vector coupling ( $`g__V>0`$). This is reasonable at least for a system with a finite $`\mu `$; because the system becomes hard to compress when the vector coupling is present, the number fluctuations will be suppressed<sup>c</sup><sup>c</sup>c The interactions due to vector fields like $`\omega `$ or neutral $`\rho `$ mesons increases the repulsion between the constituents of the system. Recall also that $`\chi _B`$ is proportional to the compressibility $`\kappa _T`$.
The comparison with the lattice data$`^\mathrm{?}`$ shows that the vector coupling is rather small in the high temperature phase. It is interesting that this suppression of the vector coupling at high temperatures is consistent with the observation that the screening masses of the vector modes obtained in the lattice simulations $`^\mathrm{?}`$ almost coincides with $`2\pi T`$, the lowest screening mass of the q-$`\overline{\mathrm{q}}`$ system in the chiral limit; a similar result is also obtained in the instanton approach$`^\mathrm{?}`$.
## 4 Estimate of $`g__V`$ and $`g__S`$ in the lattice QCD
Boyd et al $`^\mathrm{?}`$ once extracted the effective coupling constants $`g__V`$ and $`g__S`$ from the lattice data using the expressions given in the NJL model as given above; $`\chi _B=\chi _B^{(0)}/(1+g__V\chi _B^{(0)})`$ and $`\chi _S=\chi _S^{(0)}/(1g__S\chi _S^{(0)})`$. They concluded that when $`TT_C`$, $`g__S`$ is much bigger than $`g__V`$; $`g__S4g__V`$. This result is consistent with our analysis and the behavior of the screening masses.
## 5 Implications to phenomenology
Kumagai, Miyamura and Sugitate$`^\mathrm{?}`$ discussed the implications of the baryon number susceptibility and the strangeness susceptibility to the observables in the relativistic heavy-ion collisions. They argued that in the stopping region where the chemical potential is large, large fluctuations of the baryon and the strangeness numbers may be a signature of the chiral transition.
Here we wish to also indicate that the large number fluctuations cause those in the scalar channel (the sigma meson channel) at finite density.
The observability of the possible large density fluctuations caused as a critical phenomenon is discussed by several authors$`^\mathrm{?}`$.
## 6 Summary and concluding remarks
We have examined the baryon-number susceptibility $`\chi _B`$ as an observable which reflects the confinement-deconfinement and the chiral phase transitions in hot and/or dense hadronic matter.
The suppression of $`\chi _B`$ at low temperatures and steep rise around the critical temperature as shown in the lattice QCD may be roughly attributed to the confinement-deconfinement transition. Nevertheless, we have shown that such a behavior of $`\chi _B`$ is also affected by the chiral transition.
Since $`\chi _B`$ is a measure of the rate of the density fluctuation in the system, the chiral transition at finite chemical potential especially leads to an interesting phenomenological consequence to $`\chi _B`$. When the vector coupling is small, the chiral transition at low temperatures is of first order in the density direction, which implies a divergent behavior of $`\chi _B`$, accordingly a huge density fluctuations. We have emphasized that such a large enhancement of the fluctuation can be also expected for the scalar density fluctuations due to the scalar-vector mixing at finite density. Such a large enhancement may leads to an enhancement of the sigma-meson production <sup>d</sup><sup>d</sup>dAs for the significance of the sigma meson for the chiral transition at finite $`T`$ and/or $`\rho _B`$, see $`^{\mathrm{?},\mathrm{?}}`$.. The above phenomena all have relevance to experiments to be done in RHIC and LHC.
We have indicated that the nature of the chiral transition as to the first order or not etc is sensitively dependent on the strength of the vector coupling. An analysis of the lattice data suggests that the vector coupling is small in comparison with the scalar coupling at high temperature.
The susceptibility $`\chi _B`$ is nothing but the generalized susceptibility $`\chi (\omega ,k)`$ at $`\omega =k=0`$. One should examine $`\chi (\omega ,k)`$ in the whole region of $`\omega `$ and $`k`$ to get more information about the vector correlations and the density fluctuations.
In conclusion, we would like to thank the organizers of this workshop, especially, Prof. H. Suganuma, for inviting me to talk on the baryon-number susceptibility.
## References
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# Contents
## 1 Introduction
As well-known the relation between geometry and (1+1)-dimensional solitonic equations is described (the geometrical part) in terms of two-dimensional surfaces (and/or curves that is equivalent in 1+1) in a 3-dimensional $`M^3`$ (Euclidean or pseudo-Euclidean) spaces \[1-24\]. In 2+1 dimensions the situation is more complicated \[25-30, 36-38, 40\]. Usually, in this case uses the same approach, i.e., two-dimensional surfaces arbitrarily embedded in 3-dimensional space $`M^3`$ but considering also the motion (deformation) of surfaces. Of course, no doubt that it is correct but in some sense is the artificial construction. Sometimes we think that in 2+1 dimensions may be the more natural geometry (at least than the above mentioned surfaces approach) is 3-dimensional Riemann spaces (see, e.g., refs.\[26-29, 46-48\]). This line we realized in our notes \[49-51\].
This note is a sequel to the preceding notes \[49-51\]. Our main goal in this note is to discuss the interrelation between the intrinsic geometry of three-dimensional Riemann spaces and integrable systems, in particular, spin systems in 2+1 dimensions without entering into details (for details see e.g. refs.\[49-51\]). We will do this without any reference to the enveloping spaces (see, also, ref. ). We will show that three-dimensional Riemann spaces is important not only in the general relativity but is also important in the theory of multidimensional integrable systems. As well known, the theories of multidimensional integrable systems and three-dimensional Riemann spaces are not as complete as their counterparts, respectively, integrable systems in 1+1 dimensions and two-dimensional surfaces. So, the study of the connections between integrable systems in 2+1 dimensions and three-dimensional Riemann space is one of actual problems of modern mathematical physics.
## 2 Three-dimensional Riemann space
Let $`V^3`$ be the space endowed with the affine connection. In this space we introduce twoo systems of coordinates: $`(x)=(x^1,x^2,x^3)`$ and $`(y)=(y^1,y^2,y^3)`$. It is well known from the classical differential geometry that these coordinate systems are connected by the following set of equations of second order (remark: for convenience, in \[49-51\] and this note we will use the unified common numerations for formulas)
$$\frac{^2y_k}{x_ix_j}=\mathrm{\Gamma }_{ij}^l(x)\frac{y_k}{x_l}\mathrm{\Gamma }_{lm}^k(y)\frac{y_ly_m}{x_ix_j}$$
$`(155)`$
We mention that in this case the curvature tensor
$$R_{klm}^i=[\frac{\mathrm{\Gamma }_m}{x^l}\frac{\mathrm{\Gamma }_l}{x^n}+\mathrm{\Gamma }_l\mathrm{\Gamma }_m\mathrm{\Gamma }_m\mathrm{\Gamma }_l]_k^i$$
$`(156)`$
has only three independent components. Let the space $`V^3`$ is flat and let the metric tensor in $`V^3`$ in the coordinate system $`(y)`$ is diagonal:
$$ds^2=\underset{i,j=1}{\overset{3}{}}\mu _{ij}dy^idy^j$$
$`(157)`$
with $`\mu _{ij}=\pm 1`$. Hence follows that the cooresponding connection and Riemann’s curvature tensor are equal to zero:
$$\mathrm{\Gamma }_{jk}^i(y)=0,R_{klm}^i(y)=0$$
$`(158)`$
Let for the coordinate system $`(x)`$ the metric has the form
$$ds^2=\underset{i,j=1}{\overset{3}{}}g_{ij}dx_idx_j.$$
$`(159)`$
As the curvature tensor has the law of transformation as the four rank tensor
$$R_{klm}^i(x)=\frac{x^i}{y^s}\frac{y^r}{x^k}\frac{y^q}{x^l}\frac{y^p}{x^n}R_{klm}^i(y)$$
$`(160)`$
the curvature tensor for the coordinate system $`(x)`$ is also equal to zero
$$R_{klm}^i(x)=0.$$
$`(161)`$
Hence, for the coordinate system $`(x)`$ we get the following system of three equations
$$\frac{\mathrm{\Gamma }_m}{x^l}\frac{\mathrm{\Gamma }_l}{x^n}+\mathrm{\Gamma }_l\mathrm{\Gamma }_m\mathrm{\Gamma }_m\mathrm{\Gamma }_l=0$$
$`(162)`$
where $`\mathrm{\Gamma }_m(x)`$ are matrices with components
$$\mathrm{\Gamma }_{ij}^k=\frac{1}{2}g^{kl}(g_{il,j}+g_{jl,i}g_{ij,l}).$$
$`(163)`$
We note that in our case the scalar curvature is equal to zero
$$R=\underset{i,k,l,m=1}{\overset{3}{}}g^{il}g^{km}R_{iklm}=0.$$
$`(164)`$
Now the system (155) takes the form
$$\frac{^2y^k}{x^ix^j}=\mathrm{\Gamma }_{ij}^l(x)\frac{y^k}{x^l}.$$
$`(165)`$
Let $`𝐫=(y^1,y^2,y^3)=𝐫(x^1,x^2,x^3)`$ is the position vector and put $`x^1=x,x^2=y,x^3=t`$. Then as follows from (165) the position vector $`𝐫`$ satisfies the following set of equations
$$𝐫_{xx}=\mathrm{\Gamma }_{11}^1𝐫_x+\mathrm{\Gamma }_{11}^2𝐫_y+\mathrm{\Gamma }_{11}^3𝐫_t$$
$`(166a)`$
$$𝐫_{xy}=\mathrm{\Gamma }_{12}^1𝐫_x+\mathrm{\Gamma }_{12}^2𝐫_y+\mathrm{\Gamma }_{12}^3𝐫_t$$
$`(166b)`$
$$𝐫_{xt}=\mathrm{\Gamma }_{13}^1𝐫_x+\mathrm{\Gamma }_{13}^2𝐫_y+\mathrm{\Gamma }_{13}^3𝐫_t$$
$`(166c)`$
$$𝐫_{yy}=\mathrm{\Gamma }_{22}^1𝐫_x+\mathrm{\Gamma }_{22}^2𝐫_y+\mathrm{\Gamma }_{22}^3𝐫_t$$
$`(166d)`$
$$𝐫_{yt}=\mathrm{\Gamma }_{23}^1𝐫_x+\mathrm{\Gamma }_{23}^2𝐫_y+\mathrm{\Gamma }_{23}^3𝐫_t$$
$`(166e)`$
$$𝐫_{tt}=\mathrm{\Gamma }_{33}^1𝐫_x+\mathrm{\Gamma }_{33}^2𝐫_y+\mathrm{\Gamma }_{33}^3𝐫_t.$$
$`(166f)`$
We can rewrite the equation (166) in the following form
$$Z_x=A_1Z,Z_y=A_2Z,Z_t=A_3Z$$
$`(167)`$
where
$$Z=(𝐫_x,𝐫_y,𝐫_t)^T$$
$`(168)`$
and
$$A_1=\left(\begin{array}{ccc}\mathrm{\Gamma }_{11}^1& \mathrm{\Gamma }_{11}^2& \mathrm{\Gamma }_{11}^3\\ \mathrm{\Gamma }_{12}^1& \mathrm{\Gamma }_{12}^2& \mathrm{\Gamma }_{12}^3\\ \mathrm{\Gamma }_{13}^1& \mathrm{\Gamma }_{13}^2& \mathrm{\Gamma }_{13}^3\end{array}\right)A_2=\left(\begin{array}{ccc}\mathrm{\Gamma }_{12}^1& \mathrm{\Gamma }_{12}^2& \mathrm{\Gamma }_{12}^3\\ \mathrm{\Gamma }_{22}^1& \mathrm{\Gamma }_{22}^2& \mathrm{\Gamma }_{22}^3\\ \mathrm{\Gamma }_{23}^1& \mathrm{\Gamma }_{23}^2& \mathrm{\Gamma }_{23}^3\end{array}\right)$$
$$A_3=\left(\begin{array}{ccc}\mathrm{\Gamma }_{13}^1& \mathrm{\Gamma }_{13}^2& \mathrm{\Gamma }_{13}^3\\ \mathrm{\Gamma }_{23}^1& \mathrm{\Gamma }_{23}^2& \mathrm{\Gamma }_{23}^3\\ \mathrm{\Gamma }_{33}^1& \mathrm{\Gamma }_{33}^2& \mathrm{\Gamma }_{33}^3\end{array}\right).$$
$`(169)`$
The compatibility condition of these equations are given by
$$\frac{A_i}{x^j}\frac{A_j}{x^i}+[A_i,A_j]=0.$$
$`(170)`$
For our further work it is convenient use the triad of unit vectors. We introduce these vectors by the way
$$𝐞_1=\frac{𝐫_x}{H_1},𝐞_2=\frac{𝐫_y}{H_2}+c_1𝐫_x+c_2𝐫_t,𝐞_3=𝐞_1𝐞_2.$$
$`(171)`$
The explicit forms of $`c_1,c_2`$ given in and
$$H_1=𝐫_x,H_2=𝐫_y,H_3=𝐫_t.$$
$`(172)`$
The equations (156) for $`𝐞_k`$ take the form
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_x=B_1\left(\begin{array}{ccc}𝐞_1& & \\ 𝐞_2& & \\ 𝐞_3& & \end{array}\right),\left(\begin{array}{ccc}𝐞_1& & \\ 𝐞_2& & \\ 𝐞_3& & \end{array}\right)_y=B_2\left(\begin{array}{ccc}𝐞_1& & \\ 𝐞_2& & \\ 𝐞_3& & \end{array}\right)$$
$$\left(\begin{array}{ccc}𝐞_1& & \\ 𝐞_2& & \\ 𝐞_3& & \end{array}\right)_t=B_3\left(\begin{array}{ccc}𝐞_1& & \\ 𝐞_2& & \\ 𝐞_3& & \end{array}\right)$$
$`(173)`$
with
$$B_1=\left(\begin{array}{ccc}0& k& \sigma \\ \beta k& 0& \tau \\ \beta \sigma & \tau & 0\end{array}\right),B_2=\left(\begin{array}{ccc}0& m_3& m_2\\ \beta m_3& 0& m_1\\ \beta m_2& m_1& 0\end{array}\right)$$
$$B_3=\left(\begin{array}{ccc}0& \omega _3& \omega _2\\ \beta \omega _3& 0& \omega _1\\ \beta \omega _2& \omega _1& 0\end{array}\right)$$
$`(174)`$
where $`k,\tau ,\sigma ,m_i,\omega _i`$ are some real functions the explicit forms of which given in \[?\], $`\beta =𝐞_1^2=\pm 1,𝐞_2^2=𝐞_3^2=1`$. Again from the integrability condition of these equations we obtain the following set of equations
$$\frac{B_i}{x^j}\frac{B_j}{x^i}+[B_i,B_j]=0.$$
$`(175)`$
Many integrable systems in 2+1 dimensions are exact reductions of the equation (175) (see, e.g., ).
## 3 Integrable reductions
Now to find out particular integrable reductions of three-dimensional Riemann space, as in \[49-51\] , we will use multidimensional integrable spin systems (MISSs). To this end, we assume that
$$𝐞_1𝐒$$
$`(176)`$
where $`𝐒=(S_1,S_2,S_3)`$ is the spin vector, $`𝐒^2=\beta =\pm 1`$. So, the vector $`𝐞_1`$ satisfies the some given MISS. Some comments on MISSs are in order. At present there exist several MISSs (see, e.g., Appendix). They play important role both in mathematics and physics. In this note, to find out a integrable case of three-dimensional Riemann space we will use the MISS - the Ishimori equation (IE).
### 3.1 The Ishimori equation
The IE has the form
$$𝐒_t=𝐒(𝐒_{xx}+\alpha ^2S_{yy})+u_y𝐒_x+u_x𝐒_y$$
$`(177a)`$
$$u_{xx}\alpha ^2u_{yy}=2\alpha ^2𝐒(𝐒_x𝐒_y).$$
$`(177b)`$
The IE is integrable by Inverse Scattering Transform (IST) . In this case, we get
$$m_1=_x^1[\tau _y\frac{1}{2\alpha ^2}M_2^{}u],m_2=\frac{1}{2\alpha ^2k}M_2^{}u$$
$$m_3=_x^1[k_y+\frac{\tau }{2\alpha ^2k}M_2u],M_2^{}u=\alpha ^2u_{yy}u_{xx}$$
$`(178)`$
and
$$\omega _2=(k_x+\sigma \tau )\alpha ^2(m_{3y}+m_2m_1)+m_2u_x+\sigma u_y$$
$$\omega _3=(\sigma _xk\tau )+\alpha ^2(m_{2y}m_3m_1)+ku_y+m_3u_x$$
$$\omega _1=\frac{1}{k}[\sigma _t\omega _{2x}+\tau \omega _3].$$
$`(179)`$
Thus we expressed the functions $`m_k`$ and $`\omega _k`$ by the three functions $`k,\tau ,\sigma `$ and their derivatives. This means that we identified the equations (161) and (184) which define the geometry of Riemann space with the given MISS - the IE (177). In turn it means that we given off the integrable case of the three-dimensional Riemann space. Hence arises the natural question: how construct the integrable three-dimensional Riemann space using the MISS or that the same thing how find $`g_{ij}`$?. The answer as follows. Let
$$𝐫_x^2=H_1^2=\beta =\pm 1.$$
$`(180)`$
Then we have that
$$𝐫=_x^1𝐒+𝐫_0(y,t)$$
$`(181)`$
where $`_x^1=_{\mathrm{}}^x𝑑x`$. For simplicity, we put $`𝐫_0=0`$. Now we can express the coefficients of the metric (143) by $`𝐒`$. As shown in , for the (2+1)-dimensional case, we have
$$g_{11}=𝐒^2=\pm 1,g_{12}=𝐒_x^1𝐒_y,g_{13}=𝐒_x^1𝐒_t$$
$$g_{22}=(_x^1𝐒_y)^2,g_{23}=(_x^1𝐒_y)(_x^1𝐒_t),g_{33}=(_x^1𝐒_t)^2.$$
$`(182)`$
### 3.2 The Davey-Stewartson equation
It is well known from the soliton theory that between integrable spin systems and NLS-type equations take place the so-called gauge and/or L-equivalences . From this fact and from the identification the three-dimensional Riemann space and the MISS (177) in the previous subsection follows that there exist some connections with NLS-type equations. We show it in this subsection. Let us we introduce two complex functions $`q,p`$ as
$$q=a_1e^{ib_1},p=a_2e^{ib_2}$$
$`(183)`$
where $`a_j,b_j`$ are real functions. Let $`a_k,b_k`$ have the form
$$a_1^2=\frac{|a|^2}{|b|^2}\{\frac{k^2}{4}+\frac{|\alpha |^2}{4}(m_3^2+m_2^2)\frac{1}{2}\alpha _Rkm_3\frac{1}{2}\alpha _Ikm_2\}$$
$`(184a)`$
$$b_1=_x^1\{\frac{\gamma _1}{2ia_1^^2}(\overline{A}A+D\overline{D})\}$$
$`(184b)`$
$$a_2^2=\frac{|b|^2}{|a|^2}\{\frac{k^2}{4}+\frac{|\alpha |^2}{4}(m_3^2+m_2^2)+\frac{1}{2}\alpha _Rkm_3\frac{1}{2}\alpha _Ikm_2\}$$
$`(184c)`$
$$b_2=_x^1\{\frac{\gamma _2}{2ia_2^^2}(A\overline{A}+\overline{D}D)$$
$`(184d)`$
where
$$\gamma _1=i\{\frac{1}{2}k^2\tau +\frac{|\alpha |^2}{2}(m_3km_1+m_2k_y)$$
$$\frac{1}{2}\alpha _R[k^2m_1+m_3k\tau +m_2k_x]+\frac{1}{2}\alpha _I[k(2k_ym_{3x})k_xm_3]\}$$
$`(185a)`$
$$\gamma _2=i\{\frac{1}{2}k^2\tau +\frac{|\alpha |^2}{2}(m_3km_1+m_2k_y)+$$
$$\frac{1}{2}\alpha _R(k^2m_1+m_3k\tau +m_2k_x)+\frac{1}{2}\alpha _I[k(2k_ym_{3x})k_xm_3]\}.$$
$`(185b)`$
Here $`\alpha =\alpha _R+i\alpha _I`$. In this case, $`q,p`$ satisfy the DS equation \[?\]
$$iq_t+q_{xx}+\alpha ^2q_{yy}+vq=0$$
$`(186a)`$
$$ip_tp_{xx}\alpha ^2p_{yy}vp=0$$
$`(186b)`$
$$\alpha ^2v_{yy}v_{xx}=2[\alpha ^2(pq)_{yy}+(pq)_{xx}].$$
$`(186c)`$
## 4 Diagonal metrics
Now we consider the case when the metric has the diagonal form, i.e.
$$ds^2=ϵ_1H_1^2dx^2+ϵ_2H_2^2dy^2+ϵ_3H_3^2dt^2$$
$`(187)`$
where $`ϵ_i=\pm 1.`$ In this note we consider the case when $`ϵ_i=+1`$. In this case, the Christoffel symbols take the form
$$\mathrm{\Gamma }_{11}^1=\frac{H_{1x}}{H_1}=\beta _{11},\mathrm{\Gamma }_{11}^2=\frac{H_1H_{1y}}{H_2^2}=\frac{H_1}{H_2}\beta _{21},\mathrm{\Gamma }_{11}^3=\frac{H_1H_{1t}}{H_3^2}=\frac{H_1}{H_3}\beta _{31}$$
$$\mathrm{\Gamma }_{12}^1=\frac{H_{1y}}{H_1}=\frac{H_2}{H_1}\beta _{21},\mathrm{\Gamma }_{12}^2=\frac{H_{2x}}{H_2}=\frac{H_1}{H_2}\beta _{12},\mathrm{\Gamma }_{12}^3=0$$
$$\mathrm{\Gamma }_{13}^1=\frac{H_{1t}}{H_1}=\frac{H_3}{H_1}\beta _{31},\mathrm{\Gamma }_{13}^2=0,\mathrm{\Gamma }_{13}^3=\frac{H_{3x}}{H_3}=\frac{H_1}{H_3}\beta _{13}$$
$$\mathrm{\Gamma }_{22}^1=\frac{H_2H_{2x}}{H_1^2}=\frac{H_2}{H_1}\beta _{12},\mathrm{\Gamma }_{22}^2=\frac{H_{2y}}{H_2}=\beta _{22},\mathrm{\Gamma }_{22}^3=\frac{H_2H_{2t}}{H_3^2}=\frac{H_2}{H_3}\beta _{32}$$
$$\mathrm{\Gamma }_{23}^1=0,\mathrm{\Gamma }_{23}^2=\frac{H_{2t}}{H_2}=\frac{H_3}{H_2}\beta _{32},\mathrm{\Gamma }_{23}^3=\frac{H_{3y}}{H_3}=\frac{H_2}{H_3}\beta _{23}$$
$$\mathrm{\Gamma }_{33}^1=\frac{H_3H_{3x}}{H_1^2}=\frac{H_3}{H_1}\beta _{13},\mathrm{\Gamma }_{33}^2=\frac{H_3H_{3y}}{H_2^2}=\frac{H_3}{H_2}\beta _{23},\mathrm{\Gamma }_{33}^3=\frac{H_{3t}}{H_3}=\beta _{33}$$
or
$$\mathrm{\Gamma }_{ij}^k=0ijk$$
$`(188a)`$
$$\mathrm{\Gamma }_{il}^i=\frac{H_{i,l}}{H_i}=\frac{H_l}{H_i}\beta _{li}$$
$`(188b)`$
$$\mathrm{\Gamma }_{ll}^i=\frac{H_lH_{li}}{H_i^2}=\frac{H_l}{H_i}\beta _{il},il.$$
$`(188c)`$
Here
$$\beta _{ik}=\frac{H_{k,i}}{H_i}$$
$`(189)`$
are the so-called rotation coefficients. In this case we get
$$\tau =m_2=\omega _3=0$$
$`(190)`$
and the matrices $`B_i`$ take the form
$$B_1=\left(\begin{array}{ccc}0& \beta _{21}& \beta _{31}\\ \beta _{21}& 0& 0\\ \beta _{31}& 0& 0\end{array}\right),B_2=\left(\begin{array}{ccc}0& \beta _{12}& 0\\ \beta _{12}& 0& \beta _{32}\\ 0& \beta _{32}& 0\end{array}\right)$$
$$B_3=\left(\begin{array}{ccc}0& 0& \beta _{13}\\ 0& 0& \beta _{23}\\ \beta _{13}& \beta _{23}& 0\end{array}\right).$$
$`(191)`$
So we have
$$\beta _{23x}=\beta _{13}\beta _{21},\beta _{32x}=\beta _{12}\beta _{31}$$
$`(192a)`$
$$\beta _{13y}=\beta _{12}\beta _{23},\beta _{31y}=\beta _{32}\beta _{21}$$
$`(192b)`$
$$\beta _{12t}=\beta _{13}\beta _{32},\beta _{21t}=\beta _{23}\beta _{31}$$
$`(192c)`$
$$\beta _{12x}+\beta _{21y}+\beta _{31}\beta _{32}=0$$
$`(192d)`$
$$\beta _{13x}+\beta _{31t}+\beta _{21}\beta _{23}=0$$
$`(192e)`$
$$\beta _{23y}+\beta _{32t}+\beta _{12}\beta _{13}=0$$
$`(192f)`$
or
$$\frac{\beta _{ij}}{x^k}=\beta _{ik}\beta _{kj},ijk$$
$`(192g)`$
$$\frac{\beta _{ij}}{x^i}+\frac{\beta _{ji}}{x^j}+\underset{mi,j}{\overset{3}{}}\beta _{mi}\beta _{mj}=0,ij$$
$`(192h)`$
This nonlinear system is the famous Lame equation and well-known in the theory of 3-orthogonal coordinates. The problem of description of curvilinear orthogonal coordinate systems in a a (pseudo-) Euclidean space is a classical problem of differential geometry. It was studied in detail and mainly solved in the beginning of the 20th century. Locally, such coordinate systems are determinated by $`\frac{n(n1)}{2}`$ arbitrary functions of two variables. This problem in some sense is equivalent to the problem of description of diagonal flat metrics, that is, flat metrics $`g_{ij}(x)=f_i(x)\delta _{ij}`$. We mention that the Lame equation describing curvilinear orthogonal coordinate systems can be integrated by the IST (see also an algebraic-geometric approach in ). Now we would like consider some particular cases (see, also, e.g., ).
### 4.1 $`H_1=H_2=H,H_3=1`$
Let $`H=e^\psi `$. We get the following set of equations
$$\psi _{xx}+\psi _{yy}+\psi _t^2e^{2\psi }=0$$
$`(193a)`$
$$\psi _{tx}=\psi _{ty}=0,\psi _{tt}+\psi _t^2=0.$$
$`(193b)`$
From (159) we obtain $`\psi _t=\frac{1}{Ct}`$. So Eq. (193) reduced to the equation
$$\psi _{xx}+\psi _{yy}+\frac{1}{Ct}e^{2\psi }=0.$$
$`(194)`$
### 4.2 $`H_1=\mathrm{cos}\theta ,H_2=\mathrm{sin}\theta ,H_3=1`$
In this case the corresponding equation has the form
$$(\theta _t\mathrm{cos}\theta )_x=\theta _x\theta _t\mathrm{sin}\theta ,(\theta _t\mathrm{sin}\theta )_y=\theta _y\theta _t\mathrm{cos}\theta $$
$`(195a)`$
$$\theta _{xt}=\theta _{yt}=0,\theta _{xx}\theta _{yy}\theta _t^2\mathrm{sin}\theta \mathrm{cos}\theta =0$$
$`(195b)`$
$$(\theta _t\mathrm{sin}\theta )_t=(\theta _t\mathrm{cos}\theta )_t=0.$$
$`(195c)`$
### 4.3 $`H_1=\mathrm{cos}\theta ,H_2=\mathrm{sin}\theta ,H_3=\theta _t`$
In this case the corresponding equation looks like
$$(\frac{\theta _{ty}}{\mathrm{sin}\theta })_x=\frac{\theta _{xt}\theta _y}{\mathrm{cos}\theta }$$
$`(196a)`$
$$(\frac{\theta _{xt}}{\mathrm{cos}\theta })_y=\frac{\theta _x\theta _{ty}}{\mathrm{sin}\theta }$$
$`(196b)`$
$$\theta _{xx}\theta _{yy}\mathrm{sin}\theta \mathrm{cos}\theta =0$$
$`(196c)`$
$$(\frac{\theta _{xt}}{\mathrm{cos}\theta })_y(\mathrm{sin}\theta )_t\frac{\theta _y\theta _{ty}}{\mathrm{sin}\theta }=0$$
$`(196d)`$
$$(\frac{\theta _{ty}}{\mathrm{sin}\theta })_y+(\mathrm{cos}\theta )_t+\frac{\theta _x\theta _{xt}}{\mathrm{cos}\theta }=0.$$
$`(196e)`$
### 4.4 $`H_1=e^\psi ,H_2=e^\psi ,H_3=\psi _t`$
In this case we have
$$(\psi _{ty}e^\psi )_x=\psi _{tx}\psi _ye^\psi ,(\psi _{tx}e^\psi )_y=\psi _{ty}\psi _xe^\psi $$
$`(197a)`$
$$\psi _{xx}+\psi _{yy}+e^{2\psi }=0$$
$`(197b)`$
$$(\psi _{tx}e^\psi )_x+\psi _te^\psi +\psi _y\psi _{ty}e^\psi =0,(\psi _{ty}e^\psi )_y+\psi _te^\psi +\psi _x\psi _{tx}e^\psi =0.$$
$`(197c)`$
### 4.5 $`H_1=H_2=H_3=H^2`$
In this choose we get ($`H=e^\psi `$)
$$\psi _{xy}=\psi _x\psi _y,\psi _{xt}=\psi _x\psi _t,\psi _{yt}=\psi _t\psi _y$$
$`(198a)`$
$$\psi _{xx}+\psi _{yy}+4\psi _t^2=0,\psi _{xx}+\psi _{tt}+4\psi _y^2=0,\psi _{tt}+\psi _{yy}+4\psi _x^2=0.$$
$`(198b)`$
## 5 Connections with the other equations
It is remarkable that the equation (175) \[=(162)=(170)=(213)\] is related with the some well-known equations. In this section we present some of these connections.
### 5.1 Equation (175) and the Bogomolny equation
Consider the Bogomolny equation (BE)
$$\mathrm{\Phi }_t+[\mathrm{\Phi },B_3]+B_{1y}B_{2x}+[B_1,B_2]=0$$
$`(199a)`$
$$\mathrm{\Phi }_y+[\mathrm{\Phi },B_2]+B_{3x}B_{1t}+[B_3,B_1]=0$$
$`(199b)`$
$$\mathrm{\Phi }_x+[\mathrm{\Phi },B_1]+B_{2t}B_{3y}+[B_2,B_3]=0.$$
$`(199c)`$
This equation is integrable and play important role in the field theories in particular in the theory of monopols. The set of equations (175) is the particular case of the BE. In fact, as $`\mathrm{\Phi }=0`$ from (199) we obtain the system (175).
### 5.2 Equation (175) and the Self-Dual Yang-Mills equation
Equation (175) is exact reduction of the SO(3)-Self-Dual Yang-Mills equation (SDYME)
$$F_{\alpha \beta }=0,F_{\overline{\alpha }\overline{\beta }}=0,F_{\alpha \overline{\alpha }}+F_{\beta \overline{\beta }}=0$$
$`(200)`$
Here
$$F_{\mu \nu }=\frac{A_\nu }{x_\mu }\frac{A_\mu }{x_\nu }+[A_\mu ,A_\nu ]$$
$`(201)`$
and
$$\frac{}{x_\alpha }=\frac{}{z}i\frac{}{t},\frac{}{x_{\overline{\alpha }}}=\frac{}{z}+i\frac{}{t},\frac{}{x_\beta }=\frac{}{x}i\frac{}{y}$$
$$\frac{}{x_{\overline{\beta }}}=\frac{}{x}+i\frac{}{y}.$$
$`(202)`$
In fact, if in the SDYME (200) we take
$$A_\alpha =iB_3,A_{\overline{\alpha }}=iB_3,A_\beta =B_1iB_2,A_\beta =B_1+iB_2$$
$`(203)`$
and if $`B_k`$ are independent of $`z`$, then the SDYME (200) reduces to the equation (175). As known that the LR of the SDYME has the form
$$(_\alpha +\lambda _{\overline{\beta }})\mathrm{\Psi }=(A_\alpha +\lambda A_{\overline{\beta }})\mathrm{\Psi },(_\beta \lambda _{\overline{\alpha }})\mathrm{\Psi }=(A_\beta \lambda A_{\overline{\alpha }})\mathrm{\Psi }$$
$`(204)`$
where $`\lambda `$ is the spectral parameter satisfing the following set of the equations
$$\lambda _\beta =\lambda \lambda _{\overline{\alpha }},\lambda _\alpha =\lambda \lambda _{\overline{\beta }}.$$
$`(205)`$
Apropos, the simplest solution of this set has may be the following form
$$\lambda =\frac{a_1x_{\overline{\alpha }}+a_2x_{\overline{\beta }}+a_3}{a_2x_\alpha a_1x_\beta +a_4},a_j=consts.$$
$`(206)`$
From (204) we obtain the LR of the equation (175)
$$(i_t+\lambda _{\overline{\beta }})\mathrm{\Psi }=[iB_3+\lambda (B_1+iB_2)]\mathrm{\Psi }$$
$`(207a)`$
$$(_\beta i\lambda _t)\mathrm{\Psi }=[(B_1iB_2)i\lambda B_3]\mathrm{\Psi }.$$
$`(207b)`$
### 5.3 Equation (175) and the Chern-Simons equation
Consider the action of the Chern-Simons (CS) theory
$$S[J]=\frac{k}{4\pi }_Mtr(JdJ+\frac{2}{3}JJJ)$$
$`(208)`$
where $`J`$ is a 1-form gauge connection with values in the Lie algebra $`\widehat{g}`$ of a (compact or noncompact) non-Abelian simple Lie group $`\widehat{G}`$ on an oriented closed 3-dimensional manifold $`M`$, $`k`$ is the coupling constant. The classical equation of motion is the zero-curvature condition
$$dJ+JJ=0.$$
$`(209)`$
Let the 1-form $`J`$ has the form
$$J=B_1dx+B_2dy+B_3dt.$$
$`(210)`$
As shown in , subtituting the (210) into (209) we obtain the equation (175). Note that from this fact and from the results of the subsection 5.2 follows that the CS - equation of motion (209) is exact reduction of the SDYM equation (200).
### 5.4 Equation (175) as some generalization of the Lame equation
Let us the matrices $`B_i`$ (174) we rewrite in the form
$$B_1=\left(\begin{array}{ccc}0& \beta _{21}& \beta _{31}\\ \beta _{21}& 0& \tau \\ \beta _{31}& \tau & 0\end{array}\right),B_2=\left(\begin{array}{ccc}0& \beta _{12}& m_2\\ \beta _{12}& 0& \beta _{32}\\ m_2& \beta _{32}& 0\end{array}\right)$$
$$B_3=\left(\begin{array}{ccc}0& \omega _3& \beta _{13}\\ \omega _3& 0& \beta _{23}\\ \beta _{13}& \beta _{23}& 0\end{array}\right).$$
$`(211)`$
Then the equation (175) in elements takes the form
$$\beta _{23x}\tau _t=\beta _{13}\beta _{21}\omega _3\beta _{31}$$
$`(212a)`$
$$\beta _{32x}+\tau _y=\beta _{12}\beta _{31}+m_2\beta _{21}$$
$`(212b)`$
$$\beta _{13y}+m_{2t}=\beta _{12}\beta _{23}+\omega _3\beta _{32}$$
$`(212c)`$
$$\beta _{31y}m_{2x}=\beta _{32}\beta _{21}\tau \beta _{12}$$
$`(212d)`$
$$\beta _{12t}\omega _{3y}=\beta _{13}\beta _{32}m_3\beta _{23}$$
$`(212e)`$
$$\beta _{21t}+\omega _{3x}=\beta _{23}\beta _{31}+\tau \beta _{13}$$
$`(212f)`$
$$\beta _{12x}+\beta _{21y}+\beta _{31}\beta _{32}+\tau m_2=0$$
$`(212g)`$
$$\beta _{13x}+\beta _{31t}+\beta _{21}\beta _{23}+\tau \omega _3=0$$
$`(212h)`$
$$\beta _{23y}+\beta _{32t}+\beta _{12}\beta _{13}+m_2\omega _3=0$$
$`(212i)`$
Hence as $`\tau =m_2=\omega _3=0`$ we obtain the Lame equation (192). So, the equation (175) is one of the generalizations of the Lame equation.
## 6 On Lax representation of the equation (175)
As follows from the results of the section 3, Equation (175) can admits several integrable reductions. At the same time, the results of the subsections 5.1-5.2 show that the equation (175) is integrable may be and in general case. At least, it admits the LR of the form (207) and/or of the following form (see, e.g., )
$$\mathrm{\Phi }_x=U_1\mathrm{\Phi },\mathrm{\Phi }_y=U_2\mathrm{\Phi },\mathrm{\Phi }_t=U_3\mathrm{\Phi }$$
$`(213)`$
with
$$U_1=\frac{1}{2}\left(\begin{array}{cc}i\tau & (k+i\sigma )\\ ki\sigma & i\tau \end{array}\right),U_2=\frac{1}{2}\left(\begin{array}{cc}im_1& (m_3+im_2)\\ m_3im_2& im_1\end{array}\right)$$
$$U_3=\frac{1}{2}\left(\begin{array}{cc}i\omega _1& (\omega _3+i\omega _2)\\ \omega _3i\omega _2& i\omega _1\end{array}\right).$$
$`(214)`$
Systems of this type were first studied by Zakharov and Shabat . The integrability conditions on this system of overdetermined equations (211), require that
$$U_{i,j}U_{j,i}+[U_i,U_j]=0.$$
$`(215)`$
Many (and perhaps all) integrable systems in 2+1 dimensions have the LR of the form (211). In our case, the IE (177) and the DS equation (186) have also the LR of the form (211) with the functions $`m_i,\omega _i`$ given by (123) and (124). On the other hand, it is well-known that for example the DS equation (186) has the following LR of the standard form
$$\alpha \mathrm{\Psi }_y=\sigma _3\mathrm{\Psi }_x+Q\mathrm{\Psi },Q=\left(\begin{array}{cc}0& q\\ p& 0\end{array}\right)$$
$`(216a)`$
$$\mathrm{\Psi }_t=2i\sigma _3\mathrm{\Psi }_{xx}+2iQ\mathrm{\Psi }_x+\left(\begin{array}{cc}c_{11}& iq_x+i\alpha q_y\\ ip_xi\alpha p_y& c_{22}\end{array}\right)\mathrm{\Psi }$$
$`(216b)`$
with
$$c_{11x}\alpha c_{11y}=i[(pq)_x+\alpha (pq)_y],c_{22x}+\alpha c_{22y}=i[(pq)_x\alpha (pq)_y].$$
$`(217)`$
Hence arises the natural question: how connected the both LR for one and the same integrable systems (in our case for the DS equation)?. In fact, these two LR are related by the gauge transformation \[50-51\]
$$\mathrm{\Phi }=g\mathrm{\Psi }$$
$`(218)`$
where $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are some solutions of the equations (222) and (111), respectively, while $`g`$ is the some matrix.
## 7 Conclusion
In this note we have studied the relation between integrable systems in 2+1 dimensions and 3-dimensional Riemann spaces. We have shown that in this geometrical setting certain typical structures of the completely integrable (2+1)-dimensional systems arise. To find out the integrable cases of the 3-dimensional Riemann space as an examples we used the Ishimori and Davey-Stewartson equations. The connections of the equations characterizing of the 3-dimensional Riemann space with the other known equations such as the Bogomolny and Self-Dual Yang-Mills equations are considered. Such connection with the Chern-Simons equation of motion was established in \[?\]. Of course our approach needs further developments. Finally we note that the details of some calculations were given in \[50-51\].
## 8 Acknowledgments
This work was partially supported by INTAS (grant 99-1782). RM would like to thanks to V.S.Dryuma, M.Gurses, B.G. Konopelchenko, D.Levi, L.Martina and G.Soliani for very helpful discussions and especially D.Levi for the financial support and kind hospitality. He is grateful to the EINSTEIN Consortium of Lecce University and the Department of Mathematics of Bilkent University for their financial supports and warm hospitality.
## 9 Appendix: MISSs
At present there exist several MISSs. Here we will give some of them.
i) The Myrzakulov I (M-I) equation The simplest example of MISS is the M-I (Remark: about our notations please see e.g., ref. ) equation looks like
$$𝐒_t=(𝐒𝐒_y+u𝐒)_x$$
$`(219a)`$
$$u_x=𝐒(𝐒_x𝐒_y)$$
$`(219b)`$
ii) The Myrzakulov VIII (M-VIII) equation. The M-VIII equation is one of simplest MISSs in 2+1 dimensions and reads as
$$𝐒_t=𝐒𝐒_{xx}+u𝐒_x$$
$`(220a)`$
$$u_x+u_y+𝐒(𝐒_x𝐒_y)=0$$
$`(220b)`$
iii) The Ishimori equation. The famous Ishimori equation has the form
$$𝐒_t=𝐒(𝐒_{xx}+\alpha ^2𝐒_{yy})+u_x𝐒_y+u_y𝐒_x$$
$`(221a)`$
$$u_{xx}\alpha ^2u_{yy}=2\alpha ^2𝐒(𝐒_x𝐒_y)$$
$`(221b)`$
iv) The Myrzakulov IX (M-IX) equation. This equation reads as
$$𝐒_t=𝐒M_1𝐒iA_1𝐒_yiA_2𝐒_x$$
$`(222a)`$
$$M_2u=2\alpha ^2𝐒(𝐒_x𝐒_y)$$
$`(222b)`$
Here $`M_i,A_i`$ have the forms
$$M_1=\alpha ^2\frac{^2}{y^2}+4\alpha (ba)\frac{^2}{xy}+4(a^22abb)\frac{^2}{x^2}$$
$$M_2=\alpha ^2\frac{^2}{y^2}2\alpha (2a+1)\frac{^2}{xy}+4a(a+1)\frac{^2}{x^2}$$
$$A_1=i\{\alpha (2b+1)u_y2(2ab+a+b)u_x\}$$
$$A_2=i\{4\alpha ^1(2a^2b+a^2+2ab+b)u_x2(2ab+a+b)u_y\}$$
The M-IX equation contains several particular integrable cases: a) the M-VIII equation (218) as $`a=b=1`$; b) the Ishimori equation (177) as $`a=b=\frac{1}{2}`$; c) the M-XXXIV equation as $`a=b=1,y=t`$
$$𝐒_t=𝐒𝐒_{xx}+u𝐒_x$$
$`(223a)`$
$$u_t+u_x+\frac{1}{2}(𝐒_x^2)^2=0$$
$`(223b)`$
and so on. The M-XXXIV equation (221) describe the nonlinear dynamics of compressible magnets . It is the first (and, to the best of our knowledge, at present the unique) example of integrable spin system governing the nonlinear interactions of spin ($`𝐒)`$ and lattice $`(u)`$ subsystems in 1+1 dimensions.
v) The Myrzakulov XX (M-XX) equation. This equation reads as
$$𝐒_t=𝐒\{(b+1)𝐒_{xx}b𝐒_{yy}\}+bu_y𝐒_y+(b+1)u_x𝐒_x$$
$`(224a)`$
$$u_{xy}=𝐒(𝐒_x𝐒_y)$$
$`(224b)`$
vi) The (2+1)-dimensional Myrzakulov 0 (M-0) equation. The (2+1)-dimensional M-0 equation
$$𝐒_t=c_1𝐒_x+c_2𝐒_y$$
$`(223)`$
is in general not integrable but admits integrable reductions. For example, the following case of the (2+1)-dimensional M-0 equation is integrable
$$𝐒_t=\frac{\beta _{13}}{\beta _{31}}𝐒_x\frac{\beta _{13}}{\beta _{12}\beta _{31}}𝐒_y$$
$`(226a)`$
$$\frac{\beta _{ij}}{x^k}=\beta _{ik}\beta _{kj}$$
$`(226b)`$
$$\frac{\beta _{ij}}{x^i}+\frac{\beta _{ji}}{x^j}+\underset{mi,j}{\overset{3}{}}\beta _{mi}\beta _{mj}=0$$
$`(226c)`$
and so on. All of these MISSs are some (2+1)-dimensional integrable extensions of the isotropic Landau-Lifshitz (LL) equation
$$𝐒_t=𝐒𝐒_{xx}$$
$`(227)`$
and in 1+1 dimensions reduced to it. Here we would like mention that there exist the other classes MISSs which are not multidimensional generalizations of the LL equation (227), for example , the M-II, M-III and M-XXII equations and so on. Finally we note that all MISSs in 2+1 dimensions are the integrable particular cases of the M-0 equation (225).
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# Mapping Integer Order Neumann Functions To Real Orders
## Abstract
In a recent paper we unified Bessel functions of different orders .Here we extend the unification to other linairely independant solutions to Bessel equation ,Neumann’s and Hankel’s functions
mekhfi@hotmail.com
Neumann’s functions $`N_p(z)`$ and Bessel functions $`J_p(z)`$ with p integer ,are the two linairely independant solutions to the Bessel equation .
$$z^2y^{^{\prime \prime }}(z)+zy^{^{}}(z)+(z^2p^2)y(z)=0$$
(1)
When p is real it is known that $`J_p(z)`$ and $`J_p(z)`$ are also two linairely independant solutions of the above equation so consequently we can rewrite $`N_p(z)`$ as a combination of
Bessel functions of positive and negative orders and the combination is as follows
$$N_p(z)=\frac{J_p(z)cosp\pi J_p(z)}{sinp\pi }$$
(2)
Bessel and Neumann’s functions are the ”analog” respectively of the cosine and sine ,while Hankel functions $`H_p^{(1)}`$ and $`H_p^{(2)}`$ yet another linear combinations of $`J_p(z)`$ and $`J_p(z)`$ are the ”analog” of $`e^{ipz}`$ and $`e^{ipz}`$. Hankel functions are related to $`N_p(z)`$ and $`J_p(z)`$ as.
$`J_p(z)`$ $`=`$ $`{\displaystyle \frac{H_p^{(1)}(z)+H_p^{(2)}(z)}{2}}`$
$`N_p(z)`$ $`=`$ $`{\displaystyle \frac{H_p^{(1)}(z)H_p^{(2)}(z)}{2i}}`$ (3)
where $`H^{(1)}`$ and $`H^{(2)}`$ are complex conjugate to each other.
In a recent paper we showed, for the first time , that Bessel functions of real orders $`J_{n+\lambda }`$ $`0\lambda 1`$ and Bessel functions $`J_n`$ $`nZ`$ of integer orders are no longer disassociated objects , but one is the deformed<sup>1</sup><sup>1</sup>1The deformation mechanism is described in version of the other and as a consequence ,they are related to each other through the unifying formula <sup>2</sup><sup>2</sup>2We keep the notation of the partial derivative $``$ wherever and when it acts on a single variable it is identified with a simple derivative $`d.`$
$`{\displaystyle \frac{J_{n+\lambda }(z)}{z^{n+\lambda }}}`$ $`=`$ $`exp(\lambda {\displaystyle \underset{z_1}{}}){\displaystyle \frac{J_n(z)}{z^n}}`$ (4)
$`z^{n+\lambda }J_{n+\lambda }(z)`$ $`=`$ $`exp(\lambda {\displaystyle \underset{z_2}{}})z^nJ_n(z)`$ (5)
$`{\displaystyle \underset{z_1}{}}`$ $`=`$ $`{\displaystyle \underset{mϵZ/0}{}}{\displaystyle \frac{_m}{m}}`$
$`{\displaystyle \underset{z_2}{}}`$ $`=`$ $`{\displaystyle \underset{mϵZ/0}{}}(1)^m{\displaystyle \frac{_m}{m}}`$
$`_{|m|}`$ $`=`$ $`{\displaystyle \frac{2}{dz^2}}{\displaystyle \frac{2}{dz^2}}\mathrm{}{\displaystyle \frac{2}{dz^2}}.`$
$`_{|m|}`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{1}{2}}dz^2.{\displaystyle }{\displaystyle \frac{1}{2}}dz^2\mathrm{}\mathrm{}\mathrm{}\mathrm{}..{\displaystyle }{\displaystyle \frac{1}{2}}dz^2.`$
At this point it seems natural to extend this formula which applies to $`J_n,`$ to the second linairely independant solution $`N_n`$.The apparently straigthforward way to do would be to use formula 4 5 and then apply the mapping operator $`e^{\pm \lambda {\scriptscriptstyle }}`$ , but unfortunately our formula applies to the reduced Bessel function $`J_n(z)/z^n`$ rather than directly to $`J_n(z)`$ and this make our formula of no practical use for our purpose.In the following we will propose a nice way to relate $`N_{n+\lambda }`$ to $`N_n`$ .
Let us recall first that to Bessel functions one can associate polynomials $`O_n(z).`$ These are Neumann polynomials .The relationship between them is performed through the known formula.
$`{\displaystyle \frac{1}{tz}}`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}J_n(z)A_n(t)`$ (6)
$`withO_n(z)`$ $`=`$ $`A_n+(1)^nA_n`$
The first step towards the solution of our problem is to show that if , to Bessel functions one can associate Neumann polynomials $`O_n(z),`$ one can equally well associate to reduced Bessel functions something and these associated objects are precisely Neumann function $`N_n`$ $`(z)`$ . We are not sure if a formula which relates reduced Bessel functions to Neumann functions ,similar to the above one 6, already exists in the literature ,so we prefer to work it out first.In doing this we largely follow a method which dates back to Sonine and adapts it to the present need.Let us first state the result.Let $`\mathrm{\Omega }(\tau )`$ be an an arbitrary function of $`\tau `$ and if $`\mathrm{\Omega }(\tau )=x`$ ,let $`\tau =\mathrm{}(x)`$so that $`\mathrm{}`$ is the inverse function to $`\mathrm{\Omega }`$ .Denote the function and the function we associate to it respectively $`Z_n(z)`$ and $`A_n(t)`$ and define them through their integral representations as follows
$`Z_n(z)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle e^{z^2\mathrm{\Omega }+\frac{\tau }{2}}\frac{d\tau }{\tau ^{n+1}}}`$
$`A_n(t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}e^{t^2x\frac{\mathrm{}(x)}{2}}(\mathrm{}(x))^n𝑑x`$
then it is not that difficult to prove the existence of a formula relating both functions
$`{\displaystyle \frac{1}{t^2z^2}}`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}Z_n(z)A_n(t)`$ (7)
$`Rez^2`$ $`>`$ $`Ret^2`$
To prove the result 7 , suppose for any given positive value of x that $`\tau >\mathrm{}(x)`$ on a closed curve $`𝒞`$ surrounding the origin and the point z and that $`\tau <\mathrm{}(x)`$ on a closed $`c`$ surrounding the origin but not enclosing the point z .Then compute the series
$`{\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}Z_nA_n`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}{\displaystyle _𝒞}{\displaystyle _0^{\mathrm{}}}exp(z^2\mathrm{\Omega }(\tau )+t^2x+{\displaystyle \frac{1}{2}}(\tau \mathrm{}(x)){\displaystyle \frac{\mathrm{}(x)^n}{\tau ^{n+1}}}dxd\tau `$
$`+`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}{\displaystyle _c}{\displaystyle _0^{\mathrm{}}}exp(z^2\mathrm{\Omega }(\tau )+t^2x+{\displaystyle \frac{1}{2}}(\tau \mathrm{}(x)){\displaystyle \frac{\tau ^n}{\mathrm{}(x)^{n+1}}}dxd\tau `$
$`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _0^{\mathrm{}}}({\displaystyle _𝒞}{\displaystyle _c}){\displaystyle \frac{exp(z^2\mathrm{\Omega }(\tau )+t^2x+\frac{1}{2}(\tau \mathrm{}(x)}{\tau \mathrm{}(x)}}𝑑\tau 𝑑x.`$
$`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _z}{\displaystyle \frac{exp(z^2\mathrm{\Omega }(\tau )+t^2x+\frac{1}{2}(\tau \mathrm{}(x)}{\tau \mathrm{}(x)}}𝑑\tau 𝑑x.`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}e^{(t^2z^2)x}𝑑x`$
$`=`$ $`{\displaystyle \frac{1}{t^2z^2}}`$
In going to the third line,we perform the summations $`\frac{1}{\tau }_{n=0}^+\mathrm{}\left(\frac{\mathrm{}(x)}{\tau }\right)^n=\frac{1}{\tau \mathrm{}(x)}`$ ($`\tau `$ on the path $`𝒞`$ )and $`\frac{1}{\mathrm{}}_{n=0}^+\mathrm{}\left(\frac{\tau }{\mathrm{}(x)}\right)^n=\frac{1}{\tau \mathrm{}(x)}`$ ( $`\tau `$ on the path $`c`$ ) which are convergent according to the conditions above .For the case of interest we have $`Z_n(z)=\frac{J_n(z)}{z^n}`$ , $`\mathrm{\Omega }(\tau )=\frac{1}{2\tau }`$ and $`\mathrm{}(x)=\frac{1}{2x}`$ .In this case the function associated to the reduced Bessel function has the integral form .
$$A_n(t)=\frac{1}{2}_0^{\mathrm{}}e^{t^2\frac{x}{2}\frac{1}{2x}}\frac{dx}{x^n}$$
(8)
The integral defining the function $`A_n(t)`$ is easily identified to the $`K_n(t)`$ function which is the Hankel function of imaginary argument $`K_\nu (t)=\frac{\pi i}{2}e^{i\frac{\pi }{2}\nu i}H_\nu ^{(1)}(it)`$.
$`{\displaystyle \frac{1}{2}}{\displaystyle _0^{\mathrm{}}}e^{\frac{1}{2x}+\frac{t^2x}{2}}{\displaystyle \frac{dx}{x^n}}`$ $`=`$ $`{\displaystyle \frac{K_{n+1}(it)}{(it)^{n+1}}}={\displaystyle \frac{K_{n+1}(it)}{(it)^{n+1}}}`$
$`with`$ $``$ $`argt<\pi /2andRet^2>0`$
Equation 7 then relates reduced Bessel funnctions to Hankel functions which are associated to them
$$\frac{2}{\pi }\frac{i}{t^2z^2}=\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\frac{J_n(z)}{z^n}t^{n1}H_{n1}^{(1)}(t).$$
(9)
In writing this formula we have used the property of the Hankel functions that $`H_n=(1)^nH_n`$ .
At this point if we restrict ourselves to t and z reals and take the complex conjugate of the above formula we get another formula
$$\frac{2}{\pi }\frac{i}{t^2z^2}=\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\frac{J_n(z)}{z^n}t^{n1}H_{n1}^{(2)}(t).$$
(10)
where, this time, we use the property $`H_\nu ^{(1)}(z)^{}=H_\nu ^{(2)}(z^{})`$ . Adding 9 to 10 and using the definition of Neumann function in term of Hankel functions 3 we get the desired result
$$\frac{2}{\pi }\frac{1}{t^2z^2}=\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\frac{J_n(z)}{z^n}t^{n1}N_{n1}(t).$$
(11)
The particular form of the left hand side of this equation (a happy event) and the presence of reduced Bessel functions on the right hand side of the equation allow a straigthforward and quite elegant application to Neumann ’s function of the unifying formula 4 ( valid for reduced Bessel functions ) .In fact we can perform a succession of valid operations .
$`{\displaystyle \frac{2}{\pi }}exp(\lambda {\displaystyle \underset{z_1}{}}){\displaystyle \frac{1}{t^2z^2}}`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}exp(\lambda {\displaystyle \underset{z_1}{}}){\displaystyle \frac{J_n(z)}{z^n}}t^{n1}N_{n1}(t)`$ (12)
$`{\displaystyle \frac{2}{\pi }}exp(\lambda {\displaystyle \underset{t}{}}){\displaystyle \frac{1}{t^2z^2}}`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{J_{n+\lambda }(z)}{z^{n+\lambda }}}t^{n1}N_{n1}(t)`$
$`{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{1}{t^2z^2}}`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{J_{n+\lambda }(z)}{z^{n+\lambda }}}exp(\lambda {\displaystyle \underset{t}{}})t^{n1}N_{n1}(t)`$
$`{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{1}{t^2z^2}}`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{J_{n+\lambda }(z)}{z^{n+\lambda }}}t^{n+\lambda 1}N_{n+\lambda 1}(t)`$
Where the capital sigma operator with the t variable acquires an extra factor $`(1)^m`$ .
$$\underset{t}{}=\underset{mϵZ/0}{}(1)^m\frac{_m}{m}$$
(13)
This extra factor comes from the identity.
$$\frac{}{z^2}\frac{1}{t^2z^2}=\frac{}{t^2}\frac{1}{t^2z^2}$$
(14)
Note that in the last line of 12 we identified the result of the action of the operator $`exp(\lambda _t)`$ on Neumann functions of integer orders with Neumann function of orders $`n+\lambda `$ as for reduced Bessel functions.That is we put
$$t^{n+\lambda }N_{n+\lambda }(t)=exp(\lambda \underset{t}{})t^nN_n(t)$$
(15)
It remains to show 15, that is ,the result of the mapping operation by the operator $`e^\lambda `$ is indeed Neumann functions of real orders $`n+\lambda `$ and not something else.To see this let us proceed as follows
$`{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{1}{t^2z^2}}`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{J_{n+\lambda }(z)}{z^{n+\lambda }}}A_{n,\lambda }={\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{J_{n+\lambda m}(z)}{z^{n+\lambda m}}}A_{n,\lambda m}`$
$`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{J_{n+\lambda }(z)}{z^{n+\lambda }}}A_{n+m,\lambda m}`$
In the first line we use the fact that the sum is independant of $`\lambda `$ and hence change the parameter $`\lambda \lambda m`$ .In the second line we absorb the index m as we sum over all $`mϵZ`$ .From Mapping Integer Order Neumann Functions To Real Orders we infer the symmetry property of the $`A_{n,\lambda }`$
$$A_{n,\lambda }=A_{n+m,\lambda m}\lambda ,m$$
(17)
This symmetry means that $`A_{n,\lambda }`$ only depends on the combination $`n+\lambda `$ that is $`A_{n,\lambda }=A_{n+\lambda }`$ .This remark together with the fact that $`A_{n+\lambda }_{\lambda =0}=`$ $`t^nN_n(t)`$ leads to the identification 15as we anticipated in the last line of 12
Formula 15 is the main result of the paper .It shows that Neumann functions are unified irrespective of their orders being integer or real .They are unified through the same formula as reduced Bessel functions $`z^nJ_n(z)`$ in 5.
Note that we could have already applied the procedure to Hankel function’s first by starting from the identities 9 and .10. In this case we could have obtained a similar unifying formulas ( i=1,2 ).
$$t^{n+\lambda }H_{n+\lambda }^{(i)}(t)=exp(\lambda \underset{t}{})t^nH_n^{(i)}(t)$$
(18)
Let us note that our formula 18 for the evolution of Hankel functions could be used to infer the evolution equation for reduced Bessel function $`z^nJ_n`$ by using the defining equation 3 of J’<sup>s</sup> in terms of H’<sup>s</sup> . By simple inspection one can see that the result we get is indeed equation 5.
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# Rotation, Statistical Dynamics and Kinematics of Globular Clusters
## 1. Introduction
When, over ten years ago, the main focus of my work moved away from stellar dynamics, three outstanding problems were left unsolved and my efforts to interest students in these have not borne fruit. In the hope that others may be stimulated these questions are posed here.
Q1. Is it possible to have self-similar core collapse when there is a continuous distribution of stellar masses?
Q2. For a weakly rotating cluster is there a self-similar collapse of the core and how does its rotation evolve as the core radius, $`r_c`$, decreases?
Of course question 2 can be tackled for equal mass stars and then one can make it more realistic by combining both questions together in Q1 + Q2.
Q3. It is now widely accepted that just as in stars the energy generated by nuclear reactions delays thermal evolution, so in star clusters the energy generated by binary stars delays core-collapse. Eddington’s theory of the Main Sequence led to many observational consequences. Can we make a theory of the Dynamical Main Sequence of globular clusters with real predictive power as to the form these clusters should take at equilibrium?
Section 2 outlines possible solutions but full solutions involve multidimensional problems in four or five dimensions (e.g., Energy, Angular Momentum, Time and Mass).
Some years ago Inagaki and I (1990) showed that such problems could be tackled analytically using trial functions and a Local Variational Principle. I still think that method should have a great future in these problems but those wishing to use it should contact Takahashi (1992,3,5,6) who developed this approach numerically for 3 dimensional problems, but has since turned to another method.
Since few workers in N-body dynamics have even heard of Newton’s exact solution of an N-body problem for all N and all initial conditions, it is given in section 3 along with recent developments of this idea. These problems lead to some pretty dynamics which illustrate that there are exceptions with no violent relaxation and that thermodynamic equilibria can exist in a system undergoing a macroscopic radial pulsation or expansion.
Finally in section 4 we discuss the history of the discovery of tidal streams and give the methods recently invented for finding new ones. Once some globular cluster extensions or proper motions are reliably determined these methods will become far more powerful.
## 2. Core Collapse, Mass-Segregation, Rotation and the Dynamical Main Sequence
Much of globular cluster evolution follows from careful dimensional analysis, see Inagaki & Lynden-Bell (1983). Returning to our Question 1 we look for a self-similar solution with the density in the core and inner halo of the form
$$\rho (r,t)=\rho _0(t)\rho _{}\left(r_{}\right)\mathrm{where}r_{}=r/r_c(t)$$
Here $`r_c(t)`$ is the core’s radius.
Since the relaxation is far more rapid in the core than in the outer parts we may set $`/t0`$ for large $`r`$.
$$0=\left(\mathrm{}n\rho /t\right)_r=d\mathrm{}n\rho _0/dt\left(d\mathrm{}n\rho _{}/dr_{}\right)d\mathrm{}nr_c/dt$$
Therefore for large $`r`$
$$d\mathrm{}n\rho _0/d\mathrm{}nr_c=d\mathrm{}n\rho _{}/d\mathrm{}nr_{}=\alpha $$
Since the expression on the left is a function of $`t`$ alone, while that on the right is a function of $`r_{}`$ alone, it follows that both are constant (at a value we have called $`\alpha `$). Thus for all time
$$\rho _0=Ar_c^\alpha $$
(1)
and at large $`r_{}`$, $`\rho _{}=r_{}^\alpha `$.
The evolution of the core is due to the relaxation so
$$d\mathrm{}n\rho _0/dt=K_1/T_c$$
(2)
where $`K_1`$ is a dimensionless constant and $`T_c`$ is the relaxation time in the core. A standard evaluation of $`T_c`$ gives for a velocity dispersion $`\sigma `$, Binney & Tremaine (1987)
$$T_c^1=3G^2m\rho \mathrm{}n/\sigma ^3$$
(3)
and we would like to evaluate this for the core when there is a distribution of stellar masses. It would seem natural to take the distribution function to be an equilibrium Maxwell-Boltzmann there proportional to $`\mathrm{exp}\left[\beta m\left(\frac{1}{2}v^2\psi _0\right)\right]`$, but this cannot be true at energies close to the escape energy. The closeness of the escape energy is well emphasised by the following argument from the Virial Theorem.
$$2T=\underset{I}{}m_Iv_I^2=V=\frac{1}{2}\underset{I}{}m_I\psi _I$$
(4)
Now $`\psi _I=\frac{1}{2}v_e^2`$ where $`v_e`$ is the velocity of escape to infinity, hence for any cluster obeying the virial theorem
$$mv^2=\frac{1}{4}mv_e^2.$$
(5)
So that the mass weighted $`rms`$ velocity is HALF the mass weighted $`rms`$ velocity of escape. At a thermodynamic equilibrium we have equipartition of $`\frac{1}{2}mv^2`$ so the above expression implies that stars with less than a quarter of the mean mass have an rms velocity equal to the escape velocity. This cannot occur in practice. Thus the equipartition of the Maxwell-Boltzmann distribution must be modified and as Michie (1963) first showed the lowered Maxwellian of the Michie-King models \[evaluated in full by King (1966)\] does this approximately viz
$$f=B(m)\left\{\mathrm{exp}\left[\beta m\left(\frac{1}{2}v^2\psi \right)\right]\mathrm{exp}\left(\beta m\psi _e\right)\right\}$$
(6)
Although in well developed core collapse such models have insufficient temperature contrast between the centre and the halo, Lynden-Bell & Wood (1968), Lynden-Bell & Eggleton (1980), nevertheless we shall adopt such a form for the low energy stars in the core (for which relaxation is the most rapid). If $`F(m)dm`$ gives the number density of stars at mass $`m`$ in the core, so that
$$\rho _0=mF(m)𝑑m$$
(7)
then we may re-express the relaxation rate in terms of the central ‘temperature’ $`\beta ^1=\sigma ^2/m`$ and obtain from (3)
$$T_c^1=3G^2\beta ^{3/2}\rho _0\mathrm{}nm^{7/2}/m$$
(8)
where
$$m^{7/2}=m^{7/2}F(m)𝑑m/F(m)𝑑m$$
(9)
Evidently the relaxation rate, $`T_c^1`$, which determines the overall evolution of the cluster, depends on the mean seven halves power of the mass, evaluated over the core of the cluster. As evolution proceeds, the lighter stars in the core are preferentially expelled from it but simultaneously the heavy stars are gradually eliminated via stellar evolution. Before using (8) in (2) we need to re-express $`\beta `$ in terms of $`\rho _0,r_c`$ and $`m`$. If $`M_c=\frac{4}{3}\pi \rho _0r_c^3`$ is the core mass then
$$3\sigma ^2=3(\beta m)^1=GM_c/r_c=\frac{4\pi }{3}G\rho _0r_c^2$$
(10)
where we have chosen to define $`r_c`$ so that the constant of proportion is one.
$$T_c^1=K_2G^{1/2}\rho _0^{1/2}r_c^3m_c=K_2G^{1/2}A^{1/2}m_c(\rho _0/A)^{(6\alpha )/(2\alpha )}$$
(11)
where we have used (1) to express $`r_c`$ in terms of $`\rho _0`$, $`K_2`$ is the dimensionless constant $`\left(\frac{81}{8\pi ^{3/2}}\mathrm{}n\right)`$ and $`m_c(t)=m^{7/2}/m^{5/2}`$. To use expression (11) to solve (2) we must first evaluate the time dependence of $`m_c`$. This is caused by two separate effects, the progressive expulsion of lower mass stars from the core increases $`m_c`$ but it is mainly dependent on the high masses and they diminish steadily as stellar evolution and stellar death take their toll. With such a high power of the mass involved it is likely that stellar evolution plays the dominant role with $`m_c`$ behaving similarly to the age cut off. This suggests the approximate form
$$m_c=m_{}\left(t/t_{}\right)^{\delta 1}\mathrm{where}\delta 2/3,t_{}10^{10}\mathrm{yr},t>10^6\mathrm{yr}$$
and $`t`$ is time measured from the time the cluster was created. Inserting this into (11) and (11) into (2) the resulting equation for $`\rho _0`$ is readily solved to give
$$\rho _0=C|t_0^\delta t^\delta |^{2\alpha /(6\alpha )}$$
(12)
and therefore
$$r_c=(A/C)^{1/\alpha }|t_c^\delta t^\delta |^{2/(6\alpha )}$$
(13)
where $`t_c`$ is a constant of integration which gives the time of core collapse and the constant $`C`$ is
$$C=A\left[\frac{6\alpha }{2\alpha \delta }K_2G^{1/2}A^{1/2}m_{}t_{}^{1\delta }\right]^{2\alpha /(6\alpha )}.$$
As in Inagaki & Lynden-Bell (1983), we expect (12) and (13) to hold also after core collapse but then $`\rho _0`$ becomes a characteristic density rather than the central one which remains infinite if we ignore binaries and the giant gravothermal oscillations.
Of course in the absence of stellar evolution $`m_c`$ will increase and the parameterisation
$$m_c=m_{}(t/T_c)^\delta $$
might then be more appropriate but though such a model is soluble I shall not detail it here.
Expressions (12) and (13) can only be considered the solution to our problem once $`\alpha `$ is known. As Lynden-Bell and Eggleton showed $`\alpha `$ emerges as an eigen value in the full theory and for the equal mass case $`\alpha =2.22`$. It is readily seen that $`\alpha `$ must lie between the equal mass isothermal sphere with $`\alpha =2`$ and the limiting case with infinite core binding energy $`\alpha =5/2`$. With the heavy masses more concentrated to the core one may expect $`2.22<\alpha <2.5`$ and perhaps toward the upper end of that range, on the other hand once most of the lighter stars have been ejected from the core one expects the evolution to revert toward the single mass case. Our best guess is therefore
$$\alpha =2.3\pm 0.08.$$
(14)
In discussing Question 2 we note that the flattening of a cluster is of second order in $`\mathrm{\Omega }`$ the angular velocity and may be neglected to first order. In any potential of the form
$$\psi =\psi _0(r)+r^2\psi _2(\theta )$$
(15)
the expression $`I=\frac{1}{2}J^2\psi _2`$ is an exact integral of the motion where $`𝐉=(𝐫\times 𝐯)`$. In practice any non-spherical part of $`\psi `$ is second order in $`\mathrm{\Omega }`$ and the motion of a star is well approximated as lying in a precessing plane. The rate of precession can be worked out from the couple that the non-spherical potential exerts on the unperturbed rosette orbit that the star would have in the absence of asphericity. Thus in practice $`|𝐉|`$ is almost an integral of the motion while $`J_z`$ and the energy, $`ϵ`$, are exact ones in the absence of evolution. In practice it is convenient to use the radial action $`J_r=\frac{1}{2\pi }\sqrt{2[ϵ+\psi _0(r)]J^2/r^2}𝑑r`$ in place of $`ϵ`$. In potentials of the form (15) it is an exact integral but for more general forms of $`\psi `$ it is only approximate. A particular advantage of using $`|𝐉|,J_r,J_z`$ is that they are adiabatically invariant for slow changes in the potential $`\psi `$. Thus if we express the distribution function of a rotating cluster in the form $`f=f(|𝐉|,J_r,J_z,t)`$ where $`f`$ evolves little in one orbital time, then $`f/t`$ is $`(f/t)_{\mathrm{encounters}}`$ and there is no supplementary term due to the change of global potential induced by these encounters since for such changes the $`J`$ (although not $`ϵ`$) are adiabatically invariant. The $`J`$ have dimensions of $`\sigma r_c`$ and $`f`$ has the dimensions of $`\rho _0\sigma ^3`$, so in self-similar evolution the evolving $`f`$ will have the form
$$f=\rho _0\sigma ^3F_{}(|𝐉^{}|,J_r^{},J_z^{})$$
(16)
where $`F_{}`$ is a dimensionless function of its dimensionless arguments $`|𝐉^{}|=|𝐉|/(\sigma r_c),J_r^{}=J_r/(\sigma r_c),J_z^{}=J_z/(\sigma r_c).`$ As before $`\rho _0=Ar_c^\alpha `$ so $`\sigma r_cr_c^{2\alpha /2}`$. These same principles may be applied to discuss the way cluster core rotation evolves as the core contracts. Near the centre, relaxation is quite fast so we expect a rotating Maxwellian at energies little affected by escape
$$f=B\mathrm{exp}\sigma ^2(ϵ\mathrm{\Omega }_0J_z).$$
(17)
Our question is how does $`\mathrm{\Omega }_0`$ evolve during core collapse. In self-similar evolution the quantities that are dimensionless do not evolve, so $`\mathrm{\Omega }_0r_c/\sigma =\mathrm{const}`$. Hence as the core shrinks we have our prime result
$$\mathrm{\Omega }_0r_c^{\alpha /2}\alpha 2.22.$$
(18)
So as $`r_c`$ shrinks, the core should rotate faster and faster. This result is very sensibly between the $`\mathrm{\Omega }_0r_c^2=\mathrm{const}`$ that would follow if the heat flowed out of the core much more readily than angular momentum, and the $`\mathrm{\Omega }=\mathrm{const}`$ that would follow if the opposite were true. Away from the centre it is plausible that the system remembers the rotation it had when that part of the inner halo was part of the core; in reality it gains a little more angular momentum as it leaves the remaining core, because the remaining core loses it. This suggests the behaviour $`\mathrm{\Omega }r^{\alpha /2}`$ but since $`\mathrm{\Omega }^2`$ and $`G\rho `$ have the same dimensions a better prediction is that $`\mathrm{\Omega }`$ behaves approximately as $`\rho ^{1/2}`$; hence we conclude that cluster cores and inner haloes rotate as
$$\mathrm{\Omega }=\mathrm{\Omega }_0(\rho /\rho _0)^{1/2}=\mathrm{\Omega }_0\rho _{}^{1/2}$$
(19)
with $`\mathrm{\Omega }_0`$ evolving according to (18). For $`Q_1+Q_2`$ the same results will hold with $`\alpha `$ given by (14). Such predictions should be compared with the numerical work of Spurzem (2000) and the rapidly growing data from observations. A further consequence of (19) is that $`\mathrm{\Omega }^2/\pi G\rho `$ which determines the flattening will be constant in the core and inner halo thus the isophotes there should have $`b/a\mathrm{constant}`$ whereas if the core evolved with $`\mathrm{\Omega }_0r_c^2`$ the central isophotes would be more flattened. By contrast a uniformly rotating cluster has the greatest flattening in the outer isophotes. Turning now to Question 3 the theory of the Dynamical Main Sequence of Globular Clusters as yet awaits someone brave enough to make strong hypotheses such as ‘The only binaries that matter are in the core and the energy flux through the inner halo is constant and drives either the escape from the cluster or the expansion of its halo.’ Of course if core collapse ceases due to binaries then the steady core will relax to uniform rotation which will gradually extend from the core outwards.
## 3. Newton’s N-body Problem and its Generalisations
Let the force on body $`I`$ due to body $`J`$ be of the form
$$𝐅_{IJ}=km_Im_J(𝐫_J𝐫_I)$$
(20)
We sum over all $`J`$ to get the force on $`I`$ since the $`I=J`$ term is zero.
$$𝐅_I=km_I(m_Jr_JMr_I)=km_IM(\overline{𝐫}𝐫_I)$$
where $`M`$ is the total mass. Thus the total force is directed to the centre of mass $`\overline{𝐫}`$ and is proportional to the distance from it. Newton’s (1687) general solution is that each particle moves in a central ellipse about the centre of mass which itself moves uniformly in a straight line. Newton’s system has a total potential energy
$$V+\frac{1}{2}KM^2r^2$$
(21)
where
$$r^2=M^1m_I(𝐫_J\overline{𝐫})^2$$
(22)
The Virial theorem reads
$$\frac{1}{2}\ddot{I}=2T+nV$$
(23)
where $`Vr^n`$ so $`n=2`$ for the above system. To generalise. Newton’s result consider systems with total potential energies of the form
$$V=V(r)$$
with $`r`$ given by (22). This problem has a singular beauty if we consider the 3N dimensional space in which the $`𝐗_I=\sqrt{\frac{m_I}{M}}(𝐫_I\overline{𝐫})`$ are the coordinates. Letting the 3N vector $`𝐫=(𝐱_1,𝐱_2\mathrm{},𝐱_N)`$ we note that $`r^2`$ is given by (22). The initial conditions in this space are the initial values of $`𝐫`$ and $`\dot{𝐫}`$. The accelerations in this space are central since $`V=V(r)`$ so the acceleration also lies in the plane defined by the initial $`𝐫`$ and $`\dot{𝐫}`$. Thus the motion continues in that plane. In fact the whole problem now reduces to the planar orbit problem under the action of the central potential $`V(r)`$. Back in 3 space each particle feels the central force $`V^{}(r)m_IM^1(𝐫_I\overline{𝐫})/r`$ so each particle orbits in a plane. Its motion may be obtained by projecting the motion of the representative point in 3N down into the plane of the motion of particle I. As $`r`$ changes periodically around the planar orbit in 3N space it follows that $`Mr^2`$ vibrates forever so there is no violent relaxation to a Virial equilibrium. One may generalise this result still further by taking $`V=V_0(r)+r^2V_2(𝐫/r)`$ then the Virial theorem reads
$$\frac{1}{2}Md^2r^2/dt^2=2TrV_0^{}(r)+2V_2=2E2V_0rV_0^{}$$
since this last expression does not involve $`V_2`$ it follows that $`r`$ vibrates in the same way as it did when $`V_2`$ was 0. Hence once again there is no violent relaxation. Pretty N-body problems of this type are given by the forces
$$F_{IJ}=Gm_Im_J(𝐫_J𝐫_I)\left[\frac{1}{r^3}\frac{k}{(𝐫_J𝐫_I)^4}\right]$$
where $`r`$ is given by (22) or alternatively
$$F_{IJ}=G^{}m_Im_J(𝐫_J𝐫_I)\left[1k(𝐫_J𝐫_I)^4\right]$$
both of these force laws have long range attractions and short range repulsions so they may produce liquid-like and solid-like phases . In each case the vibration in $`r`$ separates from the other motions so if lattices are possible, lattices with a breathing pulsation will be too. It is possible to do a statistical mechanics of the ‘angular’ motions only while $`r`$ remains breathing but for details the reader should refer to Lynden-Bell & Lynden-Bell (1999a). The 1999b paper solves Schrodinger’s equation for both spin 0 Bosons and spin $`\frac{1}{2}`$ Fermions giving the only known non-trivial three dimensional N-body solutions for interacting Bose or Fermi gases. The degenerate Fermi case gives a white-dwarf-like solution.
## 4. Streams 1950-2000
Streams about the Galaxy have much in common with the meteor streams left behind by comets in the solar system which were discovered earlier. However the galactic orbits are rosettes rather than ellipses and their planes precess about the galactic pole due to the aspherical potential. Bertil Lindblad (1958)(1961) was one of the first to consider the theory of streams spreading from a common origin but already Eggen (1959) (1989) was on the march finding moving groups in the velocity space among both disc and halo stars. As data were refined Eggen changed his mind about the central velocities of some groups and this led to consequent changes in their membership which only increased scepticism among his critics. However Eggen was convinced he was onto something and claimed that it was much easier for others to criticise the inclusion of a star when new data showed its membership to be unlikely, than to discover the group to start with.
The subject of streams really come to life with the discovery of the Magellanic Stream stretching more than 120 around the galaxy. But this discovery was the culmination of a sequence. In 1965 Nan Dieter discovered a high velocity cloud at the South Galactic Pole. Soon afterward Hulsbosch & Raimond (1966) published their survey of such objects and showed that cloud to be strongly elongated. It was the technological innovation of Wrixon’s low noise diodes that allowed Wannier & Wrixon (1972) to see fainter 21 cm emission. Their great arc of 21 cm emission through Dieter’s cloud showed a systematic variation of radial velocity by several hundred km/sec across the sky. At Herstmonceux Kalnajs (1972) excitedly drew my attention to their paper claiming that it must be of Galactic scale. A day or two later he discovered that both the line of the arc and the radial velocity extrapolated to the Magellanic Clouds. Mathewson too was very excited by this discovery and no sooner had he got back to Australia than he, Cleary & Murray (1974) started delineating how the stream joined the Magellanic Clouds. A few years later during a non-photometric night at Sutherland I was imagining how the Magellanic Stream lay across the sky when it struck me that other satellites of the Milky Way might be members of streams and might likewise have high velocity clouds associated with them. I rapidly found that Draco & Ursa Minor are almost antipodal to the LMC and the SMC in the galactocentric sky. This meant that they too might be long lost members of the Magellanic Stream, Lynden-Bell (1976). Kunkel & Demers (1977) noticed the dwarf spheroidals were distributed non-randomly over the sky but their compromise plane was some 30 from the Magellanic Stream. In their attempt to explain the antipody of Draco & Ursa Minor, Hunter & Tremaine (1977) noticed that their elongations lay along the stream while I searched for further streams using elongations as a supplementary guide, Lynden-Bell (1982); however it was not until 1995 that this developed into a systematic method. Since tidal tearing occurs essentially in the plane of orbit, one needs to find objects scattered along a great circle in the galactocentric sky. Such great circles are hard to see in any of the projections of the sky so we need a good method of finding them. An object may be associated with the set of all great circles that pass through it. The poles of these great circles lie on another great circle which we call a polar path. (The pole of that great circle lies of course at the object we started with). The great circle through any two objects has its pole at the intersection of their polar paths. If 3 objects lie on a great circle all 3 polar paths will go through the same point. Thus we may find objects with common great circles by plotting their polar paths (in any equal area projection) and looking for multiple intersections. Figure 1a which plots 45 polar paths for dwarf spheroidals and outer globular clusters shows a result that more closely resembles a ball of wool than a useful scientific plot. More information is needed to get any definitive groupings. A known proper motion would reduce the polar path of an object to a very short segment determined by the directional error of the proper motion. If we are prepared to assume that the elongation of an object is along its orbit, as appears to be true for Draco & Ursa Minor and might be true for the outer extensions recently found around globulars (e.g., Grillmair et. al., Meylan this volume) then likewise the pole must lie in only a small sector of the polar path (see Figure 1b). But great circles are not enough. Radial velocities provide a discriminant. If all the objects in a stream follow the orbit of the progenitor then they will share the same specific energy $`E`$ and specific angular momentum $`h`$.
Taking sets of objects that share a common polar intersection and therefore lie on a galactocentric great circle one may plot $`E_r=E\frac{1}{2}h^2r^2=\frac{1}{2}v_r^2\psi (r)`$ against $`r^2`$ the inverse square of their galactocentric distances. This should be a straight line so $`E`$ and $`h^2`$ can be deduced as the intercept and gradient of such a line. In practice the observed radial velocity is not the galactocentric one so $`v_r`$ can be deduced by estimating an initial $`𝐡`$ and then iterating to find a consistent gradient of $`\frac{1}{2}h^2`$. In practice this method converges rapidly but needless to say some of the objects do not lie on straight lines and are interlopers rather than stream members. Since the method gives $`h^2`$ from the observed radial velocities, and since the poles of the great circle involved are known, we have used these methods not only to produce possible streams among the dwarf spheroidals and outer globulars but also to predict their proper motions. The method works well for the Magellanic Stream and for the Sagittarius Stream which was discovered by Ibata, Gilmore & Irwin (1995). There is now great hope that we can identify orbital streams associated with the later merging events emphasised by Searle and Zinn. Further work may take us back to the orbits of objects that took part in the initial formation itself. The ELS (1962) picture had the bulk of the metal poor globulars formed in the fragmentation of the initial collapse, but following Ambartsumian we know that most of the stars would have been formed in looser aggregates and stellar associations which subsequently broke up and spread around the Galaxy. The idea that ELS pictured a perfectly smooth collapse is not in their paper but has been added by others who have not read it. Eggen (1959) had already worked on streams in the halo, the fragmentation of collapsing systems was treated by Hunter (1962), and I was working on the Large Scale Instabilities of shape that would make the collapse non-axis symmetric, Lynden-Bell (1964). The non-linear behaviour of these instabilities was later treated in the simple non-rotating case by Lin, Mestel & Shu (1965). Those who want to read for themselves what was thought at that time would do well to read my paper to IAU Symposium No. 31 on the Formation of the Galaxy, Lynden-Bell (1967). Streams from the initial collapse are less likely to be preserved than the younger streams from late additions, but finding correlated stellar motions from the initial creation of the Galaxy was what ELS was all about and their method of using abundances and adiabatic invariants as fossilised evidence on how the Galaxy formed is probably its most lasting legacy to astronomy. The merging of galaxies was suggested as important in the very percipient paper on bridges and streams by the Toomre’s (1972) and White’s (1976) demonstration of fragmentation of a uniform overdensity showed that even within the free fall time of the whole much fragmentation and subsequent merging of those fragments occurs.
Returning to streams there is too much modern work to detail all of it, but work by Johnston (1998) holds out the prospect of more accurate determination of the Galaxy’s potential using streams of stars just escaping from a satellite galaxy. By contrast, Helmi & White (1999) take such stars to be emitted with a not inconsiderable velocity dispersion so such streams would be harder to detect and would define the potential less accurately. It is not yet clear to me which view is more realistic. Lastly, Putman et. al. have at last detected a forward stream from the Magellanic Clouds, albeit somewhat skew of the main great circle. Currently the most plausible explanation is that it was gas around the SMC that was torn but much of the forward stream ran into the LMC and the rest was diverted by it, however other interpretations are possible.
### Acknowledgments.
The theory of rotation was developed after the conference while writing up this contribution under the hospitality of IUCAA, Pune, India.
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# Exact Solution of Asymmetric Diffusion With Second-Class Particles of Arbitrary Size
## I Introduction
The similarity between the master equation describing time fluctuations in nonequilibrium problems and the Schrödinger equation describing the quantum fluctuations of quantum spin chains turns out to be fruitful for both areas of research \- . Since many quantum chains are known to be exactly integrable through the Bethe ansatz, this provides exact information on the related stochastic model. At the same time the classical physical intuition and probabilistic methods successfully applied to nonequilibrium systems give new insights into the understanding of the physical and algebraic properties of quantum chains.
An example of this fruitful interchange is the problem of asymmetric diffusion of hard-core particles on the one dimensional lattice. This model is related to the exactly integrable anisotropic Heisenberg chain (XXZ model). However if we demand this quantum chain to be invariant under a quantum group symmetry $`U_q(SU(2))`$, we have to introduce, for the equilibrium statistical system, unphysical surface terms, which on the other hand have a nice and simple interpretation for the related nonequilibrium stochastic system .
In the area of exactly integrable models it is well known that one of the possible extensions of the spin-$`\frac{1}{2}`$ XXZ chain to higher spins, is the anisotropic spin-1 Sutherland model (grading $`ϵ_1=ϵ_2=ϵ_3=1`$) . On the other hand in the area of diffusion limited reactions a simple extension of the asymmetric diffusion problem is the problem of diffusion with first and second class of particles \- . In this problem a mixture of two classes of hard-core particles diffuses on the lattice. Particles belonging to the first class ignore the presence of those in the second class, i. e., they see them in the same way as they see the holes (empty sites). In it was shown that for open boundary conditions the anisotropic spin-1 Sutherland model and this last stochastic model are exactly related, the Hamiltonian governing the quantum or time fluctuations of both models being given in terms of generators of a Hecke algebra, invariant under the quantum group $`U_qSU(3)`$. In this paper we are going to derive the exact solution of the associated quantum chain, in a closed lattice. Recently (see also ) we have shown that without losing its exact integrability, we can consider the problem of asymmetric diffusion with an arbitrary mixture of molecules with different sizes (even zero), as long they do not interchange positions, that is, there is no reactions. Motivated by these results we are going to extend the asymmetric diffusion problem with a second class of particles, to the case where the particles in each class have an arbitrary size, in units of the lattice spacing. Unlike the case of asymmetric diffusion problem, we have in this case a nested Bethe ansatz .
The paper is organized as follows. In the next section we introduce the generalized asymmetric model with second-class particles and derive the associated quantum chain. In section 3 the Bethe ansatz solution of the model is presented in a pedagogical and self contained way, which should be simple to follow by an audience not specialized in exactly integrable models. Finally in section 4 we present our conclusions, with some possible generalizations of the stochastic problem considered in this paper, and some perspectives on future work.
## II The generalized asymmetric diffusion model with first- and second-class of particles
A simple extension of the asymmetric exclusion model, where hard-core particles diffuse on the lattice, is the problem where a mixture of particles belonging to different classes (first and second class) diffuses on the lattice. This problem was used to describe shocks \- in nonequilibrium and also has a stationary probability distribution that can be expressed via the matrix-product ansatz . In this model we have $`n_1`$ and $`n_2`$ molecules belonging to the first and second class, respectively. Both classes of molecules diffuse asymmetrically, but with the same asymmetrical rates, whenever they encounter empty sites (holes) at nearest-neighbor sites. However, when molecules of different classes are at their minimum separation, the molecules of the first class exchange position with the same rate as they diffuse, and consequently the first-class molecules see no difference between molecules belonging to the second class and holes.
We now introduce a generalization of the above model, where instead of having unit size, the molecules in the first and second class have in general distinct sizes $`s_1`$ and $`s_2`$ ($`s_1,s_2=1,2,\mathrm{}`$), respectively, in units of lattice spacing. In Fig. 1 we show some examples of molecules of different sizes. We may think of a molecule of size $`s`$ as formed by $`s`$ monomers (size 1), and for simplicity, we define the position of the molecule as the center of its leftmost monomer. The molecules have a hard-core repulsion: the minimum distance $`d_{\alpha \beta }`$, in units of the lattice spacing, between molecules $`\alpha `$ and $`\beta `$, with $`\alpha `$ in the left, is given by $`d_{\alpha \beta }=s_\alpha `$. In order to describe the occupancy of a given configuration of molecules we define, at each site $`i`$ of a lattice with $`N`$ sites, a variable $`\beta _i`$ ($`i=1,2,\mathrm{},N`$), taking the values $`\beta _i=0,1`$ and 2, representing site $`i`$ empty (size $`s_0=1`$), occupied by a molecule of class 1 (size $`s_1`$) or a molecule of class 2 (size $`s_2`$), respectively. Then the allowed configurations are given by the set $`\{\beta _i\}`$ ($`i=1,\mathrm{},N`$), where for each pair $`(\beta _i,\beta _j)0`$ with $`j>i`$ we should have $`jis_{\beta _i}`$.
The time evolution of the probability distribution $`P(\{\beta \},t)`$, of a given configuration $`\{\beta \}`$ is given by the master equation
$$\frac{P(\{\beta \},t)}{t}=\mathrm{\Gamma }(\{\beta \}\{\beta ^{}\})P(\{\beta \},t)+\mathrm{\Gamma }(\{\beta ^{}\}\{\beta \})P(\{\beta ^{}\},t),$$
(1)
where $`\mathrm{\Gamma }(\{\beta \}\{\beta ^{}\})`$ is the transition rate for configuration $`\{\beta \}`$ to change to $`\{\beta ^{}\}`$. In the present model we only allow, whenever it is possible, the particles to diffuse to nearest-neighbor sites, or to exchange positions. The possible motions are diffusion to the right
$`1_i\mathrm{}_{i+s_1}`$ $``$ $`\mathrm{}_i\mathrm{\hspace{0.33em}\hspace{0.33em}1}_{i+1},(\text{rate}\mathrm{\Gamma }_R)`$ (2)
$`2_i\mathrm{}_{i+s_2}`$ $``$ $`\mathrm{}_i\mathrm{\hspace{0.33em}\hspace{0.33em}2}_{i+1},(\text{rate}\mathrm{\Gamma }_R),`$ (3)
diffusion to the left
$`\mathrm{}_i\mathrm{\hspace{0.33em}\hspace{0.33em}1}_{i+1}`$ $``$ $`1_i\mathrm{}_{i+s_1},(\text{rate}\mathrm{\Gamma }_L)`$ (4)
$`\mathrm{}_i\mathrm{\hspace{0.33em}\hspace{0.33em}2}_{i+1}`$ $``$ $`2_i\mathrm{}_{i+s_2},(\text{rate}\mathrm{\Gamma }_L),`$ (5)
and interchange of particles
$`1_i\mathrm{\hspace{0.33em}\hspace{0.33em}2}_{i+s_1}`$ $``$ $`2_i\mathrm{\hspace{0.33em}\hspace{0.33em}1}_{i+s_2},(\text{rate}\mathrm{\Gamma }_R)`$ (6)
$`2_i\mathrm{\hspace{0.33em}\hspace{0.33em}1}_{i+s_2}`$ $``$ $`1_i\mathrm{\hspace{0.33em}\hspace{0.33em}2}_{i+s_1},(\text{rate}\mathrm{\Gamma }_L).`$ (7)
As we see from (6), particles in the first class interchange positions with those of second class with the same rate as they interchange positions with the empty sites (diffusion). We should remark however that unless the second class particles have unit size ($`s_2=1`$), the net effect of the second class particles in those of the first class is distinct from the effect produced by the holes, since as the result of the exchange the first class of particles will move by $`s_2`$ lattice size units, accelerating its diffusion.
The master equation (1) can be written as a Schrödinger equation in Euclidean time (see Ref. for general application for two body processes)
$$\frac{|P>}{t}=H|P>,$$
(8)
if we interpret $`|P>P(\{\beta \},t)`$ as the associated wave function. If we represent $`\beta _i`$ as $`|\beta >_i`$ the vector $`|\beta >_1|\beta >_2\mathrm{}|\beta >_N`$ will give us the associated Hilbert space. The process (2)-(6) gives us the Hamiltonian (see Ref. for general applications)
$`H`$ $`=`$ $`D{\displaystyle \underset{j}{}}H_j`$ (9)
$`H_j`$ $`=`$ $`𝒫\{{\displaystyle \underset{\alpha =1}{\overset{2}{}}}[ϵ_+(E_j^{0\alpha }E_{j+1}^{\alpha 0}E_j^{\alpha \alpha }E_{j+1}^{00})+ϵ_{}(E_j^{\alpha 0}E_{j+1}^{0\alpha }E_j^{00}E_{j+1}^{\alpha \alpha })]`$ (11)
$`+ϵ_+(E_j^{21}E_{j+s_2}^{10}E_{j+s_1}^{02}E_j^{11}E_{j+s_2}^{00}E_{j+s_1}^{22})+ϵ_{}(E_j^{12}E_{j+s_1}^{20}E_{j+s_2}^{01}E_j^{22}E_{j+s_1}^{00}E_{j+s_2}^{11})\}𝒫`$
with
$$D=\mathrm{\Gamma }_R+\mathrm{\Gamma }_L,ϵ_+=\frac{\mathrm{\Gamma }_R}{\mathrm{\Gamma }_R+\mathrm{\Gamma }_L},ϵ_{}=\frac{\mathrm{\Gamma }_L}{\mathrm{\Gamma }_R+\mathrm{\Gamma }_L}(ϵ_++ϵ_{}=1),$$
(12)
and periodic boundary conditions. The matrices $`E^{\alpha ,\beta }`$ are $`3\times 3`$ matrices with a single nonzero element $`(E^{\alpha ,\beta })_{i,j}=\delta _{\alpha ,i}\delta _{\beta ,j}`$ ($`\alpha ,\beta ,i,j\mathrm{Z})`$. The projector $`𝒫`$ in (9), projects out from the associated Hilbert space the vectors $`|\{\beta \}>`$ which represent forbidden positions of the molecules due to their finite size, which mathematically means that for all $`i,j`$ with $`\beta _i,\beta _j0,|ij|s_{\beta _i}(j>i)`$. The constant D in (9) fixes the time scale and for simplicity we chose $`D=1`$. A particular simplification of (9) occurs when molecules of the first and second class have the same size $`s_1=s_2=s`$. In this case the Hamiltonian can be expressed as an anisotropic nearest-neighbor interaction spin-1 $`SU(3)`$ chain. Moreover in the case where their sizes are unity ($`s=1`$) the model may be related to the $`SU(3)`$ Sutherland model with twisted boundary conditions .
## III The Bethe ansatz equations
We present in this section the exact solution of the general quantum chain (9). Since the present paper is going to appear in a special issue where the main subject is nonequilibrium physics, we are going to present a pedagogical and self-contained derivation of the exact solution, that can clearly be followed by a nonexpert in the arena of exactly solvable models in statistical mechanics.
Due to the conservation of particles in the diffusion and interchange processes the total number of particles $`n_1`$ and $`n_2`$ of particles of class 1 and 2 are good quantum numbers and consequently we can split the associated Hilbert space into block disjoint sectors labeled by the numbers $`n_1`$ and $`n_2`$ ($`n_1=0,1,\mathrm{},n;n_2=nn_1;n=n_1+n_2`$). We therefore consider the eigenvalue equation
$$H|n_1,n_2>=E|n_1,n_2>,$$
(13)
where
$$|n_1,n_2>=\underset{\{Q\}}{}\underset{\{x\}}{}f(x_1,Q_1;\mathrm{};x_n,Q_n)|x_1,Q_1;\mathrm{};x_n,Q_n>.$$
(14)
Here $`|x_1,Q_1;\mathrm{};x_n,Q_n>`$ means the configuration where a particle of class $`Q_i(Q_i=1,2)`$ is at position $`x_i`$ ($`x_i=1,\mathrm{},N`$). The summation $`\{Q\}=\{Q_1,\mathrm{},Q_n\}`$ extends over all permutations of $`n`$ numbers in which $`n_1`$ terms are 1 and $`n_2`$ terms are 2, while the summation $`\{x\}=\{x_1,\mathrm{},x_n\}`$ runs, for each permutation $`\{Q\}`$, in the set of the $`n`$ nondecreasing integers satisfying
$$x_{i+1}x_i+s_{Q_i},i=1,\mathrm{},n1,s_{Q_1}x_nx_1Ns_{Q_n}.$$
(15)
Before getting the results for general values of $`n`$ let us consider initially the cases where we have 1 or 2 particles.
n = 1. For one particle on the chain (class 1 or 2), as a consequence of the translational invariance of (9) it is simple to verify directly that the eigenfunctions are the momentum-$`k`$ eigenfunctions
$$|1,0>=\underset{x=1}{\overset{N}{}}f(x,1)|x,1>,\text{or}|0,1>=\underset{x=1}{\overset{N}{}}f(x,2)|x,2>,$$
(16)
with
$$f(x,1)=f(x,2)=e^{ikx},k=\frac{2\pi l}{N},l=0,1,\mathrm{},N1,$$
(17)
and energy given by
$$E=e(k)(ϵ_{}e^{ik}+ϵ_+e^{ik}1).$$
(18)
n =2. For two particles of classes $`Q_1`$ and $`Q_2`$ ($`Q_1,Q_2=1,2`$) on the lattice, the eigenvalue equation (13) gives us two distinct relations depending on the relative location of the particles. The first relation applies to the case in which a particle of class $`Q_1`$ (size $`s_{Q_1}`$) is at position $`x_1`$ and a particle $`Q_2`$ (size $`s_{Q_2}`$) is at position $`x_2`$, where $`x_2>x_1+s_{Q_1}`$. We obtain in this case the relation
$`Ef(x_1,Q_1;x_2,Q_2)=ϵ_+f(x_11,Q_1;x_2,Q_2)ϵ_{}f(x_1,Q_1;x_2+1,Q_2)`$ (19)
$``$ $`ϵ_{}f(x_1+1,Q_1;x_2,Q_2)ϵ_+f(x_1,Q_1;x_21,Q_2)+2f(x_1,Q_1;x_2,Q_2),`$ (20)
where we have used the relation $`ϵ_++ϵ_{}=1`$. This last equation can be solved promptly by the ansatz
$$f(x_1,Q_1;x_2,Q_2)=e^{ik_1x_1}e^{ik_2x_2},$$
(21)
with energy
$$E=e(k_1)+e(k_2).$$
(22)
Since this relation is symmetric under the interchange of $`k_1`$ and $`k_2`$, we can write a more general solution of (19) as
$`f(x_1,Q_1;x_2,Q_2)`$ $`=`$ $`{\displaystyle \underset{P}{}}A_{P_1,P_2}^{Q_1,Q_2}e^{i(k_{P_1}x_1+k_{P_2}x_2)}`$ (23)
$`=`$ $`A_{1,2}^{Q_1,Q_2}e^{i(k_1x_1+k_2x_2)}+A_{2,1}^{Q_1,Q_2}e^{i(k_2x_1+k_1x_2)}`$ (24)
with the same energy as in (22). In (23) the summation is over the permutations $`P=P_1,P_2`$ of(1,2). The second relation applies when $`x_2=x_1+s_{Q_1}`$. In this case instead of (19) we have
$`Ef(x_1,Q_1`$ ; $`x_1+s_{Q_1},Q_2)=ϵ_+f(x_11,Q_1;x_1+s_{Q_2},Q_2)ϵ_{}f(x_1,Q_1;x_1+s_{Q_1}+1,Q_2)`$ (25)
$``$ $`\stackrel{~}{ϵ}_{Q_1,Q_2}f(x_1,Q_2;x_1+s_{Q_2},Q_1)+(1+\stackrel{~}{ϵ}_{Q_1,Q_2})f(x_1,Q_1;x_1+s_{Q_1},Q_2),`$ (26)
where
$$\stackrel{~}{ϵ}_{1,1}=\stackrel{~}{ϵ}_{2,2}=0,\stackrel{~}{ϵ}_{1,2}=ϵ_+\text{and}\stackrel{~}{ϵ}_{2,1}=ϵ_{}.$$
(27)
If we now substitute the ansatz (23) with the energy (22), the constants $`A_{12}^{Q_1,Q_2}`$ and $`A_{21}^{Q_1,Q_2}`$, initially arbitrary, should now satisfy
$$\underset{P}{}\{\left[𝒟_{P_1,P_2}+e^{ik_{P_2}}(1\stackrel{~}{ϵ}_{Q_1,Q_2})\right]e^{ik_{P_2}(s_{Q_1}1)}A_{P_1,P_2}^{Q_1,Q_2}+\stackrel{~}{ϵ}_{Q_1,Q_2}e^{ik_{P_2}s_{Q_2}}A_{P_1,P_2}^{Q_2,Q_1}\}=0$$
(28)
where
$$𝒟_{i,j}=(ϵ_++ϵ_{}e^{i(k_i+k_j)}).$$
(29)
At this point it is convenient to consider separately the case where $`Q_1=Q_2`$ from those where $`Q_1Q_2`$. If $`Q_1=Q_2=Q`$ ($`Q=1,2`$) eq. (28) gives
$$\underset{P}{}\left(𝒟_{P_1,P_2}+e^{ik_{P_2}}\right)e^{ik_{P_2}(s_Q1)}A_{P_1,P_2}^{Q_1,Q_2}=0$$
(30)
and the cases $`Q_1Q_2`$ give us the equations
$`{\displaystyle \underset{P}{}}\left[\begin{array}{cc}𝒟_{P_1,P_2}+e^{ik_{P_2}}ϵ_{}& ϵ_+e^{ik_{P_2}}\\ ϵ_{}e^{ik_{P_2}}& 𝒟_{P_1,P_2}+ϵ_+e^{ik_{P_2}}\end{array}\right]\left[\begin{array}{c}e^{ik_{P_2}(s_11)}A_{P_1,P_2}^{1,2}\\ e^{ik_{P_2}(s_21)}A_{P_1,P_2}^{2,1}\end{array}\right]=0.`$ (35)
Performing the above summation we obtain, after lengthy but straightforward algebra, the following relation among the amplitudes
$`\left[\begin{array}{c}A_{1,2}^{1,2}e^{ik_2(s_11)}\\ A_{1,2}^{2,1}e^{ik_2(s_21)}\end{array}\right]={\displaystyle \frac{𝒟_{1,2}+e^{ik_1}}{𝒟_{1,2}+e^{ik_2}}}\left\{1\mathrm{\Phi }(k_1,k_2)\left[\begin{array}{cc}ϵ_+& ϵ_+\\ ϵ_{}& ϵ_{}\end{array}\right]\right\}\left[\begin{array}{c}A_{2,1}^{12}e^{ik_1(s_11)}\\ A_{2,1}^{21}e^{ik_1(s_21)}\end{array}\right],`$ (42)
where
$$\mathrm{\Phi }(k_1,k_2)=\frac{e^{ik_1}e^{ik_2}}{𝒟_{1,2}+e^{ik_1}}.$$
(43)
Equations (30) and (42) can be written in a compact form
$$A_{P_1,P_2}^{Q1,Q2}=\mathrm{\Xi }_{P_1,P_2}\underset{Q_1^{},Q_2^{}=1}{\overset{2}{}}S_{Q_1^{},Q_2^{}}^{Q_1,Q_2}(k_{P_1},k_{P_2})A_{P_2,P_1}^{Q_2^{},Q_1^{}},(Q_1,Q_2=1,2)$$
(44)
with
$$\mathrm{\Xi }_{l,j}=\frac{𝒟_{l,j}+e^{ik_l}}{𝒟_{l,j}+e^{ik_j}}=\frac{ϵ_++ϵ_{}e^{i(k_l+k_j)}e^{ik_l}}{ϵ_++ϵ_{}e^{i(k_l+k_j)}e^{ik_j}},$$
(45)
where we have introduced the $`S`$ matrix. From 30 and (42), the only non zero elements of $`S_{Q_1^{},Q_2^{}}^{Q_1,Q_2}`$ are:
$`S_{1,1}^{1,1}(k_1,k_2)`$ $`=`$ $`e^{i(k_1k_2)(s_11)},s_{2,2}^{2,2}(k_1,k_2)=e^{i(k_1k_2)(s_21)},`$ (46)
$`S_{2,1}^{1,2}(k_1,k_2)`$ $`=`$ $`\left[1ϵ_+\mathrm{\Phi }(k_1,k_2)\right]e^{i(k_1k_2)(s_11)},`$ (47)
$`S_{1,2}^{2,1}(k_1,k_2)`$ $`=`$ $`\left[1ϵ_{}\mathrm{\Phi }(k_1,k_2)\right]e^{i(k_1k_2)(s_21)},`$ (48)
$`S_{1,2}^{1,2}(k_1,k_2)`$ $`=`$ $`ϵ_+\mathrm{\Phi }(k_1,k_2)e^{ik_1(s_21)}e^{ik_2(s_11)},`$ (49)
$`S_{2,1}^{2,1}(k_1,k_2)`$ $`=`$ $`ϵ_{}\mathrm{\Phi }(k_1,k_2)e^{ik_1(s_11)}e^{ik_2(s_21)}.`$ (50)
A graphical representation of the S matrix is shown in Fig. 2a. Equations (44) do not fix the “wave numbers” $`k_1`$ and $`k_2`$. In general, these numbers are complex, and are fixed due to the cyclic boundary condition
$$f(x_1,Q_1;x_2,Q_2)=f(x_2,Q_2;x_1+N,Q_1),$$
(51)
which from (23) gives the relation
$$A_{1,2}^{Q_1Q_2}=e^{ik_1N}A_{2,1}^{Q_2,Q_1},A_{2,1}^{Q_1,Q_2}=e^{ik_2N}A_{2,1}^{Q_2,Q_1}.$$
(52)
This last equation, when solved by exploiting (44)-(46), gives us the possible values of $`k_1`$ and $`k_2`$, and from (22) the eigenenergies in the sector with 2 particles. Instead of solving these equations for the particular case $`n=2`$ let us now consider the case of general $`n`$.
General n. The above calculation can be generalized for arbitrary values of $`n_1`$ and $`n_2`$ of particles of classes 1 and 2, respectively ($`n_1+n_2=n`$). The ansatz for the wave function (14) becomes
$$f(x_1,Q_1;\mathrm{};x_n,Q_n)=\underset{P}{}A_{P_1,\mathrm{},P_n}^{Q_1,\mathrm{},Q_n}e^{i(k_{P_1}x_1+\mathrm{}+k_{P_n}x_n)},$$
(53)
where the sum extends over all permutations $`P`$ of the integers $`1,2,\mathrm{},n`$. For the components $`|x_1,Q_1;\mathrm{};x_n,Q_n>`$ where $`x_{i+1}x_i>s_{Q_i}`$ for $`i=1,2,\mathrm{},n`$, it is simple to see that the eigenvalue equation (13) is satisfied by the ansatz (53) with energy
$$E=\underset{j=1}{\overset{n}{}}e(k_j).$$
(54)
On the other hand if a pair of particles of class $`Q_i,Q_{i+1}`$ is at positions $`x_i,x_{i+1}`$, where $`x_{i+1}=x_i+s_{Q_i}`$, equation (13) with the ansatz (53) and the relation (54) give us the generalization of relation (44), namely
$$A_{\mathrm{},P_i,P_{i+1},\mathrm{}}^{\mathrm{},Q_i,Q_{i+1},\mathrm{}}=\mathrm{\Xi }_{P_i,P_{i+1}}\underset{Q_1^{},Q_2^{}}{\overset{2}{}}S_{Q_1^{},Q_2^{}}^{Q_i,Q_{i+1}}(k_{P_i},k_{P_{i+1}})A_{\mathrm{},P_{i+1},P_i,\mathrm{}}^{\mathrm{},Q_2^{},Q_1^{},\mathrm{}}Q_i,Q_{i+1}=1,2,$$
(55)
with $`S`$ given by eq. (46). Inserting the ansatz (53) in the boundary condition
$$f(x_1,Q_1;\mathrm{};x_n,Q_n)=f(x_2,Q_2;\mathrm{};x_n,Q_n;x_1+N,Q_1)$$
(56)
we obtain the additional relation
$$A_{P_1,\mathrm{},P_n}^{Q_1,\mathrm{},Q_n}=e^{ik_{P_1}N}A_{P_2,\mathrm{},P_n,P_1}^{Q_2,\mathrm{},Q_n,Q_1},$$
(57)
which together with (55) should give us the energies.
Successive applications of (55) give us in general distinct relations between the amplitudes. For example $`A_{\mathrm{},k_1,k_2,k_3,\mathrm{}}^{\mathrm{},\alpha ,\beta ,\gamma ,\mathrm{}}`$ relate to $`A_{\mathrm{},k_3,k_2,k_1,\mathrm{}}^{\mathrm{},\gamma ,\beta ,\alpha ,\mathrm{}}`$ by performing the permutations $`\alpha \beta \gamma \beta \alpha \gamma \beta \gamma \alpha \gamma \beta \alpha `$ or $`\alpha \beta \gamma \alpha \gamma \beta \gamma \alpha \beta \gamma \beta \alpha `$, and consequently the $`S`$-matrix should satisfy the Yang-Baxter equation
$`{\displaystyle \underset{\gamma ,\gamma ^{},\gamma ^{\prime \prime }=1}{\overset{2}{}}}S_{\gamma ,\gamma ^{}}^{\alpha ,\alpha ^{}}(k_1,k_2)S_{\beta ,\gamma ^{\prime \prime }}^{\gamma ,\alpha ^{\prime \prime }}(k_1,k_3)`$ $`S_{\beta ^{},\beta ^{\prime \prime }}^{\gamma ^{},\gamma ^{\prime \prime }}(k_2,k_3)=`$ (59)
$`{\displaystyle \underset{\gamma ,\gamma ^{},\gamma ^{\prime \prime }=1}{\overset{2}{}}}S_{\gamma ^{},\gamma ^{\prime \prime }}^{\alpha ^{},\alpha ^{\prime \prime }}(k_2,k_3)S_{\gamma ,\beta ^{\prime \prime }}^{\alpha ,\gamma ^{\prime \prime }}(k_1,k_3)S_{\beta ,\beta ^{}}^{\gamma ,\gamma ^{}}(k_1,k_2),`$
for $`\alpha ,\alpha ^{},\alpha ^{\prime \prime },\beta ,\beta ^{},\beta ^{\prime \prime }=1,2`$ and $`S`$ given by (46). In Fig. 2b we show graphically this equation. Actually the relation (59) is a necessary and sufficient condition to obtain a non-trivial solution for the amplitudes in Eq. (55).
We can verify by a long and straightforward calculation that for arbitrary values of $`s_1`$ and $`s_2`$, the $`S`$ matrix (46), satisfies the Yang-Baxter equation (59), and consequently we may use relations (55) and (57) to obtain the eigenenergies of the Hamiltonian (9). Applying relation (55) $`n`$ times on the right of equation (57) we obtain a relation between the amplitudes with the same momenta. Using the graphical representation in Fig. 3a for the $`S`$ matrix, we illustrate in Fig. 4 the result of such applications. We then obtain
$`A_{P_1,\mathrm{},P_n}^{Q_1,\mathrm{},Q_n}=e^{ik_{P_1}N}A_{P_2,\mathrm{},P_n,P_1}^{Q_2,\mathrm{},Q_n,Q_1}=\left({\displaystyle \underset{i=2}{\overset{n}{}}}\mathrm{\Xi }_{P_i,P_1}\right)e^{ik_{P_1}N}{\displaystyle \underset{Q_1^{},\mathrm{},Q_n^{}}{}}{\displaystyle \underset{Q_1^{\prime \prime },\mathrm{},Q_n^{\prime \prime }}{}}`$ (60)
$`S_{Q_1^{},Q_1^{\prime \prime }}^{Q_1,Q_2^{\prime \prime }}(k_{P_1},k_{P_1})S_{Q_2^{},Q_2^{\prime \prime }}^{Q_2,Q_3^{\prime \prime }}(k_{P_2},k_{P_1})\mathrm{}S_{Q_{n1}^{},Q_{n1}^{\prime \prime }}^{Q_{n1},Q_n^{\prime \prime }}(k_{P_{n1}},k_{P_1})S_{Q_n^{},Q_n^{\prime \prime }}^{Q_n,Q_1^{\prime \prime }}(k_{P_n},k_{P_1})A_{P_1,\mathrm{},P_n}^{Q_1^{},\mathrm{},Q_n^{}},`$ (61)
where we have introduced the harmless extra sum
$$1=\underset{Q_1^{\prime \prime },Q_2^{\prime \prime }=1}{\overset{2}{}}\delta _{Q_2^{\prime \prime },Q_1^{}}\delta _{Q_1^{\prime \prime },Q_1}=\underset{Q_1^{\prime \prime },Q_2^{\prime \prime }=1}{\overset{2}{}}S_{Q_1^{},Q_1^{\prime \prime }}^{Q_1,Q_2^{\prime \prime }}(k_{P_1},k_{P_1}).$$
(62)
In order to fix the values of $`\{k_j\}`$ we should solve (60), i.e., we should find the eigenvalues $`\mathrm{\Lambda }(k)`$ of the matrix
$$T(k)_{\{Q^{}\}}^{\{Q\}}=\underset{Q_1^{\prime \prime },\mathrm{},Q_n^{\prime \prime }=1}{\overset{2}{}}\left\{\left(\underset{l=1}{\overset{n1}{}}S_{Q_l^{},Q_l^{\prime \prime }}^{Q_l,Q_{l+1}^{\prime \prime }}(k_{P_l},k)\right)S_{Q_n^{},Q_n^{\prime \prime }}^{Q_n,Q_1^{\prime \prime }}(k_{P_n},k)\right\},$$
(63)
and the Bethe-ansatz equations which fix the set $`\{k_l\}`$ will be given from (60) by
$$e^{ik_jN}=(1)^{n1}\left(\underset{l=1}{\overset{n}{}}\mathrm{\Xi }_{l,j}\right)\mathrm{\Lambda }(k_j),j=1,\mathrm{},n.$$
(64)
Using the graphical representation for $`S`$ given in Fig. 3a the matrix $`T`$ in (63) can be represented graphically as in Fig. 5, and we identify $`T(k)`$ as the transfer matrix of an inhomogeneous 6-vertex model, on a periodic lattice, with Boltzmann weights $`S_{Q_1^{},Q_2^{}}^{Q_1,Q_2}(k_{P_l},k)`$ ($`l=1,\mathrm{},n`$). Consequently, in order to obtain the eigenenergies of the quantum chain (9) we should diagonalize the above transfer matrix $`T(k)`$.
Diagonalization of $`T(k)`$
The simplest way to diagonalize $`T`$ is through the introduction of the monodromy matrix $`(k)`$ , which is a transfer matrix of the inhomogeneous vertex model under consideration, where, instead of being periodic, the first and last link in the horizontal direction are fixed to the values $`\mu _1`$ and $`\mu _{n+1}`$ ($`\mu _1,\mu _{n+1}=1,2`$), that is
$$_{\{Q^{}\},\mu _1}^{\{Q\},\mu _{n+1}}(k)=\underset{\mu _2,\mathrm{},\mu _n}{}S_{Q_1^{},\mu _1}^{Q_1,\mu _2}(k_{P_1},k)S_{Q_2^{},\mu _2}^{Q_2,\mu _3}(k_{P_2},k)\mathrm{}S_{Q_{n1}^{},\mu _{n1}}^{Q_{n1},\mu _n}(k_{P_{n1}},k)S_{Q_n^{},\mu _n}^{Q_n,\mu _{n+1}}(k_{P_n},k).$$
(65)
The monodromy matrix $`_{\{Q^{}\},\mu _1}^{\{Q\},\mu _{n+1}}(k)`$ is represented in Fig. 6a and has coordinates $`\{Q\},\{Q^{}\}`$ in the vertical space ($`2^n`$ dimensions) and coordinates $`\mu _1,\mu _{n+1}`$ in the horizontal space (4 dimensions). This matrix satisfies the following important relations
$$S_{\nu _1,\mu _1}^{\nu _1^{},\mu _1^{}}(k^{},k)_{\{\alpha _l\},\mu _1^{}}^{\{\gamma _l\},\mu _{n+1}}(k)_{\{\gamma _l\},\nu _1^{}}^{\{\beta _l\},\nu _{n+1}}(k^{})=_{\{\alpha _l\},\nu _1}^{\{\gamma _l\},\nu _{n+1}^{}}(k^{})_{\{\gamma _l\},\mu _1}^{\{\beta _l\},\mu _{n+1}^{}}(k)S_{\nu _{n+1}^{},\mu _{n+1}^{}}^{\nu _{n+1},\mu _{n+1}}(k^{},k).$$
(66)
This relation is shown graphically in Fig. 6b; its validity follows directly from successive applications of the Yang-Baxter equations (59) (see Fig. 2b).
In order to exploit relation (66) let us denote the components of the monodromy matrix in the horizontal space by
$`A(k)`$ $`=`$ $`A(k)_{\{\alpha _l\}}^{\{\gamma _l\}}=_{\{\alpha _l\},1}^{\{\gamma _l\},1}(k),B(k)=B(k)_{\{\alpha _l\}}^{\{\gamma _l\}}=_{\{\alpha _l\},1}^{\{\gamma _l\},2}(k),`$ (67)
$`C(k)`$ $`=`$ $`C(k)_{\{\alpha _l\}}^{\{\gamma _l\}}=_{\{\alpha _l\},2}^{\{\gamma _l\},1}(k),D(k)=D(k)_{\{\alpha _l\}}^{\{\gamma _l\}}=_{\{\alpha _l\},2}^{\{\gamma _l\},2}(k),`$ (68)
and clearly the transfer matrix $`T(k)`$ of the periodic inhomogeneous lattice, we want to diagonalize, is given by
$$T(k)=A(k)+D(k).$$
(69)
As a consequence of (66) the matrices $`A`$, $`B`$, $`C`$ and $`D`$ in (67) obey some algebraic relations. The elements $`(\nu _1,\mu _1,\nu _{n+1},\mu _{n+1})=(1,1,2,2)`$, $`(1,1,1,2)`$ and $`(1,2,2,2)`$ give us the relations
$$[B(k),B(k^{})]=0,$$
(70)
$$A(k)B(k^{})=\frac{S_{1,1}^{1,1}(k,k^{})}{S_{1,2}^{1,2}(k,k^{})}B(k^{})A(k)\frac{S_{2,1}^{1,2}(k,k^{})}{S_{1,2}^{1,2}(k,k^{})}B(k)A(k^{}),$$
(71)
$$D(k)B(k^{})=\frac{S_{2,2}^{2,2}(k^{},k)}{S_{1,2}^{1,2}(k^{},k)}B(k^{})D(k)\frac{S_{1,2}^{2,1}(k^{},k)}{S_{1,2}^{1,2}(k^{},k)}B(k)D(k^{}),$$
(72)
respectively. The diagonalization of $`T(k)`$ in (69) will be done by exploiting the above relations. This procedure is known in the literature as the algebraic Bethe ansatz . The first step in this method follows from the identification of a reference state $`|\mathrm{\Omega }>`$, which should be an eigenstate of $`A(k)`$ and $`D(k)`$, and hence $`T(k)`$, but not of $`B(k)`$. In our problem it is simple to see that a possible reference state is $`|\mathrm{\Omega }>=|\{\alpha _l=1\}>_{l=1,\mathrm{},n}`$, which corresponds to a state with first-class particles only. It is simple to calculate
$`A(k)|\mathrm{\Omega }>`$ $`=`$ $`a(k)|\mathrm{\Omega }>,D(k)|\mathrm{\Omega }>=d(k)|\mathrm{\Omega }>,`$ (73)
$`C(k)|\mathrm{\Omega }>`$ $`=`$ $`0,B(k)|\mathrm{\Omega }>={\displaystyle \underset{i=1}{\overset{n}{}}}b_i(k)|\mathrm{\Omega }>,`$ (74)
where
$`a(k)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}S_{1,1}^{1,1}(k_{P_i},k),d(k)={\displaystyle \underset{i=1}{\overset{n}{}}}S_{2,1}^{1,2}(k_{P_i},k),`$ (75)
$`b_i(k)`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{i1}{}}}S_{1,1}^{1,1}(k_{P_l},k){\displaystyle \underset{l=j}{\overset{n}{}}}S_{1,2}^{1,2}(k_{P_l},k).`$ (76)
In Fig. 7 we show a graphical representation of the above relation. The idea in the algebraic Bethe ansatz is that $`B(k)`$ acts as a creation operator in the reference (“vacuum”) state, i.e., it creates particles of second class in a sea of particles of first class $`|\mathrm{\Omega }>`$, and it is expected that a general eigenvector of $`T(k)`$ in the sector with $`n_2`$ second-class particles will be given by the ansatz
$$\mathrm{\Psi }(k_1^{(1)},k_2^{(1)},\mathrm{},k_{n_2}^{(1)})=B(k_1^{(1)})B(k_2^{(1)})\mathrm{}B(k_{n_2}^{(1)})|\mathrm{\Omega }>.$$
(77)
The numbers $`\{k_i^{(1)}\},i=1,\mathrm{},n_2`$ here play the role of the “wave number” $`\{k_i\}`$ in the coordinate Bethe ansatz (53), and are going to be fixed by the eigenvalue equation
$$T(k)\mathrm{\Psi }(k_1^{(1)},\mathrm{},k_{n_2}^{(1)})=\left(A(k)+D(k)\right)\mathrm{\Psi }(k_1^{(1)},\mathrm{},k_{n_2}^{(1)})=\mathrm{\Lambda }(k)\mathrm{\Psi }(k_1^{(1)},\mathrm{},k_{n_2}^{(1)}).$$
(78)
Before deriving the relation for general values of $`n_2`$, let us consider initially the cases where $`n_2=1`$ and $`n_2=2`$.
$`𝐧_\mathrm{𝟐}=\mathrm{𝟏}`$. Using relation (71) and (72) in (78) we obtain
$`T(k)\mathrm{\Psi }(k_1^{(1)})`$ $`=`$ $`\left[{\displaystyle \frac{S_{1,1}^{1,1}(k,k_1^{(1)})}{S_{1,2}^{1,2}(k,k_1^{(1)})}}+d(k){\displaystyle \frac{S_{2,2}^{2,2}(k_1^{(1)},k)}{S_{1,2}^{1,2}(k_1^{(1)},k)}}\right]\mathrm{\Psi }(k_1^{(1)})`$ (79)
$``$ $`\left[{\displaystyle \frac{S_{2,1}^{1,2}(k,k_1^{(1)})}{S_{1,2}^{1,2}(k,k_1^{(1)})}}+d(k_1^{(1)}){\displaystyle \frac{S_{1,2}^{2,1}(k_1^{(1)},k)}{S_{1,2}^{1,2}(k_1^{(1)},k)}}\right]B(k)|\mathrm{\Omega }>=\mathrm{\Lambda }(k)\mathrm{\Psi }(k_1^{(1)}),`$ (80)
which is clearly satisfied if the coefficient of $`B(k)|\mathrm{\Omega }>`$ (unwanted term) vanishes. This condition, which fixes $`k_1^{(1)}`$, is given by
$$\underset{j=1}{\overset{n}{}}S_{1,2}^{1,2}(k_{P_j},k_1^{(1)})=1,$$
(81)
where we used the relation
$$\frac{S_{2,1}^{1,2}(k,k^{})}{S_{1,2}^{1,2}(k,k^{})}=\frac{S_{1,2}^{2,1}(k^{},k)}{S_{1,2}^{1,2}(k^{},k)},$$
(82)
valid for the $`S`$ matrix (46). The eigenvalue is given by
$$\mathrm{\Lambda }(k)=\frac{S_{1,1}^{1,1}(k,k_1^{(1)})}{S_{1,2}^{1,2}(k,k_1^{(1)})}+\frac{S_{2,2}^{2,2}(k_1^{(1)},k)}{S_{1,2}^{1,2}(k_1^{(1)},k)}\underset{j=1}{\overset{n}{}}S_{1,2}^{1,2}(k_{P_j},k),$$
(83)
provided (81)is satisfied.
$`𝐧_\mathrm{𝟐}=\mathrm{𝟐}`$. In this case the application of $`A(k)`$ and $`D(k)`$ in the ansatz (77) give us
$`A(k)\mathrm{\Psi }(k_1^{(1)},k_2^{(1)})`$ $`=`$ $`{\displaystyle \frac{S_{1,1}^{1,1}(k,k_1^{(1)})S_{1,1}^{1,1}(k,k_2^{(1)})}{S_{1,2}^{1,2}(k,k_1^{(1)})S_{1,2}^{1,2}(k,k_2^{(1)})}}a(k)\mathrm{\Psi }(k_1^{(1)},k_2^{(1)})`$ (84)
$``$ $`{\displaystyle \frac{S_{2,1}^{1,2}(k,k_1^{(1)})S_{1,1}^{1,1}(k_1^{(1)},k_2^{(1)})}{S_{1,2}^{1,2}(k,k_1^{(1)})S_{1,2}^{1,2}(k_1^{(1)},k_2^{(1)})}}a(k_1^{(1)})B(k)B(k_2^{(1)})|\mathrm{\Omega }>`$ (85)
$``$ $`{\displaystyle \frac{S_{2,1}^{1,2}(k,k_2^{(1)})S_{1,1}^{1,1}(k_2^{(1)},k_1^{(1)})}{S_{1,2}^{1,2}(k,k_2^{(1)})S_{1,2}^{1,2}(k_2^{(1)},k_1^{(1)})}}a(k_2^{(1)})B(k)B(k_1^{(1)})|\mathrm{\Omega }>,`$ (86)
and
$`D(k)\mathrm{\Psi }(k_1^{(1)},k_2^{(1)})`$ $`=`$ $`{\displaystyle \frac{S_{2,2}^{2,2}(k_1^{(1)},k)S_{2,2}^{2,2}(k_2^{(1)},k)}{S_{1,2}^{1,2}(k_1^{(1)},k)S_{1,2}^{1,2}(k_2^{(1)},k)}}d(k)\mathrm{\Psi }(k_1^{(1)},k_2^{(1)})`$ (87)
$``$ $`{\displaystyle \frac{S_{1,2}^{2,1}(k_1^{(1)},k)S_{2,2}^{2,2}(k_2^{(1)},k_1^{(1)})}{S_{1,2}^{1,2}(k_1^{(1)},k)S_{1,2}^{1,2}(k_2^{(1)},k_1^{(1)})}}d(k_1^{(1)})B(k)B(k_2^{(1)})|\mathrm{\Omega }>`$ (88)
$``$ $`{\displaystyle \frac{S_{1,2}^{2,1}(k_2^{(1)},k)S_{2,2}^{2,2}(k_1^{(1)},k_2^{(1)})}{S_{1,2}^{1,2}(k_2^{(1)},k)S_{1,2}^{1,2}(k_1^{(1)},k_2^{(1)})}}d(k_2^{(1)})B(k)B(k_1^{(1)})|\mathrm{\Omega }>,`$ (89)
where besides the relation (71) and (72) we have used (70), which imply $`\mathrm{\Psi }(k_1^{(1)},k_2^{(1)})=\mathrm{\Psi }(k_2^{(1)},k_1^{(1)})`$. From (84) and (87) the condition (78) gives us
$$\mathrm{\Lambda }(k)=\underset{i=1}{\overset{n}{}}S_{1,1}^{1,1}(k_{P_i},k)\underset{j=1}{\overset{2}{}}\frac{S_{1,1}^{1,1}(k,k_j^{(1)})}{S_{1,2}^{1,2}(k,k_j^{(1)})}+\underset{i=1}{\overset{n}{}}S_{1,2}^{1,2}(k_{P_i},k)\underset{j=1}{\overset{2}{}}\frac{S_{2,2}^{2,2}(k_j^{(1)},k)}{S_{1,2}^{1,2}(k_j^{(1)},k)},$$
(90)
under the condition, which fixes $`k_1^{(1)}`$ and $`k_2^{(1)}`$,
$$\underset{i=1}{\overset{n}{}}S_{1,2}^{1,2}(k_{P_i},k_\alpha ^{(1)})=\underset{\beta =1(\alpha )}{\overset{2}{}}\frac{S_{1,1}^{1,1}(k_\alpha ^{(1)},k_\beta ^{(1)})S_{1,2}^{1,2}(k_\beta ^{(1)},k_\alpha ^{(1)})}{S_{1,2}^{1,2}(k_\alpha ^{(1)},k_\beta ^{(1)})S_{2,2}^{2,2}(k_\beta ^{(1)},k_\alpha ^{(1)})}\underset{j=1}{\overset{n}{}}S_{1,1}^{1,1}(k_{P_j},k_\alpha ^{(1)})\alpha =1,2.$$
(91)
In deriving this last expression the relation (82) was also used.
General $`𝐧_\mathrm{𝟐}`$. The previous procedure can be iterated straightforwadly for arbitrary numbers $`n_2`$ of second-class particles, which gives
$`A(k)\mathrm{\Psi }(k_1^{(1)},\mathrm{},k_{n_2}^{(1)})`$ $`=`$ $`\mathrm{\Lambda }^{(A)}(k)\mathrm{\Psi }(k_1^{(1)},\mathrm{},k_{n_2}^{(1)})`$ (92)
$``$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\mathrm{\Lambda }_i^{(A)}(k)\mathrm{\Psi }(k,k_1^{(1)},\mathrm{},k_{i1}^{(1)},k_{i+1}^{(1)},\mathrm{},k_{n_2}^{(1)}),`$ (93)
$`D(k)\mathrm{\Psi }(k_1^{(1)},\mathrm{},k_{n_2}^{(1)})`$ $`=`$ $`\mathrm{\Lambda }^{(D)}(k)\mathrm{\Psi }(k_1^{(1)},\mathrm{},k_{n_2}^{(1)})`$ (94)
$``$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\mathrm{\Lambda }_i^{(D)}(k)\mathrm{\Psi }(k,k_1^{(1)},\mathrm{},k_{i1}^{(1)},k_{i+1}^{(1)},\mathrm{},k_{n_2}^{(1)}),`$ (95)
where
$`\mathrm{\Lambda }^{(A)}(k)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}S_{1,1}^{1,1}(k_{P_j},k){\displaystyle \underset{i=1}{\overset{n_2}{}}}{\displaystyle \frac{S_{1,1}^{1,1}(k,k_i^{(1)})}{S_{1,2}^{1,2}(k,k_i^{(1)})}},`$ (96)
$`\mathrm{\Lambda }^{(D)}(k)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}S_{1,2}^{1,2}(k_{P_j},k){\displaystyle \underset{i=1}{\overset{n_2}{}}}{\displaystyle \frac{S_{2,2}^{2,2}(k_i^{(1)},k)}{S_{1,2}^{1,2}(k_i^{(1)},k)}},`$ (97)
$`\mathrm{\Lambda }_j^{(A)}(k)`$ $`=`$ $`{\displaystyle \frac{S_{2,1}^{1,2}(k,k_j^{(1)})}{S_{1,2}^{1,2}(k,k_j^{(1)})}}a(k_j^{(1)}){\displaystyle \underset{l=1(j)}{\overset{n_2}{}}}{\displaystyle \frac{S_{1,1}^{1,1}(k_j^{(1)},k_l^{(1)})}{S_{1,2}^{1,2}(k_j^{(1)},k_l^{(1)})}},`$ (98)
$`\mathrm{\Lambda }_j^{(D)}(k)`$ $`=`$ $`{\displaystyle \frac{S_{1,2}^{2,1}(k_j^{(1)},k)}{S_{1,2}^{1,2}(k_j^{(1)},k)}}d(k_j^{(1)}){\displaystyle \underset{l=1(j)}{\overset{n_2}{}}}{\displaystyle \frac{S_{2,2}^{2,2}(k_l^{(1)},k_j^{(1)})}{S_{1,2}^{1,2}(k_l^{(1)},k_j^{(1)})}}.`$ (99)
Then $`\mathrm{\Psi }(k_1^{(1)},\mathrm{},k_{n_2}^{(1)})`$ is an eigenvector of $`T`$ with eigenvalue $`\mathrm{\Lambda }(k)=\mathrm{\Lambda }^{(A)}(k)+\mathrm{\Lambda }^{(D)}(k)`$, i. e.,
$$\mathrm{\Lambda }(k)=\underset{i=1}{\overset{n}{}}S_{1,1}^{1,1}(k_{P_i},k)\underset{\alpha =1}{\overset{n_2}{}}\frac{S_{1,1}^{1,1}(k,k_\alpha ^{(1)})}{S_{1,2}^{1,2}(k,k_\alpha ^{(1)})}+\underset{i=1}{\overset{n}{}}S_{1,2}^{1,2}(k_{P_i},k]\underset{\alpha =1}{\overset{n_2}{}}\frac{S_{2,2}^{2,2}(k_\alpha ^{(1)},k)}{S_{1,2}^{1,2}(k_\alpha ^{(1)},k)},$$
(101)
if the following conditions, which fix $`\{k_1^{(1)},\mathrm{},k_{n_2}\}`$ are satisfied:
$$\mathrm{\Lambda }_i^{(A)}(k)+\mathrm{\Lambda }_i^{(D)}(k)=0,i=1,\mathrm{},n_2.$$
(102)
Using the relations (82), this last condition can be written as
$$\underset{j=1}{\overset{n}{}}\frac{S_{1,2}^{1,2}(k_{P_j},k_\alpha ^{(1)})}{S_{1,1}^{1,1}(k_{P_j},k_\alpha ^{(1)})}=\underset{\beta =1(\alpha )}{\overset{n_2}{}}\frac{S_{1,1}^{1,1}(k_\alpha ^{(1)},k_\beta ^{(1)})S_{1,2}^{1,2}(k_\beta ^{(1)},k_\alpha ^{(1)})}{S_{1,2}^{1,2}(k_\alpha ^{(1)},k_\beta ^{(1)})S_{2,2}^{2,2}(k_\beta ^{(1)},k_\alpha ^{(1)})},\alpha =1,\mathrm{},n_2,$$
(103)
which concludes the diagonalization of $`T(k)`$.
Now let us return to our original problem of finding the eigenvalues of the Hamiltonian (9). The Bethe-ansatz equations will be obtained by inserting in (64) the eigenvalues evaluated at $`k_j`$, i. e., $`\mathrm{\Lambda }(k_j)`$, given in (101), with the condition (103). Taking into account that $`S_{1,2}^{1,2}(k,k)=0`$, we obtain
$$e^{ik_jN}=()^{n1}\underset{l=1}{\overset{n}{}}\left(\mathrm{\Xi }_{l,j}S_{1,1}^{1,1}(k_l,k_j)\right)\underset{\alpha =1}{\overset{n_2}{}}\frac{S_{1,1}^{1,1}(k_j,k_\alpha ^{(1)})}{S_{1,2}^{1,2}(k_j,k_\alpha ^{(1)})},j=1,\mathrm{},n,$$
(104)
with the condition
$$\underset{j=1}{\overset{n}{}}\frac{S_{1,2}^{1,2}(k_j,k_\alpha ^{(1)})}{S_{1,1}^{1,1}(k_j,k_\alpha ^{(1)})}=\underset{\beta =1(\alpha )}{\overset{n_2}{}}\frac{S_{1,1}^{1,1}(k_\alpha ^{(1)},k_\beta ^{(1)})S_{1,2}^{1,2}(k_\beta ^{(1)},k_\alpha ^{(1)})}{S_{1,2}^{1,2}(k_\alpha ^{(1)},k_\beta ^{(1)})S_{2,2}^{2,2}(k_\beta ^{(1)},k_\alpha ^{(1)})},\alpha =1,\mathrm{},n_2.$$
(105)
Writing explicitly the $`S`$-matrix elements (46) in (104) and (105) we can state that the energies, in the sector with $`n_1`$ first-class particles and $`n_2`$ second-class particles are given by
$$E=\underset{j=1}{\overset{n}{}}(ϵ_{}e^{ik_j}+ϵ_+e^{ik_j}1),$$
(106)
where $`\{k_j,j=1,\mathrm{},n\}`$ are given by the solutions of
$`e^{ik_j(Nn_2(s_2s_1))}`$ $`=`$ $`()^{n_11}{\displaystyle \underset{l=1}{\overset{n}{}}}\left[\left({\displaystyle \frac{e^{ik_j}}{e^{ik_l}}}\right)^{s_11}{\displaystyle \frac{ϵ_++ϵ_{}e^{i(k_l+k_j)}e^{ik_j}}{ϵ_++ϵ_{}e^{i(k_l+k_j)}e^{ik_l}}}\right]`$ (107)
$`\times `$ $`{\displaystyle \underset{\alpha =1}{\overset{n_2}{}}}{\displaystyle \frac{ϵ_+(e^{ik_j}e^{ik_\alpha ^{(1)}})}{ϵ_++ϵ_{}e^{i(k_j+k_\alpha ^{(1)})}e^{ik_\alpha ^{(1)}}}}(j=1,\mathrm{},n),`$ (108)
$`e^{ik_j(s_2s_1)n}`$ $`(ϵ_+)^n`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{e^{ik_j}e^{ik_\alpha ^{(1)}}}{ϵ_++ϵ_{}e^{i(k_j+k_\alpha ^{(1)})}e^{ik_j}}}=`$ (109)
$`()^{n_11}`$ $`{\displaystyle \underset{\beta =1}{\overset{n_2}{}}}{\displaystyle \frac{ϵ_++ϵ_{}e^{i(k_\alpha ^{(1)}+k_\beta ^{(1)})}e^{ik_\alpha ^{(1)}}}{ϵ_++ϵ_{}e^{i(k_\alpha ^{(1)}+k_\beta ^{(1)})}e^{ik_\beta ^{(1)}}}}(\alpha =1,\mathrm{},n_2).`$ (110)
It is interesting to observe that in the particular case where $`n_2=0`$ we obtain the Bethe-ansatz equations, recently derived , for the asymmetric diffusion problem with particles of size $`s_1`$. Also the case $`s_1=s_2=1`$ give us the corresponding Bethe ansatz equations for the standard problem of second class particles. The Bethe-ansatz equations for the fully asymmetric problem is obtained by setting $`ϵ_+=1`$ and $`ϵ_{}=0`$.
## IV Conclusions and generalizations
We obtained through the Bethe ansatz the exact solution of the problem of first- and second-class particles diffusing and interchanging positions on the lattice. We show that the solution can be derived in the general case where the particles have arbitrary sizes.
Some extensions of our results can be made. A first generalization is the problem where $`M>2`$ hierarchical ordered classes of particles (1st, 2nd, …,Mth classes), with sizes $`s_i`$ ($`i=1,\mathrm{},M`$), besides diffusing on the lattice, interchange positions in such way that particles of higher hierarchical classes see the lower ones in the same way as they see the holes (0 class). The allowed processes and the Hamiltonians are simple generalizations of relations (2)-(6) and (9), respectively. The Hamiltonian obtained in a open lattice will be $`U_q(SU(M))`$ symmetric and in the case $`s_1=s_2=\mathrm{}=s_M=1`$ it is related with the spin $`S=\frac{M}{2}`$ Sutherland model . The model is certainly integrable and its eigenspectra will be given in terms of $`M`$-nested Bethe ansatz equations which are generalizations of (107)-(109).
A further quite interesting generalization of our model happens when we consider molecules in one or both classes with size $`s=0`$. Molecules of size zero do not occupy space on the lattice, having no hard-core exclusion effect. Consequently we may have, at a given lattice point, an arbitrary number of them. The Bethe-ansatz solution presented in the previous section is extended directly in this case (the equations are the same) and the eigenenergies are given by fixing in (106)-(109) the appropriate sizes of the molecules. It is interesting to remark that particles of second class with size $`s_2=0`$, contrary to the case $`s_2>1`$, where they “accelerate” the diffusion of the first class particles, now they “retard” the diffusive motion of these particles. The quantum Hamiltonian in the cases where the particles have size zero are written in terms of spin $`s=\mathrm{}`$ quantum chains.
The Bethe-ansatz equations in the case of the asymmetric diffusion, with particles of unit size , or with arbitrary size , was used to obtain the finite-size corrections of the mass gap $`G_N`$ of the associated quantum chain. The real part of these finite-size corrections are governed by the dynamical critical exponent $`z`$, i. e.,
$$\text{Re}(G_N)N^z.$$
(111)
The calculation of the exponent $`z`$ for the model presented in this paper, with particles of arbitrary sizes, is presently in progress .
###### Acknowledgements.
This work was supported in part by Conselho Nacional de Desenvolvimento Científico e Tecnológico - CNPq - Brazil and by the Russian Foundation of Fundamental Investigation ( Grant 99-02-17646).
Figure Captions
Figure 1 - Example of configurations of molecules with distinct sizes $`s`$ in a lattice of size $`N=6`$. The coordinates of the molecules are denoted by the black squares.
Figure 2 - Graphical representations of: a) the $`S`$ matrix (46) and b) the Yang-Baxter equations (59).
Figure 3 - Graphical representation of the relation (55). The circles represents the particle’s indices $`\{Q\}`$ and $`\{Q^{}\}`$, and the lines the “wave numbers” $`\{k_{P_i}\}`$. The variables represented in the full circles are summed.
Figure 4 - a) Graphical representation of the result of $`n`$ iterations of relation (55) on the right of equation (57). The diagramatic representation of Fig. 3 was used. b) A graphical representation of equation (62).
Figure 5 - Graphical representation of the transfer matrix $`T(k)`$, of the inhomogeneous 6-vertex model, in a periodic lattice of size $`n`$.
Figure 6 - Graphical representation of the relation (66) between the $`S`$ matrix and the monodromy matrix $``$.
Figure 7 - Graphical representation of the relations (73).
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# Conductivity sum rule, implication for in-plane dynamics and 𝑐-axis response
## Abstract
Recently observed $`c`$-axis optical sum rule violations indicate non-Fermi liquid in-plane behavior. For coherent $`c`$-axis coupling, the observed flat, nearly frequency independent $`c`$-axis conductivity $`\sigma _1(\omega )`$ implies a large in-plane scattering rate $`\mathrm{\Gamma }`$ around $`(0,\pi )`$ and therefore any pseudogap that might form at low frequency in the normal state will be smeared. On the other hand incoherent $`c`$-axis coupling places no restriction on the value of $`\mathrm{\Gamma }`$ and gives a more consistent picture of the observed sum rule violation which, we find in some cases, can be less than half.
$`C`$-axis electrodynamics is important in distinguishing high-$`T_c`$ cuprates from conventional superconductors and for understanding mechanism in the cuprates. Since the cuprates have layered structures, interlayer coupling between two adjacent CuO<sub>2</sub> planes plays an essential role, and its exact nature impacts on the recent experimental observation of a significant violation of the conventional optical sum rule of Ferrel, Glove and Tinkham. The normalized missing spectral weight (NMSW) is different from one and of order of a half in several high-$`T_c`$ cuprates. NMSW is the difference between the area under the real part of the conductivity $`\sigma ^{n(s)}(\omega )\left[=\sigma _1^{n(s)}(\omega )+i\sigma _2^{n(s)}(\omega )\right]`$ in the normal minus superconducting state, devided by the superfluid density which is obtained from $`\sigma _2^s(\omega )`$. We have pointed out in Ref. that in order to explain NMSW of less than one, it is necessary to consider non-Fermi liquid models for the in-plane dynamics, regardless of the nature of the interlayer coupling. A possible such model introduces a pseudogap in the normal state above $`T_c`$ as is observed. Another is the ”mode” coupling model determined from consideration of ARPES data by Norman et al. This model has been used to describe kinetic, as opposed to potential energy driven superconductivity.
In this paper, we show that coherent $`c`$-axis coupling, even with a pseudogap, cannot easily explain recent experimental findings, but that incoherent coupling describes them including the observed low frequency behavior of the effective NMSW.
The interlayer coupling is represented by a Hamiltonian $`H_c`$:
$$H_c=\underset{ij,\sigma }{}[t_{ij}c_{i1\sigma }^+c_{i2\sigma }+H.c.],$$
(1)
where the hopping matrix $`t_{ij}`$ describes weak interlayer tunneling and $`c_{i1}^+`$ creates an electron with spin $``$ at the site $`i`$ in the plane $`1`$. We classify interlayer couplings in two classes because the results are remarkably different depending on their nature: i) One is coherent coupling $`(t_{ij}=t_{})`$, which originates from an overlap of electronic wave functions between the two planes. For an in-plane Fermi liquid, it was shown in Ref. that the superfluid density $`(\rho _s)`$ is equal to the missing spectral weight $`(N_n(\omega _c)N_s(\omega _c)=8_{0^+}^{\omega _c}d\omega ^{}[\sigma _{1c}^n(\omega ^{})\sigma _{1c}^s(\omega ^{})]`$, where $`\omega _c`$ is a cutoff frequency of the order of a bandwidth). Thus coherent coupling (unless the density of states has a strong variation with energy) can explain the experimental results on optimally doped YBCO and over-doped Tl$`2201`$. ii) The other is incoherent coupling, for which $`t_{ij}=V_i\delta _{ij}`$, where $`V_i`$ is an impurity potential. In this case, we showed in Ref. that NMSW $`1.58`$. The characteristic difference between coherent and incoherent coupling is whether or not electron momentum is conserved in the interlayer transfer.
In the presence of an external vector potential $`A_z`$, $`H_c`$ is modified to $`H_c(A_z)`$ by the phase factor $`\mathrm{exp}(\pm ieA_z)`$. It is sufficient to expand $`H_c(A_z)`$ to second order in $`A_z`$. The current $`j_c=\delta H_c(A_z)/\delta A_z=j_p+j_d`$, where $`j_p=ied_{i\sigma }t_{}\left[c_{i1\sigma }^+c_{i2\sigma }c_{i2\sigma }^+c_{i1\sigma }\right]`$ and $`j_d=e^2d^2H_cA_z`$ with $`d`$ the interlayer spacing. In linear response theory, $`j_c=\left[\mathrm{\Pi }+e^2d^2H_c\right]A_z`$, where $`\mathrm{\Pi }`$ is the current-current correlation function associated with $`j_p`$ and $`H_c`$ is the perturbation of $`j_d`$ due to $`H_c`$.
From the $`c`$-axis conductivity sum rule the superfluid density $`\rho _s`$ can be written as
$$\rho _s=8_{0^+}^{\omega _c}𝑑\omega \left[\sigma _{1c}^n(\omega )\sigma _{1c}^s(\omega )\right]4\pi e^2d^2\left[H_c^sH_c^n\right],$$
(2)
where $`\omega _c`$ is the cutoff frequency for the interband transitions that $`H_c`$ does not describe, and we use units such that $`\mathrm{}=c=k_B=1`$ and set the volume of the system to be unity.
The penetration depth $`\lambda _c`$ can be calculated in two ways. Based on the Kramers-Kronig relation for the conductivity, we obtain $`\lambda _c`$ as $`1/4\pi \lambda _c^2=lim_{\omega 0}[\omega \text{I}\text{m}\sigma _c(0,\omega )]`$. Alternatively, using Eq.(2) we can also calculate $`\lambda _c(=1/\sqrt{\rho _s})`$. Combining these two equations for $`\lambda _c`$, we obtain the formula:
$$\frac{(N_nN_s)}{\rho _s}=\frac{1}{2}+\frac{1}{2}\frac{\underset{\omega }{}\underset{𝐤,𝐩}{}|t_{𝐤𝐩}|^2[G(𝐤,\omega )G(𝐩,\omega )G_0(𝐤,\omega )G_0(𝐩,\omega )]}{_\omega _{𝐤,𝐩}|t_{𝐤𝐩}|^2F(𝐤,\omega )F^+(𝐩,\omega )},$$
(3)
where $`G(𝐤,\omega )`$ and $`F(𝐤,\omega )`$ are superconducting Green functions and $`G_0(𝐤,\omega )`$ is in the normal state. For coherent $`c`$-axis coupling $`(|t_{𝐤𝐩}|^2=t_{}^2\delta _{𝐤𝐩})`$ electron momentum parallel to the plane is conserved and $`t_{}^2`$ may depend on the in-plane momentum. For incoherent coupling, an impurity configuration average is implied over a potential $`V_i`$ $`(|t_{𝐤𝐩}|^2=|V_{𝐤𝐩}|^2)`$, and electron momentum is not conserved, and $`𝐤`$ and $`𝐩`$ remain unconstraint.
Coherent coupling. As the simplest case, one can consider a normal state spectral function $`A_0(𝐤,\omega )=\mathrm{\Gamma }/\left[(\omega \xi _𝐤)^2+\mathrm{\Gamma }^2\right]`$ based on Fermi-liquid theory. The in-plane scattering rate $`\mathrm{\Gamma }`$ is expected to depend strongly on direction of the momentum $`𝐤`$ in the two dimensional Brillouin zone. Cold spot exists along $`(\pi ,\pi )`$ and hot spot along $`(0,\pi )`$. The coherent $`c`$-axis matrix element $`t_{}`$, however, is itself dependent on direction and $`t_{}^2`$ varies as $`\mathrm{cos}^4(2\varphi )`$, where $`\varphi `$ is the angle defining the in-plane direction of $`𝐤`$. This factor makes the $`c`$-axis conductivity sensitive mainly to the hot spot, although it remains Drude-like. Here we can ignore both $`\varphi `$ dependences and interprete $`\mathrm{\Gamma }`$ as the scattering rate coming from $`(0,\pi )`$ region (hot spot). However, except for YBCO at optimum doping, $`\sigma _1(\omega )`$ vs $`\omega `$ is found to be almost flat over an energy range of a few $`100`$ meV. This implies that $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }_s`$ need to be several $`100`$ meV. Non-Fermi liquid models with large $`\mathrm{\Gamma }`$ have appeared in the literature based on ARPES data in Bi$`2212`$. Norman et al. have derived a simple phenomenological model - called the mode model in which the normal state is described by a large frequency independent $`\mathrm{\Gamma }`$ and the superconducting state by $`\mathrm{\Gamma }(\omega )=\mathrm{\Gamma }\theta (\omega \mathrm{\Delta }(0)\omega _{mode})`$, i.e. a Fermi liquid is recovered in this state as $`\mathrm{\Gamma }(\omega )=0`$ for $`\omega <\mathrm{\Delta }(0)+\omega _{mode}`$, where $`\mathrm{\Delta }(0)`$ is the gap and $`\omega _{mode}`$ is the frequency of some collective mode. In addition, a pseudogap $`\stackrel{~}{\mathrm{\Delta }}`$ can be introduced in the normal state as in the preformed pair model with no long range phase coherence above $`T_c`$. We take the pseudogap to have the same $`d`$-wave symmetry as does the superconducting gap in agreement with experiment but it may not have the same amplitude. A possible phenomenological form of spectral function $`\stackrel{~}{A}(𝐤,\omega )`$ is $`\stackrel{~}{A}_+(𝐤,\omega )+\stackrel{~}{A}_{}(𝐤,\omega )`$, where $`\stackrel{~}{A}_\pm (𝐤,\omega )=\mathrm{\Gamma }(1\pm \xi _𝐤/\stackrel{~}{E}_𝐤)/\left[(\omega \stackrel{~}{E}_𝐤)^2+\mathrm{\Gamma }^2\right]`$. Here $`\stackrel{~}{E}_𝐤=\sqrt{\xi _𝐤^2+\stackrel{~}{\mathrm{\Delta }}_𝐤^2}`$ and the factors $`(1\pm \xi _𝐤/\stackrel{~}{E}_𝐤)`$ do not imply the phase coherence of BCS theory.
If the frequency cut-off $`\omega _c`$ is much larger than any energy scale in our consideration, then under the assumption of a cylindrical Fermi surface, it can be shown that as $`T0`$, $`T_{\omega ,𝐤}G(𝐤,\omega )^2=1/2+𝐊\left(i\mathrm{\Delta }(0)/\mathrm{\Gamma }_s\right)/\pi `$ and $`T_{\omega ,𝐤}F(𝐤,\omega )^2=1/2𝐊\left(i\mathrm{\Delta }(0)/\mathrm{\Gamma }_s\right)/\pi `$ for superconducting state with an in-plane scattering rate $`\mathrm{\Gamma }_s`$, which is assumed, for simplicity, to be frequency independent, and $`T_{\omega ,𝐤}G_0(𝐤,\omega )^2=1/2+𝐊\left(i\stackrel{~}{\mathrm{\Delta }}/\mathrm{\Gamma }\right)/\pi `$ for normal pseudogap state. Here $`𝐊`$ is the complete elliptic integral of the first kind and we can set the density of states $`N(0)`$ at the Fermi level equal to be unity because it does not appear in NMSW. The in-plane angle dependence of $`t_{}^2`$ such as $`\mathrm{cos}^4(2\varphi )`$ does not change the results as long as $`\omega _c\mathrm{\Delta }(0)`$ and $`\stackrel{~}{\mathrm{\Delta }}`$. Therefore, NMSW becomes
$$\frac{(N_nN_s)}{\rho _s}=\frac{1}{2}+\frac{𝐊\left(i\frac{\stackrel{~}{\mathrm{\Delta }}}{\mathrm{\Gamma }}\right)𝐊\left(i\frac{\mathrm{\Delta }(0)}{\mathrm{\Gamma }_s}\right)}{\pi 2𝐊\left(i\frac{\mathrm{\Delta }(0)}{\mathrm{\Gamma }_s}\right)}.$$
(4)
$`\mathrm{\Gamma }_s`$ may be different from $`\mathrm{\Gamma }`$ and much smaller as ARPES data imply. Such data are more consistent with the existence of quasiparticles in the superconducting state than in the normal state.
It is worthwhile illustrating the implications of Eq. (4) in some detail. i) Suppose the in-plane scattering rate $`\mathrm{\Gamma }_s`$ is negligible compared with $`\mathrm{\Delta }(0)`$, then Eq. (4) reduces to $`(N_nN_s)/\rho _s=1/2+𝐊\left(i\stackrel{~}{\mathrm{\Delta }}/\mathrm{\Gamma }\right)/\pi `$, which depends only on the normal state and NMSW $`1/2`$. In this case, the conventional sum rule is recovered when $`\stackrel{~}{\mathrm{\Delta }}/\mathrm{\Gamma }1`$ because the spectral function $`\stackrel{~}{A}(𝐤,\omega )`$ becomes insensitive to the pseudogap. ii) If $`\mathrm{\Delta }(0)/\mathrm{\Gamma }_s1`$ and $`\stackrel{~}{\mathrm{\Delta }}/\mathrm{\Gamma }1`$, then NMSW $`1/2`$. To get $`1/2`$, it is not necessary that $`\mathrm{\Delta }(0)\stackrel{~}{\mathrm{\Delta }}`$ and $`\mathrm{\Gamma }_s\mathrm{\Gamma }`$ simultaneously; in other words, a mismatch between $`\stackrel{~}{\mathrm{\Delta }}`$ and $`\mathrm{\Delta }(0)`$ does not matter as long as $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }_s`$ are both negligible. But a small value of $`\mathrm{\Gamma }`$ gives a frequency dependence of $`\sigma _{1c}^n(\omega )`$, which disagrees with the observation that it is flat for $`\omega `$ up to a few $`100`$ meV. Thus, the large $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }_s`$ limit must be applied, and consequently any pseudogap region at low $`\omega \stackrel{~}{\mathrm{\Delta }}`$ will be filled in and the model cannot describe the observation made in underdoped YBCO. A cancelation between $`_{\omega ,𝐤}G(𝐤,\omega )^2`$ and $`_{\omega ,𝐤}G_0(𝐤,\omega )^2`$ in Eq. (3) can still arise for large $`\mathrm{\Gamma }`$ if there is a match between $`\mathrm{\Delta }(0)`$ and $`\stackrel{~}{\mathrm{\Delta }}`$ and between $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }_s`$ as can be seen from Fig.1, where we plot $`T_{\omega ,𝐤}G_0(𝐤,\omega )^2`$ (or $`T_{\omega ,𝐤}G(𝐤,\omega )^2`$) for three values of the pseudogap (or gap) as a function of $`\mathrm{\Gamma }`$ (or $`\mathrm{\Gamma }_s`$). We also show $`T_{\omega ,𝐤}F(𝐤,\omega )^2`$. Clearly the large $`\mathrm{\Gamma }`$ region can also give a value of the sum rule bigger or smaller than $`1/2`$. Should the pseudogap be larger than the superconducting gap (point $`𝐛`$ in Fig. 1) and $`\mathrm{\Gamma }_s`$ not too much smaller than $`\mathrm{\Gamma }`$ (point $`𝐚`$), the sum rule will be less than $`1/2`$ while if $`\mathrm{\Gamma }_s`$ is much less than $`\mathrm{\Gamma }`$ such as for point $`𝐝`$, it will be larger than $`1/2`$ but less than $`1`$. Other combinations could also be used. To be more explicit, suppose that $`\stackrel{~}{\mathrm{\Delta }}/\mathrm{\Gamma }=\mathrm{\Delta }(0)/\mathrm{\Gamma }_s+\gamma `$ with $`\mathrm{\Delta }(0)/\mathrm{\Gamma }_s\gamma `$, then $`𝐊(i\stackrel{~}{\mathrm{\Delta }}/\mathrm{\Gamma })/\pi 𝐊(i\mathrm{\Delta }(0)/\mathrm{\Gamma }_s)/\pi (\gamma /4)\mathrm{\Delta }(0)/\mathrm{\Gamma }_s`$. Since $`1/2𝐊(i\mathrm{\Delta }(0)/\mathrm{\Gamma }_s)/\pi (1/8)(\mathrm{\Delta }(0)/\mathrm{\Gamma }_s)^2`$ for $`\mathrm{\Delta }(0)/\mathrm{\Gamma }_s1`$, NMSW$`1/2\gamma \mathrm{\Gamma }_s/\mathrm{\Delta }(0)`$ for $`\gamma \mathrm{\Delta }(0)/\mathrm{\Gamma }_s1`$. Based on ARPES data, we may assume that $`\stackrel{~}{\mathrm{\Delta }}\mathrm{\Delta }(0)`$ and $`\mathrm{\Gamma }\mathrm{\Gamma }_s0`$. In this case NMSW$`1/2+(𝐚𝐝)/(2𝐜)`$, which is obviously greater than $`1/2`$, but small $`\mathrm{\Gamma }_s`$ will not give the observed flat response $`\sigma _{1c}^s(\omega )`$ in the superconducting state.
So far we have assumed that both superconducting and normal state are at $`T0`$. Because it is not possible to access the normal state at low $`T`$, Basov et al. have used $`T=T_c`$ instead. In the inset of Fig. 1, we plot $`T_{\omega ,𝐤}G_0^2`$ with $`\stackrel{~}{\mathrm{\Delta }}=4T_c`$ and $`\mathrm{\Gamma }=T_c`$ as a function of $`T`$. (We emphasize that $`T_{\omega ,𝐤}G_0^2`$ is nothing but NMSW if $`\mathrm{\Gamma }_s\mathrm{\Delta }(0)`$.) It is clear that including the $`T`$ dependence in the NMWS does not qualitatively change the physics of the conductivity sum rule.
Incoherent coupling. For incoherent $`c`$-axis coupling, we need a specific model for the impurity scattering potential $`|V_{𝐤𝐩}|^2`$ and need to average over impurity configuration. Part of the disorder scattering can be due to a mismatch of overlap matrix elements $`t_{}`$ between planes. We use a simple model for the scattering potential $`|V_{𝐤𝐩}|^2=|V_0|^2+|V_1|^2\mathrm{cos}(2\varphi _k)\mathrm{cos}(2\varphi _p)`$ and assume $`|V_0|^2=|V_1|^2`$ for simplicity. The NMSW is now
$$\frac{(N_nN_s)}{\rho _s}=\frac{1}{2}+\frac{1}{2}\frac{\underset{\omega }{}\left[\stackrel{~}{\kappa }^2𝐊^2(\stackrel{~}{\kappa })\kappa ^2𝐊^2(\kappa )\right]}{_\omega \left[(\kappa ^2/\kappa )𝐊(\kappa )𝐄(\kappa )/\kappa \right]^2},$$
(5)
where $`\stackrel{~}{\kappa }=\stackrel{~}{\mathrm{\Delta }}/\sqrt{\stackrel{~}{\mathrm{\Delta }}^2+(\omega +\mathrm{\Gamma }\text{sgn}\omega )^2}`$, $`\stackrel{~}{\kappa }^{}=\sqrt{1\stackrel{~}{\kappa }^2}`$, $`\kappa =\mathrm{\Delta }(T)/\sqrt{\mathrm{\Delta }(T)^2+(\omega +\mathrm{\Gamma }_s\text{sgn}\omega )^2}`$, $`\kappa ^{}=\sqrt{1\kappa ^2}`$, and $`𝐄`$ is the complete elliptic integral of the second kind. As one can easily see, $`\stackrel{~}{\kappa }^2𝐊^2(\stackrel{~}{\kappa })`$ becomes $`(\pi /2)^2`$ if $`\stackrel{~}{\mathrm{\Delta }}=0`$ regardless of $`\mathrm{\Gamma }`$. It can also be inferred from Eq. (5) that i) as $`\stackrel{~}{\mathrm{\Delta }}`$ increases or $`\mathrm{\Gamma }`$ decreases, NMSW becomes smaller because $`\stackrel{~}{\kappa }`$ is closer to $`1`$ for a given $`\omega `$, and ii) as $`\mathrm{\Gamma }_s`$ increases NMSW becomes smaller because $`\kappa `$ is closer to $`0`$. In the incoherent coupling case, the normal state response $`\sigma _{1c}^n|V_0|^2`$ is independent of $`\omega `$ and there is no restriction on $`\mathrm{\Gamma }`$. In the inset of Fig. 2, we plot NMSW as a function of $`\mathrm{\Gamma }_s`$ with fixed values of $`\alpha =0`$ and $`1`$, where $`\alpha \stackrel{~}{\mathrm{\Delta }}/\mathrm{\Delta }(0)`$, and $`\mathrm{\Gamma }=\mathrm{\Delta }(0)`$. As expected, NMSW with $`\alpha =1`$ decreases with increasing $`\mathrm{\Gamma }_s`$. Note that NMSW with $`\alpha =0`$ is almost unchanged with increasing $`\mathrm{\Gamma }_s`$ because the numerator and the denominator of the second term in Eq. (5) are both decreasing in similar manner. In the case $`\stackrel{~}{\mathrm{\Delta }}=\mathrm{\Delta }(0)(1+\delta )`$ $`\left[\delta 1\right]`$ with $`\mathrm{\Gamma }=\mathrm{\Gamma }_s=0`$, Eq. (5) reduces to NMSW$`1/21.083\delta `$, which is less than $`1/2`$. For $`\mathrm{\Gamma }=\mathrm{\Gamma }_s0`$, the coefficient of $`\delta `$ will be changed. In the general case we need to evaluate the sums in Eq. 5 for given $`\stackrel{~}{\mathrm{\Delta }}`$, $`\mathrm{\Delta }(0)`$, $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }_s`$. In the main frame of Fig. 2, we plot NMSW as a function of $`\mathrm{\Gamma }`$ with a few values of $`\alpha `$ assuming $`\mathrm{\Gamma }_s=0`$ since we have seen that $`\mathrm{\Gamma }_s0`$ decreases NMSW. As one sees, NMSW is greater than 1 and independent of $`\mathrm{\Gamma }`$ if $`\alpha =0`$. Note that NMSW with $`\alpha =0`$ is an asymptote of NMSW with nonzero $`\alpha `$ because the effect of $`\stackrel{~}{\mathrm{\Delta }}`$ on NMSW vanishes when $`\mathrm{\Gamma }\stackrel{~}{\mathrm{\Delta }}`$. As we expect, NMSW decreases with increasing $`\alpha `$ and decreasing $`\mathrm{\Gamma }`$. If $`\alpha 1.46`$ and $`\mathrm{\Gamma }\mathrm{\Delta }(0)`$, then NMSW$`<0`$. Generally speaking, in the incoherent coupling model, we can explain NMSW$`1/2`$ with $`\stackrel{~}{\mathrm{\Delta }}\mathrm{\Delta }(0)`$, $`\mathrm{\Gamma }>\mathrm{\Delta }(0)`$ and $`\mathrm{\Gamma }_s<\mathrm{\Delta }(0)`$.
Additional insight into the sum rule can be obtained from consideration of $`\sigma _{1c}^{n(s)}(\omega )`$. In the top panel of Fig. 3, we compare $`\sigma _{1c}^n(\omega )`$ (solid curve), with pseudogap $`\stackrel{~}{\mathrm{\Delta }}=\mathrm{\Delta }(0)`$, with $`\sigma _{1c}^s(\omega )`$ (dashed curve). In both cases, $`\mathrm{\Gamma }`$ is taken to be zero and therefore both curves go to zero at $`\omega =0`$. At small $`\omega \mathrm{\Delta }(0)`$ the solid and dashed curves are almost the same. However, they start to differ at a higher $`\omega `$, with $`\sigma _{1c}^s(\omega )`$ falling below $`\sigma _{1c}^n(\omega )`$. We already know that for these parameter the sum rule is exactly $`1/2`$. But in this example, the effective NMSW defined as $`𝒩(\omega )=\left[N_n(\omega )N_s(\omega )\right]/\rho _s`$ is nealy zero up to $`\omega \mathrm{\Delta }(0)`$ as shown in the dashed curve in the bottom panel of Fig. 3. This does not agree with the findings of Basov et al., where $`𝒩(\omega )`$ increases almost linearly out of zero and rapidly move toward its saturated value. It is clear that, even for the incoherent $`c`$-axis coupling case, the argument presented in Ref. that the preformed pair model agrees with a sum rule of $`1/2`$ due to a simple cancelation between the normal and superconducting state is deficient when its detailed approach toward its saturated value is considered, i.e. when we compare the low $`\omega `$ behavior of $`𝒩(\omega )`$ with experiment.
A more reasonable model is obtained for finite normal state $`\mathrm{\Gamma }`$. Returning to the solid curve of the top panel and including a finite $`\mathrm{\Gamma }`$ in the pseudogap state fills in the region below $`2\stackrel{~}{\mathrm{\Delta }}`$ by transfering spectral weight out of the region above this. In the middle panel we show results (solid curve) for a pseudogap $`\stackrel{~}{\mathrm{\Delta }}=1.4\mathrm{\Delta }(0)`$ and $`\mathrm{\Gamma }=0.5\mathrm{\Delta }(0)`$ for illustration only. Other choices could also have been made showing that our explanations are robust. Comparison of the solid curve (normal state) with the superconducting (dashed) curve $`(\mathrm{\Gamma }_s=0)`$ shows a very different low $`\omega `$ behavior. Now $`𝒩(\omega )`$ will grow rapidly as $`\omega `$ increases from $`\omega =0`$. This behavior is shown as the solid curve of the bottom panel. $`𝒩(\omega )`$ grows linearly out of $`\omega =0`$ and nearly reaches its saturated value$`(0.52)`$ by $`\omega 67\mathrm{\Delta }(0)`$ in much better agreement with experiment. One can also conceive that $`\stackrel{~}{\mathrm{\Delta }}`$ and $`\mathrm{\Gamma }`$ are larger than considered above and that the normal state may not exhibit much of a pseudogap because of a large $`\mathrm{\Gamma }`$, but that the sum rule is still nevertheless less than 1. This could explain optimally doped Tl$`2201`$ with NMSW$`0.6`$ where, however, no in-plane pseudogap has been reported although in the data of Basov et al., $`\sigma _{1c}^n(\omega )`$ is droping with decreasing $`\omega `$ at small $`\omega `$ consistent with a smeared pseudogap.
Conclusion. We conclude that incoherent $`c`$-axis coupling can more easily explain the observed violations of the optical sum rule than can coherent coupling, but that non-Fermi liquid in-plane behavior is required. It is possible that the normalized missing spectral weight is less than $`1/2`$. A sharp increase in accumulated normalized spectral weight as a function of cut-off frequency rising to a value close to its saturated value within a few times the gap as is observed can also be understood within our phenomenological model but not within a pure preformed pair model.
Work supported by the Natural Sciences and Engineering Research Council of Canada (NSERC). W. K. acknowledges N. Whelan and P. J. Hirschfeld for discussions and correspondence, respectively. J. P. C thanks D. Basov for enlightening discussions.
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