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# 1 Introduction
## 1 Introduction
The $`q`$-state Potts antiferromagnet (AF) exhibits nonzero ground state entropy, $`S_0>0`$ (without frustration) for sufficiently large $`q`$ on a given lattice $`\mathrm{\Lambda }`$ or, more generally, on a graph $`G=(V,E)`$ defined by its set of vertices $`V`$ and edges joining these vertices $`E`$. This is equivalent to a ground state degeneracy per site $`W>1`$, since $`S_0=k_B\mathrm{ln}W`$. Such nonzero ground state entropy is important as an exception to the third law of thermodynamics . There is a close connection with graph theory here, since the zero-temperature partition function of the above-mentioned $`q`$-state Potts antiferromagnet on a graph $`G`$ satisfies
$$Z(G,q,T=0)_{PAF}=P(G,q)$$
(1.1)
where $`P(G,q)`$ is the chromatic polynomial expressing the number of ways of coloring the vertices of the graph $`G`$ with $`q`$ colors such that no two adjacent vertices have the same color (for reviews, see -). The minimum number of colors necessary for such a coloring of $`G`$ is called the chromatic number, $`\chi (G)`$. Thus
$$W(\{G\},q)=\underset{n\mathrm{}}{lim}P(G,q)^{1/n}$$
(1.2)
where $`n=|V|`$ is the number of vertices of $`G`$ and $`\{G\}=lim_n\mathrm{}G`$. At certain special points $`q_s`$ (typically $`q_s=0,1,..,\chi (G)`$), one has the noncommutativity of limits
$$\underset{qq_s}{lim}\underset{n\mathrm{}}{lim}P(G,q)^{1/n}\underset{n\mathrm{}}{lim}\underset{qq_s}{lim}P(G,q)^{1/n}$$
(1.3)
and hence it is necessary to specify the order of the limits in the definition of $`W(\{G\},q_s)`$ . Denoting $`W_{qn}`$ and $`W_{nq}`$ as the functions defined by the different order of limits on the left and right-hand sides of (1.3), we take $`WW_{qn}`$ here; this has the advantage of removing certain isolated discontinuities that are present in $`W_{nq}`$.
Using the expression for $`P(G,q)`$, one can generalize $`q`$ from $`_+`$ to $``$. The zeros of $`P(G,q)`$ in the complex $`q`$ plane are called chromatic zeros; a subset of these may form an accumulation set in the $`n\mathrm{}`$ limit, denoted $``$ , which is the continuous locus of points where $`W(\{G\},q)`$ is nonanalytic. <sup>1</sup><sup>1</sup>1For some families of graphs $``$ may be null, and $`W`$ may also be nonanalytic at certain discrete points. The maximal region in the complex $`q`$ plane to which one can analytically continue the function $`W(\{G\},q)`$ from physical values where there is nonzero ground state entropy is denoted $`R_1`$. The maximal value of $`q`$ where $``$ intersects the (positive) real axis is labelled $`q_c(\{G\})`$. This point is important since it separates the interval $`q>q_c(\{G\})`$ on the positive real $`q`$ axis where the Potts model (with $`q`$ extended from $`_+`$ to $``$) exhibits nonzero ground state entropy (which increases with $`q`$, asymptotically approaching $`S_0=k_B\mathrm{ln}q`$ for large $`q`$, and which for a regular lattice $`\mathrm{\Lambda }`$ can be calculated approximately via large–$`q`$ series expansions) from the interval $`0qq_c(\{G\})`$ in which $`S_0`$ has a different analytic form. Early calculations of chromatic polynomials for $`L_y=2`$ strips of the square lattice with periodic longitudinal boundary conditions were performed in (see also the related works -).
Here we present exact calculations of the chromatic polynomials for strips of the square and triangular lattice with transverse width $`L_y=4`$ (i.e. transverse cross sections forming squares) and arbitrarily great length $`L_x`$ with the following boundary conditions: (i) $`(PBC_y,PBC_x)=`$ toroidal, and (ii) $`(PBC_y,TPBC_x)=`$ Klein bottle, where $`PBC_i`$ denotes periodic boundary conditions in the $`i`$’th direction and $`TPBC_x`$ denotes periodic longitudinal boundary conditions with an orientation-reversal (twist).<sup>2</sup><sup>2</sup>2The boundary conditions $`(PBC_y,PBC_x)`$ and $`(PBC_y,TPBC_x)`$ can be implemented in a manner that is uniform in the length $`L_x`$; as noted before , the boundary conditions $`(TPBC_y,PBC_x)`$ (different type of Klein bottle) and $`(TPBC_y,TPBC_x)`$ (projective plane) require different identifications as $`L_x`$ varies and will not be considered here. These extend our previous calculations of chromatic polynomials for width $`L_y=3`$ on the square and triangular lattices with torus and Klein bottle boundary conditions.
A major motivation for using boundary conditions that are fully periodic or fully periodic with reversed orientation (here, toroidal and Klein bottle) is the well-known fact that if one imposes periodic boundary conditions in a certain direction, this removes edge effects in that direction. Clearly the most complete removal of such edge effects is achieved if one imposes fully periodic boundary conditions (including the possibility of orientation reversal). This also has an important related consequence pertaining to the uniformity of the lattice. To discuss this, we first recall two definitions from mathematical graph theory. The degree $`\mathrm{\Delta }`$ of a vertex of a graph is the number of edges connected to it. A $`\mathrm{\Delta }`$-regular graph is a graph in which all vertices have the same degree, $`\mathrm{\Delta }`$. An infinite regular lattice has the property that each vertex (site) on the lattice has the same degree, i.e., coordination number. For the two types of lattices considered here, namely square and triangular, the coordination number is 4 and 6, respectively. It is advantageous to deal with finite sections of regular lattices having boundary conditions that preserve the $`\mathrm{\Delta }`$-regular property of the infinite lattice. Fully periodic periodic boundary conditions, and the reversed-orientation periodic boundary conditions considered here, have the merit of preserving this property of $`\mathrm{\Delta }`$-regularity; in contrast, this is not the case if one uses boundary conditions that are free in one or more directions. In previous studies with families of lattice strip graphs of arbitrarily great length with periodic or reversed-orientation periodic longitudinal boundary conditions and free transverse boundary conditions (i.e., cyclic or Möbius strips), it was shown that, in the $`L_x\mathrm{}`$ limit, the resultant locus $``$ exhibits, for finite width $`L_y`$, a number of properties expected to hold for the locus $``$ on the infinite 2D lattice, including (i) passing through $`q=0`$, (ii) passing through $`q=2`$, (iii) passing through a maximal real point, $`q_c`$, and (iv) enclosing one or more regions including the interval $`0<q<q_c`$ , -. In contrast, if one uses free longitudinal boundary conditions, it was found in - that properties (i) and (iv) do not hold, and properties (ii) and (iii) do not, in general, hold; rather, one anticipates that these would be approached in the limit $`L_y\mathrm{}`$. It was thus inferred that the key condition to guarantee that these properties hold is the presence of periodic (or reversed-orientation periodic) longitudinal boundary conditions . This thus provides a third motivation for calculations with doubly periodic boundary conditions, since one expects that the resultant loci $``$ will exhibit the features (i)-(iv) already for finite $`L_y`$, and this was confirmed by the study of $`L_y=3`$ strips of the square and triangular lattices. As will be seen, our exact results for $`L_y=4`$ again support this inference. A fourth motivation for this study is that, as was shown in the earlier calculations of chromatic polynomials for strips of the square and triangular lattices with width $`L_y=3`$ and is again true for width $`L_y=4`$, the use of Klein bottle, as opposed to torus, boundary conditions has the effect of simplifying the structure of the resultant chromatic polynomial. This thus elucidates the effect of the topology of the surface on which the family of strip graphs is embedded with the structure of the chromatic polynomial. In addition to those listed, some previous related calculations of chromatic polynomials for families of graphs with periodic longitudinal boundary conditions are in Refs. -.
In general, the $`L_y\times L_x`$ strips of the square and triangular lattice have $`n=|V|=L_yL_x`$ vertices and, for the number of edges $`|E|=(\mathrm{\Delta }/2)n`$ the values $`|E|=2n`$ and $`|E|=3n`$ respectively. (For $`L_x=2`$, some of these strip graphs involve multiple edges joining pairs of vertices and hence are multigraphs rather than proper graphs; we shall be interested primarily in the cases $`L_x3`$ where there are no multiple edges.)
We label a particular type of strip graph as $`G_s`$ or just $`G`$ and the specific graph of width $`L_y`$ and length $`L_x`$ vertices as $`(G_s,L_y\times L_x,BC_y,BC_x)`$. A generic form for chromatic polynomials for recursively defined families of graphs, of which strip graphs $`G_s`$ are special cases, is
$$P(G_s,L_y\times L_x,BC_y,BC_x,q)=\underset{j=1}{\overset{N_{G_s,\lambda }}{}}c_{G_s,j}(q)(\lambda _{G_s,j}(q))^m$$
(1.4)
where $`c_{G_s,j}(q)`$ and the $`N_{G_s,\lambda }`$ terms $`\lambda _{G_s,j}(q)`$ depend on the type of strip graph $`G_s`$ but are independent of $`m`$. The $`\lambda _{G_s,j}`$ are the (nonzero) eigenvalues of the coloring matrix . We shall denote the total number of different eigenvalues of the coloring matrix for a recursive family of graphs $`G_s`$ as $`N_{G_s,\lambda ,tot}`$. Clearly $`N_{G_s,\lambda ,tot}=N_{G_s,\lambda }`$ if there is no zero eigenvalue, and $`N_{G_s,\lambda ,tot}=N_{G_s,\lambda }+1`$ if there is a zero eigenvalue. Our results illustrate both of these possibilities.
For a given type of strip graph $`G_s`$, we denote the sum of the coefficients $`c_{G_s,j}`$ as
$$C_{G_s}C(G_s)=\underset{j=1}{\overset{N_{G_s,\lambda }}{}}c_{G_s,j}.$$
(1.5)
According to a general theorem, for a strip $`G_s`$ of the square or triangular lattice with torus boundary conditions ,
$$C(G_s,L_y\times L_x,PBC_y,PBC_x)=P(C_{L_y},q),G_s=sq,tri$$
(1.6)
where $`C_n`$ denotes the circuit graph with $`n`$ vertices and $`P(C_n,q)=(q1)^n+(q1)(1)^n`$. Further, for a strip of the square or triangular lattice with Klein bottle boundary conditions
$$C(G_s,L_y\times L_x,PBC_y,TPBC_x)=0,G_s=sq,tri.$$
(1.7)
## 2 $`L_y=4`$ Strip of the Square Lattice with $`(PBC_y,PBC_x)`$
In general, for a strip of the square lattice of size $`L_y\times L_x`$ with $`(PBC_y,PBC_x)`$, i.e., toroidal boundary conditions, for $`L_y2`$ and $`L_x2`$, the chromatic number is given by
$$\chi (sq,L_y\times L_x,PBC_y,PBC_x)=\{\begin{array}{cc}2\hfill & \text{if }L_y\text{ is even and }L_x\text{ is even}\hfill \\ 3\hfill & \text{otherwise}\hfill \end{array}$$
(2.1)
Thus, in the present case with $`L_y=4`$, it follows that $`\chi =2`$ for even $`L_x`$ and $`\chi =3`$ for odd $`L_x`$. We calculate the chromatic polynomial $`P`$ by a systematic, iterative use of the deletion-contraction theorems as in our earlier work and a coloring matrix method . For the $`L_y=4`$ strip graphs of the square lattice with torus boundary conditions (labelled $`st4`$), we find $`N_{st4,\lambda }=33`$ and
$$P(sq,4\times L_x,PBC_y,PBC_x,q)=\underset{j=1}{\overset{33}{}}c_{st4,j}(\lambda _{st4,j})^{L_x}$$
(2.2)
where
$$\lambda _{st4,1}=1$$
(2.3)
$$\lambda _{st4,2}=1q$$
(2.4)
$$\lambda _{st4,3}=2q$$
(2.5)
$$\lambda _{st4,4}=3q$$
(2.6)
$$\lambda _{st4,5}=4q$$
(2.7)
$$\lambda _{st4,6}=5q$$
(2.8)
$$\lambda _{st4,7}=q^25q+5$$
(2.9)
$$\lambda _{st4,8}=q^25q+7$$
(2.10)
$$\lambda _{st4,9}=(q1)(q3)$$
(2.11)
$`\lambda _{st4,(10,11)}={\displaystyle \frac{1}{2}}[q^48q^3+29q^255q+46`$ (2.12)
(2.13)
$`\pm (q^816q^7+118q^6526q^5+1569q^43250q^3+4617q^24136q+1776)^{1/2}]`$ (2.14)
(2.15)
(2.16)
$`\lambda _{st4,(12,13)}={\displaystyle \frac{1}{2}}[(q^37q^2+18q17)`$ (2.17)
(2.18)
$`\pm (q^614q^5+81q^4250q^3+442q^2436q+193)^{1/2}]`$ (2.19)
$$\lambda _{st4,(14,15)}=\frac{1}{2}\left[q^27q+9\pm \left(q^410q^3+35q^250q+33\right)^{1/2}\right]$$
(2.20)
$$\lambda _{st4,(16,17)}=\frac{1}{2}\left[(q3)^2\pm \left(q^48q^3+26q^248q+41\right)^{1/2}\right]$$
(2.21)
$$\lambda _{st4,(18,19)}=\frac{1}{2}\left[q^26q+11\pm (q3)\left(q^22q+9\right)^{1/2}\right]$$
(2.22)
$$\lambda _{st4,(20,21)}=3q\pm \sqrt{3}.$$
(2.23)
The remaining twelve $`\lambda _{st4,j}`$’s for $`22j33`$ are roots of four cubic equations,
$`\xi ^3+(q^36q^2+16q14)\xi ^2(q1)(q^49q^3+31q^255q+43)\xi `$ (2.24)
(2.25)
$`(q3)(q1)^2(q^36q^2+12q10)=0`$ (2.26)
with roots $`\lambda _{st4,j}`$ for $`j=22,23,24`$,
$`\xi ^3+(q4)(q^26q+12)\xi ^2(q3)(q^411q^3+45q^281q+59)\xi `$ (2.27)
(2.28)
$`(q^615q^5+91q^4285q^3+488q^2442q+170)=0`$ (2.29)
with roots $`\lambda _{st4,j}`$ for $`j=25,26,27`$,
$`\xi ^32(q^26q+12)\xi ^2+(q^413q^3+59q^2113q+83)\xi `$ (2.30)
(2.31)
$`+(q^513q^4+62q^3135q^2+141q60)=0`$ (2.32)
with roots $`\lambda _{st4,j}`$ for $`j=28,29,30`$, and
$`\xi ^32(q^26q+10)\xi ^2+(q^413q^3+59q^2113q+75)\xi `$ (2.33)
(2.34)
$`+(q^513q^4+64q^3149q^2+167q72)=0`$ (2.35)
with roots $`\lambda _{st4,j}`$ for $`j=31,32,33`$.
The corresponding coefficients are
$$c_{st4,1}=q^48q^3+20q^215q+1$$
(2.36)
$$c_{st4,2}=\frac{1}{2}c_{st4,4}=c_{st4,6}=\frac{1}{3}(q1)(q^25q+3)$$
(2.37)
$$c_{st4,3}=\frac{1}{6}(q2)(4q^213q3)$$
(2.38)
$$c_{st4,j}=\frac{2}{3}q(q2)(q4)\mathrm{for}j=5,20,21$$
(2.39)
$$c_{st4,j}=\frac{1}{2}(q1)(q2)\mathrm{for}j=7,18,19$$
(2.40)
$$c_{st4,j}=(q1)(q2)\mathrm{for}j=31,32,33$$
(2.41)
$$c_{st4,j}=\frac{1}{2}q(q3)\mathrm{for}j=8,14,15,28,29,30$$
(2.42)
$$c_{st4,j}=q(q3)\mathrm{for}j=16,17$$
(2.43)
$$c_{st4,j}=1\mathrm{for}j=9,10,11$$
(2.44)
$$c_{st4,j}=2(q1)\mathrm{for}j=12,13$$
(2.45)
$$c_{st4,j}=q1\mathrm{for}22j27.$$
(2.46)
The sum of these coefficients is equal to $`P(C_4,q)=q(q1)(q^23q+3)`$, as dictated by the $`L_y=4`$ special case of our general result (1.6).
The singular locus $``$ for the $`L_x\mathrm{}`$ limit of the strip of the square lattice with $`L_y=4`$ and toroidal boundary conditions is shown in Fig. 1. For comparison, chromatic zeros are calculated and shown for length $`L_x=30`$ (i.e., $`n=120`$ vertices). The locus $``$ crosses the real axis at the points $`q=0`$, $`q=2`$, and at the maximal point $`q=q_c`$, where
$$q_c=2.7827657\mathrm{}\mathrm{for}\{G\}=(sq,4\times \mathrm{},PBC_y,PBC_x).$$
(2.47)
As is evident from Fig. 1, the locus $``$ separates the $`q`$ plane into different regions including the following: (i) $`R_1`$, containing the semi-infinite intervals $`q>q_c`$ and $`q<0`$ on the real axis and extending outward to infinite $`|q|`$; (ii) $`R_2`$ containing the interval $`2<q<q_c`$; (iii) $`R_3`$ containing the real interval $`0<q<2`$; and (iv) the complex-conjugate pair $`R_4,R_4^{}`$ centered approximately at $`q=2.9\pm 1.3i`$. The (nonzero) density of chromatic zeros has the smallest values on the curve separating regions $`R_1`$ and $`R_3`$ in the vicinity of the point $`q=0`$ and on the curve separating regions $`R_2`$ and $`R_3`$ in the vicinity of the point $`q=2`$.
In region $`R_1`$, $`\lambda _{st4,10}`$ is the dominant $`\lambda _{G,j}`$, so
$$W=(\lambda _{st4,10})^{1/4},qR_1.$$
(2.48)
This is the same as $`W`$ for the corresponding $`L_x\mathrm{}`$ limit of the strip of the square lattice with the same width $`L_y=4`$ and cylindrical $`(PBC_y,FBC_x)`$ boundary conditions, calculated in . This equality of the $`W`$ functions for the $`L_x\mathrm{}`$ limit of two strips of a given lattice with the same transverse boundary conditions and different longitudinal boundary conditions in the more restrictive region $`R_1`$ defined by the two boundary conditions is a general result .
In region $`R_2`$, the largest root of the cubic equation (2.32) is dominant; we label this as $`\lambda _{st4,28}`$ so that
$$|W|=|\lambda _{st4,28}|^{1/4},qR_2$$
(2.49)
(in regions other than $`R_1`$, only $`|W|`$ can be determined unambiguously ). Thus, $`q_c`$ is the relevant solution of the equation of degeneracy in magnitude $`|\lambda _{st4,10}|=|\lambda _{st4,28}|`$. In region $`R_3`$,
$$|W|=|\lambda _{st4,25}|^{1/4},qR_3.$$
(2.50)
In regions $`R_4,R_4^{}`$,
$$|W|=|\lambda _{st4,22}|^{1/4},qR_4,R_4^{}.$$
(2.51)
In we have listed values of $`W`$ for a range of values of $`q`$ for the $`L_x\mathrm{}`$ limit of various strips of the square lattice, including $`(sq,4\times \mathrm{},PBC_y,FBC_x)`$. Since $`W`$ is independent of $`BC_x`$ for $`q`$ in the more restrictive region $`R_1`$ defined by $`FBC_x`$ and $`(T)PBC_x`$ (which is the $`R_1`$ defined by $`PBC_x`$ here), it follows, in particular, that
$$W(4\times \mathrm{},PBC_y,(T)PBC_x,q)=W(4\times \mathrm{},PBC_y,FBC_x,q)\mathrm{for}qq_c$$
(2.52)
where $`q_c`$ was given above in (2.47). For low integral values of $`q`$ we list the values of $`|W(q)|`$ for this strip in Table 1, together with corresponding values given in for $`W`$ in the $`L_x\mathrm{}`$ limit of the $`L_y=3`$ strip with $`(PBC_y,(T)PBC_x)`$.
For various lengths $`L_x`$, some of the chromatic zeros (those near to the origin) have support for $`Re(q)<0`$, but the locus $``$ itself only has support for $`Re(q)0`$. We have encountered this type of situation in earlier work . The property that $``$ only has support for $`Re(q)0`$ can be demonstrated by carrying out a Taylor series expansion of the degeneracy equation $`|\lambda _{st4,10}|=|\lambda _{st4,25}|`$ near the origin, which is, numerically,
$$|44.152.0q+27.6q^2+O(q^3)|=|44.132.2q+8.8q^2+O(q^3)|.$$
(2.53)
More generally, consider a degeneracy equation determining a curve on $``$ which, in the vicinity of the origin $`q=0`$, has the form
$$|a_0+a_1q+a_2q^2+O(q^3)|=|a_0+b_1q+b_2q^2+O(q^3)|$$
(2.54)
where the coefficients $`a_i`$ and $`b_i`$ are real and nonzero, $`a_1b_1`$, and, without loss of generality, we can take $`a_0>0`$. Writing $`q`$ in polar coordinates as $`q=re^{i\theta }`$ and expanding for small $`r`$, eq. (2.54) reduces, to order $`r`$, to the equation $`a_0(a_1b_1)r\mathrm{cos}\theta =0`$, which has as its solution $`\theta =\pm \pi /2`$. Thus the curve $``$ defined by a degeneracy equation of the form (2.54) passes through the origin vertically. In order to determine in which direction (right or left) the curve bends away from the vertical as one moves away from the origin, let us write $`q=q__R+iq__I`$ where $`q__R`$ and $`q__I`$ are real, with $`q__R^2+q__I^2=r^21`$ Substituting, expanding, and using the fact that the curve $``$ passes vertically through the origin so that near this point $`|q_R|`$ is small compared with $`|q_I|`$, we find, to this order,
$$q__R=\frac{[2a_0(a_2b_2)+b_1^2a_1^2]q__I^2}{2a_0(a_1b_1)}.$$
(2.55)
Thus if the right-hand side of this equation is positive (negative), the curve $``$ bends to the right (left) into the half-plane with $`Re(q)>0`$ ($`Re(q)<0`$) as one moves away from the origin. For the degeneracy equation (2.53), the right-hand side of eq. (2.55) is positive, so $``$ bends to the right near the origin. As is evident from Fig. 1, as one moves farther away from the origin, the curve $``$ bends farther to the right, so that $``$ has no support for $`Re(q)<0`$. This is to be contrasted with the situation for (the $`L_x\mathrm{}`$ limit of) sufficiently wide strips with cyclic or Möbius boundary conditions (the $`L_x\mathrm{}`$ limits of a given strip with cyclic boundary conditions is the same as the limit with Möbius boundary conditions), where it was found that for widths $`L_y=3,4`$ for the square lattice and for width $`L_y=4`$ for the triangular lattice , $``$ did have some support for $`Re(q)<0`$. A comparison of some properties of $``$ in the present case and for other strips with periodic or orientation-reversing periodic longitudinal boundary conditions is given in Table 2.
## 3 $`L_y=4`$ Strip of the Square Lattice with $`(PBC_y,TPBC_x)`$
In general, for the strip graph of the square lattice with even width $`L_y`$ and $`(PBC_y,TPBC_y)`$, i.e., Klein bottle boundary conditions, we find that $`\chi =4`$ if $`L_x=2`$ and, for $`L_x3`$,
$$\chi (sq,L_y\times L_x,PBC_y,TPBC_x)=\{\begin{array}{cc}2\hfill & \text{if }L_x\text{ is odd}\hfill \\ 3\hfill & \text{if }L_x\text{ is even}\hfill \end{array}$$
(3.1)
For this strip (labelled $`sk4`$) we calculate that $`N_{sk4,\lambda }=22`$ and
$$P(sq,4\times L_x,PBC_y,TPBC_x,q)=\underset{j=1}{\overset{22}{}}c_{sk4,j}.(\lambda _{sk4,j})^{L_x}$$
(3.2)
The nonzero terms $`\lambda _{sk4,j}`$ are identical to a subset of the terms $`\lambda _{st4,j}`$’s for the same strip with torus boundary conditions. The 11 terms that occur in the chromatic polynomial (2.2) for toroidal boundary conditions but are absent in the chromatic polynomial (3.2) for the Klein bottle case are
$$\lambda _{st4,j},j=4,5,12,13,16,17,20,21,31,32,33.$$
(3.3)
We have
$$\lambda _{sk4,j}=\lambda _{st4,j}\mathrm{for}1j3$$
(3.4)
$$\lambda _{sk4,j}=\lambda _{st4,j+2}\mathrm{for}4j9$$
(3.5)
$$\lambda _{sk4,j}=\lambda _{st4,j+4}\mathrm{for}j=10,11$$
(3.6)
$$\lambda _{sk4,j}=\lambda _{st4,j+6}\mathrm{for}j=12,13$$
(3.7)
$$\lambda _{sk4,j}=\lambda _{st4,j+8}\mathrm{for}14j22.$$
(3.8)
The corresponding coefficients are
$$c_{sk4,1}=1$$
(3.9)
$$c_{sk4,2}=q1$$
(3.10)
$$c_{sk4,3}=\frac{1}{2}(q1)(q2)$$
(3.11)
$$c_{sk4,4}=(q1)$$
(3.12)
$$c_{sk4,5}=c_{st4,7}=\frac{1}{2}(q1)(q2)$$
(3.13)
$$c_{sk4,6}=c_{st4,8}=\frac{1}{2}q(q3)$$
(3.14)
$$c_{sk4,7}=c_{st4,9}=1$$
(3.15)
$$c_{sk4,j}=c_{st4,j+2}=1\mathrm{for}j=8,9$$
(3.16)
$$c_{sk4,j}=c_{st4,j+4}=\frac{1}{2}q(q3)\mathrm{for}j=10,11$$
(3.17)
$$c_{sk4,j}=c_{st4,j+6}=\frac{1}{2}(q1)(q2)\mathrm{for}j=12,13$$
(3.18)
$$c_{sk4,j}=c_{st4,j+8}=(q1)\mathrm{for}14j16$$
(3.19)
$$c_{sk4,j}=c_{st4,j+8}=q1\mathrm{for}17j19$$
(3.20)
$$c_{sk4,j}=c_{st4,j+8}=\frac{1}{2}q(q3)\mathrm{for}20j22.$$
(3.21)
The sum of these coefficients is zero, as dictated by the $`L_y=4`$ special case of the general result (1.7) above.
Because none of the terms $`\lambda _{st4,j}`$ in (3.3) that is present in (2.2) and absent in (3.2) is dominant, it follows that in the limit $`L_x\mathrm{}`$, the $`W`$ functions are the same for both of these boundary conditions, and hence, so is the singular locus $``$. Below we shall prove in general that this must be the case; that is, in the limit $`L_x\mathrm{}`$, a strip of the square (or triangular) lattice of width $`L_y`$ with $`(PBC_y,PBC_x)`$ (torus) boundary conditions yields the same $`W`$ function and singular locus $``$ as the corresponding strip with $`(PBC_y,TPBC_x)`$ (Klein bottle) boundary conditions.
## 4 $`L_y=4`$ Strip of the Triangular Lattice with $`(PBC_y,PBC_x)`$
By similar methods, we have calculated the chromatic polynomials for strips of the triangular lattice with width $`L_y=4`$, arbitrarily great length $`L_x`$, and torus boundary conditions (labelled $`tt4`$). In general, for a strip of the triangular lattice of size $`L_y\times L_x`$ with toroidal boundary conditions, for $`L_y3`$ and $`L_x3`$, the chromatic number is given by
$$\chi (tri,L_y\times L_x,PBC_y,PBC_x)=\{\begin{array}{cc}3\hfill & \text{if }L_y=0\text{ mod 3 and }L_x=0\text{ mod 3}\hfill \\ 4\hfill & \text{otherwise}\hfill \end{array}$$
(4.1)
Thus, in the present case, $`\chi =4`$, independent of $`L_x`$. In the notation of eq. (1.4) we find $`N_{tt4,\lambda }=37`$ and
$$P(tri,4\times L_x,PBC_y,PBC_x,q)=\underset{j=1}{\overset{37}{}}c_{tt4,j}(\lambda _{tt4,j})^{L_x}$$
(4.2)
where
$$\lambda _{tt4,1}=2$$
(4.3)
$$\lambda _{tt4,(2,3)}=\sqrt{2}e^{\pm i\pi /4}$$
(4.4)
$$\lambda _{tt4,4}=2(3q)$$
(4.5)
$$\lambda _{tt4,5}=3q$$
(4.6)
$$\lambda _{tt4,6}=2(2q9)$$
(4.7)
$$\lambda _{tt4,7}=2(q3)^2$$
(4.8)
$`\lambda _{tt4,(8,9)}`$ $`=`$ $`{\displaystyle \frac{(q3)}{2}}[q^39q^2+33q48`$ (4.11)
$`\pm (q4)(q^410q^3+43q^2106q+129)^{1/2}]`$
$$\lambda _{tt4,(10,11)}=\pm i\sqrt{3}(q3)$$
(4.12)
$$\lambda _{tt4,(12,13)}=\pm i(q2)\sqrt{2(q3)(q4)}.$$
(4.13)
The $`\lambda _{tt4,j}`$’s for $`14j19`$ are roots of two cubic equations,
$`\xi ^3+2(q^312q^2+51q75)\xi ^2`$ (4.14)
(4.15)
$`4(q3)^3(q^27q+13)\xi 8(q3)^4(q^25q+5)=0`$ (4.16)
with roots $`\lambda _{tt4,j}`$, $`j=14,15,16`$, and
$`\xi ^32(2q^217q+39)\xi ^2+2(q^417q^3+100q^2244q+214)\xi `$ (4.17)
(4.18)
$`+4(q3)(q^411q^3+44q^276q+46)=0`$ (4.19)
with roots $`\lambda _{tt4,j}`$, $`j=17,18,19`$. The $`\lambda _{tt4,j}`$’s for $`20j31`$ are roots of three quartic equations,
$`\xi ^4+2(q^39q^2+29q34)\xi ^3+2(q^39q^2+29q34)^2\xi ^2`$ (4.20)
(4.21)
$`+4(q3)^2(q^25q+5)(q^39q^2+29q34)\xi +4(q3)^4(q^25q+5)^2=0`$ (4.22)
(4.23)
(4.24)
with roots $`\lambda _{tt4,j}`$, $`20j23`$,
$`\xi ^42(q^27q+14)\xi ^3+2(q^27q+14)^2\xi ^2`$ (4.25)
(4.26)
$`+4(q3)(q^27q+14)(q^26q+7)\xi +4(q3)^2(q^26q+7)^2=0`$ (4.27)
(4.28)
(4.29)
with roots $`\lambda _{tt4,j}`$, $`24j27`$, and
$`\xi ^4+2(3q11)\xi ^3+2(3q11)^2\xi ^2`$ (4.30)
(4.31)
$`+2(3q11)(3q^218q+23)\xi +(3q^218q+23)^2=0`$ (4.32)
(4.33)
(4.34)
with roots $`\lambda _{tt4,j}`$, $`28j31`$. Finally, the $`\lambda _{tt4,j}`$’s for $`32j37`$ are roots of an equation of degree six:
$`\xi ^62(q5)(2q7)\xi ^5+2(q5)^2(2q7)^2\xi ^4`$ (4.35)
(4.36)
$`+8(q4)^2(3q^329q^2+89q85)\xi ^3+4(3q^328q^2+84q79)^2\xi ^2`$ (4.37)
(4.38)
$`+8(q3)^2(q^25q+5)(3q^328q^2+84q79)\xi +8(q3)^4(q^25q+5)^2=0.`$ (4.39)
(4.40)
(4.41)
Each of the three quartic equations above has roots of the form $`a_{\mathrm{}}e^{\pm i\pi /4}`$, $`b_{\mathrm{}}e^{\pm i\pi /4}`$, where $`\mathrm{}=1,2,3`$ indexes the quartic equation, so
$$\lambda _{tt4,j}=a_1e^{\pm i\pi /4}\mathrm{for}j=20,21$$
(4.42)
$$\lambda _{tt4,j}=b_1e^{\pm i\pi /4}\mathrm{for}j=22,23$$
(4.43)
$$\lambda _{tt4,j}=a_2e^{\pm i\pi /4}\mathrm{for}j=24,25$$
(4.44)
$$\lambda _{tt4,j}=b_2e^{\pm i\pi /4}\mathrm{for}j=26,27$$
(4.45)
$$\lambda _{tt4,j}=a_3e^{\pm i\pi /4}\mathrm{for}j=28,29$$
(4.46)
$$\lambda _{tt4,j}=b_3e^{\pm i\pi /4}\mathrm{for}j=30,31$$
(4.47)
where the values of $`a_{\mathrm{}}`$ and $`b_{\mathrm{}}`$, $`\mathrm{}=1,2,3`$ are determined by these quartic equations. Similarly, the roots of the sixth-order equation are of the form $`c_{\mathrm{}}e^{\pm i\pi /4}`$, $`\mathrm{}=1,2,3`$, i.e.,
$$\lambda _{tt4,j}=c_1e^{\pm i\pi /4}\mathrm{for}j=32,33$$
(4.48)
$$\lambda _{tt4,j}=c_2e^{\pm i\pi /4}\mathrm{for}j=34,35$$
(4.49)
$$\lambda _{tt4,j}=c_3e^{\pm i\pi /4}\mathrm{for}j=36,37$$
(4.50)
where the values of $`c_{\mathrm{}}`$, $`\mathrm{}=1,2,3`$ follow from eq. (4.41). Below we shall comment further on these phase factors.
The corresponding coefficients are
$$c_{tt4,1}=\frac{1}{4}q(q2)(q3)^2$$
(4.51)
$$c_{tt4,2}=c_{tt4,3}=\frac{1}{12}(q1)(q2)(3q^211q6)$$
(4.52)
$$c_{tt4,j}=\frac{1}{2}(q1)(q2)\mathrm{for}j=4,7\mathrm{and}32j37$$
(4.53)
$$c_{tt4,5}=\frac{2}{3}q(q2)(q4)$$
(4.54)
$$c_{tt4,j}=\frac{1}{3}q(q2)(q4)\mathrm{for}j=10,11\mathrm{and}28j31$$
(4.55)
$$c_{tt4,6}=\frac{1}{3}(q1)(q^25q+3)$$
(4.56)
$$c_{tt4,8}=c_{tt4,9}=1$$
(4.57)
$$c_{tt4,j}=\frac{1}{2}q(q3)\mathrm{for}j=12,13,17,18,19\mathrm{and}24j27$$
(4.58)
$$c_{tt4,j}=q1\mathrm{for}j=14,15,16\mathrm{and}20j23.$$
(4.59)
Formally, we have also found a zero eigenvalue,
$$\lambda _{tt4,38}=0$$
(4.60)
with coefficient (multiplicity)
$$c_{tt4,38}=\frac{1}{12}q(q1)(3q^217q+40).$$
(4.61)
Although this term does not contribute to the chromatic polynomial (1.4), the corresponding coefficient does contribute to the sum of multiplicities, i.e. to the total dimension of the space of coloring configurations, given by (1.6). The sum of all of the coefficients, including that corresponding to the zero eigenvalue, is equal to $`P(C_4,q)=q(q1)(q^23q+3)`$, which is an $`L_y=4`$ special case of (1.6).
The singular locus $``$ for the $`L_x\mathrm{}`$ limit of the strip of the triangular lattice with $`L_y=4`$ and toroidal boundary conditions is shown in Fig. 2. For comparison, chromatic zeros are calculated and shown for length $`L_x=30`$ (i.e., $`n=120`$ vertices). The locus $``$ crosses the real axis at the points $`q=0`$, $`q=2`$, and at the maximal point $`q=q_c`$, where
$$q_c=4\mathrm{for}\{G\}=(tri,4\times L_x,PBC_y,PBC_x).$$
(4.62)
At this point there are several degeneracies of magnitudes of eigenvalues; these occur for $`\lambda _j`$ with $`j=1,4,6,7,8,9`$ and $`14j19`$.
As is evident from Fig. 2, the locus $``$ separates the $`q`$ plane into different regions including the following (we use the same symbols as for the $`L_y=4`$ toroidal strip of the square lattice, but it is understood that the regions are specific to this section): (i) $`R_1`$, containing the semi-infinite intervals $`q>4`$ and $`q<0`$ on the real axis and extending outward to infinite $`|q|`$, (ii) $`R_2`$ containing the interval $`2<q<4`$, and (iii) $`R_3`$ containing the real interval $`0<q<2`$ Again, the (nonzero) density of chromatic zeros has the smallest values on the curve separating regions $`R_1`$ and $`R_3`$ in the vicinity of the point $`q=0`$ and on the curve separating regions $`R_2`$ and $`R_3`$ in the vicinity of the point $`q=2`$.
In region $`R_1`$, $`\lambda _{tt4,8}`$ is the dominant $`\lambda _{G,j}`$, so
$$W=(\lambda _{tt4,8})^{1/4},qR_1.$$
(4.63)
This is the same as $`W`$ for the corresponding $`L_x\mathrm{}`$ limit of the strip of the triangular lattice with the same width $`L_y=4`$ and cylindrical $`(PBC_y,FBC_x)`$ boundary conditions, calculated in .
In region $`R_2`$,
$$|W|=|\lambda _{tt4,17}|^{1/4},qR_2$$
(4.64)
where $`\lambda _{tt3,17}`$ is the root of the cubic equation (4.19) that has the maximal magnitude for $`2<q<4`$. In region $`R_3`$,
$$|W|=|\lambda _{tt4,14}|^{1/4},qR_3.$$
(4.65)
There are no other regions containing nonzero intervals of the real axis besides $`R_j`$, $`j=1,2,3`$. However, our previous calculations for various families of graphs have shown that $``$ can include pairs of extremely small complex-conjugate sliver regions. We have not made an exhaustive search for these in the present case.
Corresponding to eq. (2.52) for the toroidal or Klein bottle and cylindrical strips of the square lattice, we have
$$W(tri,4\times \mathrm{},PBC_y,(T)PBC_x,q)=W(tri,4\times \mathrm{},PBC_y,FBC_x,q)\mathrm{for}q4.$$
(4.66)
Hence the values of $`W(tri,4\times \mathrm{},PBC_y,FBC_x,q)`$ for various values of $`q4`$ given in (see also ) are also applicable here. For low integral values of $`q`$ we list the values of $`|W(q)|`$ for this strip in Table 3, together with corresponding values given in for $`W`$ in the $`L_x\mathrm{}`$ limit of the $`L_y=3`$ strip with $`(PBC_y,(T)PBC_x)`$.
The locus $``$ only has support for $`Re(q)0`$. This can be demonstrated by carrying out a Taylor series expansion of the degeneracy equation $`|\lambda _{tt4,8}|=|\lambda _{tt4,14}|`$ near the origin, which is, numerically,
$$|140.15141.25q+57.6q^2+O(q^3)|=|140.1593.4q+21.9q^2+O(q^3)|.$$
(4.67)
This equation is of the form (2.54), and, using eq. (2.55), we verify that $``$ bends to the right as one moves away from the origin. Farther away from the origin, one can see from Fig. 2 that $``$ continues to move into the half-plane with $`Re(q)>0`$, so that the conclusion stated above follows, that this locus has no support for $`Re(q)<0`$.
## 5 $`L_y=4`$ Strip of the Triangular Lattice with $`(PBC_y,TPBC_x)`$
The strip of the triangular lattice with width $`L_y=4`$, arbitrarily great length $`L_x`$, and $`(PBC_y,TPBC_x)=`$ Klein bottle boundary conditions, labelled $`tk4`$, has (for $`L_x2`$) chromatic number
$$\chi (tri,4\times L_x,PBC_y,TPBC_x)=\{\begin{array}{cc}4\hfill & \text{if }L_x\text{ is even}\hfill \\ 5\hfill & \text{if }L_x\text{ is odd}\hfill \end{array}$$
(5.1)
In the notation of eq. (1.4) we find $`N_{tk4,\lambda }=12`$ and
$$P(tri,4\times L_x,PBC_y,TPBC_x,q)=\underset{j=1}{\overset{12}{}}c_{tk4,j}(\lambda _{tk4,j})^{L_x}$$
(5.2)
where
$$\lambda _{tk4,1}=\lambda _{tt4,1}=2$$
(5.3)
$$\lambda _{tk4,2}=\lambda _{tt4,4}=2(3q)$$
(5.4)
$$\lambda _{tk4,3}=\lambda _{tt4,6}=2(2q9)$$
(5.5)
$$\lambda _{tk4,4}=\lambda _{tt4,7}=2(q3)^2$$
(5.6)
$$\lambda _{tk4,j}=\lambda _{tt4,j+3}\mathrm{for}j=5,6$$
(5.7)
$$\lambda _{tk4,j}=\lambda _{tt4,j+7}\mathrm{for}7j12.$$
(5.8)
The corresponding coefficients are
$$c_{tk4,1}=\frac{1}{2}q(q2)(q3)$$
(5.9)
$$c_{tk4,2}=c_{tt4,4}=\frac{1}{2}(q1)(q2)$$
(5.10)
$$c_{tk4,3}=(q1)$$
(5.11)
$$c_{tk4,4}=c_{tt4,7}=\frac{1}{2}(q1)(q2)$$
(5.12)
$$c_{tk4,5}=c_{tk4,6}=c_{tt4,8}=c_{tt4,9}=1$$
(5.13)
$$c_{tk4,j}=c_{tt4,j+7}=q1\mathrm{for}j=7,8,9$$
(5.14)
$$c_{tk4,j}=c_{tt4,j+7}=\frac{1}{2}q(q3)\mathrm{for}j=10,11,12.$$
(5.15)
The coloring matrix also has another eigenvalue, namely,
$$\lambda _{tk4,13}=0$$
(5.16)
with multiplicity
$$c_{tk4,13}=\frac{1}{2}q(q1)^2.$$
(5.17)
Hence, the total number of distinct eigenvalues of the coloring matrix for this strip is $`N_{tk4,\lambda ,tot}=N_{tk4,\lambda }+1=13`$. The sum of all of the coefficients, including that for the zero eigenvalue, is zero; this is an $`L_y=4`$ special case of (1.7).
## 6 Cyclic and Toroidal Crossing-Subgraph Strips of the Square Lattice
### 6.1 General
It is worthwhile to include here some results on certain related families of strip graphs since these give insight into the structure of the chromatic polynomials for the various strips with longitudinal boundary conditions which are periodic or periodic with reversed orientation. Let us consider first a strip of the square lattice of fixed width $`L_y`$ and arbitrarily great length $`L_x`$ constructed as follows. As before, the longitudinal (horizontal) direction on the strip to be $`x`$ and the transverse (vertical) direction to be $`y`$. Label the vertices of two successive transverse slices of the strip, starting at the top as $`(1,2,..,L_y)`$ and $`(1^{},2^{},\mathrm{},L_y^{})`$. First, consider the case of free transverse boundary conditions, for which these transverse slices of the strip are line (path) graphs with $`L_y`$ vertices. Connect these with edges linking vertices 1 to $`L_y^{}`$, 2 to $`(L_y1)^{}`$, 3 to $`(L_y2)^{}`$, and so forth. For example, for $`L_y=2`$, we connect 1 to $`2^{}`$ and 2 to $`1^{}`$; for $`L_y=3`$, we connect 1 to $`3^{}`$, 2 to $`2^{}`$, and 3 to $`1^{}`$, etc. for other values of $`L_y`$. An example of this crossing-subgraph strip of the square lattice of width $`L_y=3`$ is given in Fig. 3(a).
We impose periodic longitudinal boundary conditions. We shall denote this crossing-subgraph strip (labelled $`cg`$) of the square ($`sq`$) lattice as $`cg(sq,L_y\times L_x,FBC_y,PBC_x)`$. We observe that
$$cg(sq,L_y\times L_x,FBC_y,PBC_x)=\{\begin{array}{cc}(sq,L_y\times L_x,FBC_y,PBC_x)\hfill & \text{if }L_x\text{ is even}\hfill \\ (sq,L_y\times L_x,FBC_y,TPBC_x)\hfill & \text{if }L_x\text{ is odd}\hfill \end{array}$$
(6.1.1)
That is, for even (odd) $`L_x`$, this crossing-subgraph strip reduces to the cyclic (Möbius) strip of the square lattice. Secondly, consider the case where we impose periodic transverse boundary conditions; we denote this toroidal crossing-subgraph strip of the square lattice as $`cg(sq,L_y\times L_x,PBC_y,PBC_x)`$. In this case the transverse slices are circuit graphs with $`L_y`$ vertices. An example of this toroidal crossing-subgraph strip of the square lattice with width $`L_y=3`$ is shown in Fig. 3(b). We have
$$cg(sq,L_y\times L_x,PBC_y,PBC_x)=\{\begin{array}{cc}(sq,L_y\times L_x,PBC_y,PBC_x)\hfill & \text{if }L_x\text{ is even}\hfill \\ (sq,L_y\times L_x,PBC_y,TPBC_x)\hfill & \text{if }L_x\text{ is odd}\hfill \end{array}$$
(6.1.2)
That is, for even (odd) $`L_x`$, this crossing-subgraph strip reduces to the strip of the square lattice with torus (Klein bottle) boundary conditions. Given the relations (6.1.1) and (6.1.2), it follows that a knowledge of the chromatic polynomial for the cyclic crossing-subgraph strip of width $`L_y`$ of the square lattice is equivalent to a knowledge of the chromatic polynomials for the strip of this lattice with width $`L_y`$ with both cyclic and Möbius boundary conditions and, similarly, a knowledge of the chromatic polynomial for the toroidal crossing-subgraph strip of the square lattice is equivalent to a knowledge of the chromatic polynomials for the strip of this lattice with both torus and Klein bottle boundary conditions.
The sum of the coefficients for the $`L_y\times L_x`$ cyclic crossing-subgraph strip of the square lattice is
$$C_{cg(sq,L_y\times L_x,FBC_y,PBC_x)}=P(T_{L_y},q)$$
(6.1.3)
where $`T_n`$ is the tree graph on $`n`$ vertices, and $`P(T_n,q)=q(q1)^{n1}`$. The sum of the coefficients for the $`L_y\times L_x`$ toroidal crossing-subgraph strip of the square or triangular lattice is
$$C_{cg(G_s,L_y\times L_x,PBC_y,PBC_x)}=P(C_{L_y},q)\mathrm{for}G_s=sq,tri$$
(6.1.4)
as in (1.6).
We have carried out explicit calculations of chromatic polynomials for a number of crossing-subgraph strips (labelled $`cg`$) and have related the results to those for strips with cyclic/Möbius and torus/Klein bottle boundary conditions. We concentrate here on strips of the square lattice and discuss those of the triangular lattice below. We find that a certain subset of the terms in (1.4) for the crossing-subgraph strips occur in opposite-sign pairs of the form $`\pm \lambda _{cg,j}`$. Consider the coefficients for the $`\pm \lambda _{cg,j}`$’s in each pair: in some cases, these are different, while in others they are the same. Let us denote the number of $`\lambda _{cg,j}`$’s comprising opposite-sign pairs such that the members of each pair have different (the same) coefficients as $`N_{cg,opd,\lambda }`$ ($`N_{cg,ops,\lambda }`$), respectively. The number of remaining $`\lambda _{cg,j}`$’s that are not members of an opposite-sign pair is denoted $`N_{cg,up,\lambda }`$, where $`up`$ means “unpaired”. Clearly
$$N_{cg,\lambda }=N_{cg,up,\lambda }+N_{cg,opd,\lambda }+N_{cg,ops,\lambda }.$$
(6.1.5)
For the cyclic and Möbius strips that we have studied, we find $`N_{cg,ops,\lambda }=0`$, while for torus and Klein bottle strips, $`N_{cg,ops,\lambda }`$ is, in general, nonzero. For even $`L_x`$, where, according to the identities (6.1.1) and (6.1.2), the cyclic (toroidal) crossing-subgraph strip reduces to the cyclic (toroidal) strip of the square lattice, two such opposite-sign terms reduce to a single term as follows (the subscripts $`jp,jm`$ denote $`j,\pm `$)
$`c_{cg,cyc,jp}(\lambda _{cg,cyc,j})^{L_x}+c_{cg,cyc,jm}(\lambda _{cg,cyc,j})^{L_x}`$ $`=`$ $`(c_{cg,cyc,jp}+c_{cg,cyc,jm})(\lambda _{cg,cyc,j})^{L_x}`$ (6.1.6)
$`=`$ $`c_{sq,cyc,j}(\lambda _{cg,cyc,j})^{L_x}.`$ (6.1.8)
For odd $`L_x=m`$ the cyclic (toroidal) crossing-subgraph strip reduces to the Möbius (Klein bottle) strip, and the pair of opposite-sign terms reduces to a single term as follows:
$`c_{cg,cyc,jp}(\lambda _{cg,cyc,j})^{L_x}+c_{cg,cyc,jm}(\lambda _{cg,cyc,j})^{L_x}`$ $`=`$ $`(c_{cg,cyc,jp}c_{cg,cyc,jm})(\lambda _{cg,cyc,j})^{L_x}`$ (6.1.9)
$`=`$ $`c_{sq,Mb,j}(\lambda _{cg,cyc,j})^{L_x}`$ (6.1.11)
where the subscript $`Mb`$ denotes Möbius. In particular, if $`\lambda _{cg,j}`$ is one of the $`N_{cg,ops,\lambda }`$ terms with $`c_{cg,jp}=c_{cg,jm}`$, then the terms in (6.1.11) cancel each other, leaving no contribution. As noted above, in our studies, we have found that this can happen for Klein bottle strips, since $`N_{cg,ops,\lambda }0`$ for these, but not for Möbius strips, since $`N_{cg,ops,\lambda }=0`$ for these. The inverse relations connecting the coefficients for the terms $`\pm \lambda _{cg,cyc,j}`$ in the chromatic polynomial for the crossing-subgraph cyclic strip to the coefficients $`c_{sq,cyc,j}`$ and $`c_{sq,Mb,j}`$ in the cyclic and Möbius strips are thus
$$c_{cg,cyc,jp}=\frac{1}{2}(c_{sq,cyc,j}+c_{sq,Mb,j})$$
(6.1.12)
$$c_{cg,cyc,jm}=\frac{1}{2}(c_{sq,cyc,j}c_{sq,Mb,j})$$
(6.1.13)
and similarly, for the coefficients for the terms $`\pm \lambda _{cg,torus,j}`$ in the chromatic polynomial for the toroidal crossing-subgraph strip in terms of the coefficients $`c_{sq,torus,j}`$ and $`c_{sq,Kb,j}`$ in the torus and Klein bottle ($`Kb`$) strips,
$$c_{cg,torus,jp}=\frac{1}{2}(c_{sq,torus,j}+c_{sq,Kb,j})$$
(6.1.14)
$$c_{cg,torus,jm}=\frac{1}{2}(c_{sq,torus,j}c_{sq,Kb,j}).$$
(6.1.15)
From these considerations, we derive the following general formula:
$$N_{sq,L_y,cyc,\lambda }=N_{cg,sq,L_y,cyc,\lambda }\frac{1}{2}N_{cg,sq,L_y,cyc,opd,\lambda }$$
(6.1.16)
and, since $`N_{cg,ops,\lambda }=0`$ for the cyclic crossing-graph strips of the square lattice that we have studied, the same formula applies to the corresponding Möbius strips with $`N_{sq,L_y,cyc,\lambda }`$ replaced by $`N_{sq,L_y,Mb,\lambda }`$. Further,
$$N_{sq,L_y,torus,\lambda }=N_{cg,sq,L_y,torus,\lambda }\frac{1}{2}N_{cg,sq,L_y,torus,opd,\lambda }\frac{1}{2}N_{cg,sq,L_y,torus,ops,\lambda }$$
(6.1.17)
$$N_{sq,L_y,Kb,\lambda }=N_{cg,sq,L_y,torus,\lambda }\frac{1}{2}N_{cg,sq,L_y,torus,opd,\lambda }N_{cg,sq,L_y,torus,ops,\lambda }.$$
(6.1.18)
Thus,
$$N_{sq,L_y,torus,\lambda }N_{sq,L_y,Kb,\lambda }=\frac{1}{2}N_{cg,sq,L_y,torus,ops,\lambda }.$$
(6.1.19)
In the limit $`L_x\mathrm{}`$, the cyclic or toroidal crossing-subgraph strip of a given width $`L_y`$ yields a $`W`$ function via (1.2) and hence a singular locus $``$. An important theorem can be proved from this by observing that we can take this limit using even or odd values of $`L_x`$; in the even-$`L_x`$ case, we obtain the function $`W`$ and locus $``$ for the strip with torus boundary conditions, while in the odd-$`L_x`$ case, we obtain the $`W`$ function and $``$ for the strip with Klein bottle boundary conditions. Since the original limit exists, all three of these limits must be the same. This proves the following theorem:
Theorem 1: The $`W`$ function and singular locus $``$ are the same for the $`L_x\mathrm{}`$ limit of strip of the square lattice with width $`L_y`$ and length $`L_x`$ whether one imposes $`(PBC_y,PBC_x)`$ or $`(PBC_y,TPBC_x)`$, i.e. torus or Klein bottle boundary conditions.
Since the chromatic polynomials for these two sets of boundary conditions involve different numbers of terms, this was not, a priori obvious. This feature was first noticed in and was shown there to be a consequence of the fact that none of the terms $`\lambda _{st3,j}`$ for the torus case that were absent in the Klein bottle case was dominant; here we have succeeded in explaining why this had to be true; if it were not, then the respective loci $``$ would be different, but this is impossible, as a consequence of our present theorem. Thus, a corollary to the theorem is
Corollary 1: Consider an $`L_y\times L_x`$ strip of the square lattice with torus or Klein bottle boundary conditions, and denote the set of nonzero eigenvalues that contribute to (1.4) for these two respective strips as $`\lambda _{stL_y,j}`$, $`j=1,..,N_{stL_y,\lambda }`$ and $`\lambda _{skL_y,j}`$, $`j=1,..,N_{skL_y,\lambda }`$. Focus on the set of eigenvalues $`\lambda _{stL_y,j}`$ that do not occur among the set $`\lambda _{skL_y,j}`$ (the number of these is given by eq. (6.1.19)); none of these can be dominant eigenvalues.
In a similar manner, one can prove that the locus $``$ for the $`L_x\mathrm{}`$ limits of the strips of the square lattice with cyclic and Möbius boundary conditions are the same without using as input the identity of terms $`\lambda _{sq,L_y,cyc,j}=\lambda _{sq,L_y,Mb,j}`$ that we have observed in our studies . Below we shall also prove a similar theorem for strips of the triangular lattice with torus and Klein bottle boundary conditions.
We recall that in we observed from our work that the coefficients $`c_{G,j}`$ in (1.4) for cyclic and Möbius strips of the square lattice are Chebyshev polynomials; in particular, for a given degree $`d`$ polynomial, there is a unique coefficient with this degree, and it is given by
$$c^{(d)}=U_{2d}\left(\frac{\sqrt{q}}{2}\right)$$
(6.1.20)
where $`U_n(x)`$ is the Chebyshev polynomial of the second kind, defined by
$$U_n(x)=\underset{j=0}{\overset{[\frac{n}{2}]}{}}(1)^j\left(\genfrac{}{}{0pt}{}{nj}{j}\right)(2x)^{n2j}$$
(6.1.21)
where $`[\frac{n}{2}]`$ means the integral part of $`\frac{n}{2}`$. The first few of these coefficients are $`c^{(0)}=1`$, $`c^{(1)}=q1`$, $`c^{(2)}=q^23q+1`$, $`c^{(3)}=q^35q^2+6q1`$, etc. We also found that the eigenvalues $`\lambda _{G,j}`$ were the same for cyclic and Möbius strips of the square (and triangular) lattices of a given width that we considered. We established the transformation rules specifying how a coefficient of a given degree changes when one switches from the cyclic to Möbius strip of the square lattice :
$$c^{(0)}\pm c^{(0)}$$
(6.1.22)
$$c^{(2k)}\pm c^{(k1)},1k\left[\frac{L_y}{2}\right]$$
(6.1.23)
$$c^{(2k+1)}\pm c^{(k+1)},0k\left[\frac{L_y1}{2}\right].$$
(6.1.24)
Following the notation of , denote the number of terms (eigenvalues) $`\lambda _{G,j}`$ with coefficients $`c^{(d)}`$ as $`n_P(L_y,d)`$. We concentrate on the case of the square strip here and suppress the $`sq`$ in the notation. This satisfies
$$N_{L_y,\lambda }=\underset{d=0}{\overset{L_y}{}}n_P(L_y,d)$$
(6.1.25)
where in the notation used here, $`N_{L_y,\lambda }`$ refers to the quantity denoted $`N_{P,L_y,\lambda }`$ in . We gave general formulas for $`n_P(L_y,d)`$ and $`N_{L_y,\lambda }`$; in particular, here we shall need the following ones:
$$n_P(L_y,0)=M_{L_y1}$$
(6.1.26)
and
$$n_P(L_y,1)=M_{L_y}$$
(6.1.27)
where the Motzkin number $`M_n`$ is given by
$$M_n=\underset{j=0}{\overset{n}{}}(1)^jC_{n+1j}\left(\genfrac{}{}{0pt}{}{n}{j}\right)$$
(6.1.28)
and
$$C_n=\frac{1}{n+1}\left(\genfrac{}{}{0pt}{}{2n}{n}\right)$$
(6.1.29)
is the Catalan number. For the total number of terms, we obtained the result
$$N_{L_y,\lambda }=2(L_y1)!\underset{j=0}{\overset{[\frac{L_y}{2}]}{}}\frac{(L_yj)}{(j!)^2(L_y2j)!}.$$
(6.1.30)
Now our transformation formulas (6.1.22)-(6.1.24) imply that the only cases where the coefficients in (1.4) for the cyclic and Möbius strips can be the same, up to sign, are for $`c^{(0)}=1`$ and $`c^{(1)}`$. If these coefficients are the same, then, by eq. (6.1.13), $`c_{cg,jm}=0`$, while if they are opposite in sign, then, by eq. (6.1.12), $`c_{cg,jp}=0`$; hence, in either case, there is only one term from the possible pair $`\pm \lambda _{cg,sq,cyc,j}`$ contributing to (1.4) for the cyclic crossed-subgraph strip. For all of the other coefficients $`c_{sq,cyc,j}`$ in (1.4) for the cyclic strips, our transformation formulas (6.1.22)-(6.1.24) imply that $`c_{sq,Mb,j}\pm c_{sq,cyc,j}`$, so that both $`c_{cg,jp}`$ and $`c_{cg,jm}`$ are nonzero and both of the corresponding pair $`\pm \lambda _{cg,sq,cyc,j}`$ contribute to (1.4) for the cyclic crossed-subgraph strip. It follows that the total number of terms for the cyclic crossed-subgraph strip of the square lattice is given by
$$N_{cg,sq,L_y,\lambda }=2N_{sq,cyc,L_y,\lambda }n_P(L_y,0)n_P(L_y,1).$$
(6.1.31)
Substituting our results from for each of the quantities on the right-hand side, we obtain, for the number of unpaired terms for the crossed-subgraph strip, the relation
$$N_{cg,sq,L_y,up,\lambda }=M_{L_y1}+M_{L_y},$$
(6.1.32)
for the number of terms comprising members of opposite-sign pairs, the relation
$`N_{cg,sq,L_y,opd,\lambda }`$ $`=`$ $`2(N_{cg,sq,L_y,\lambda }N_{sq,cyc,L_y,\lambda })`$ (6.1.33)
$`=`$ $`2\left[2(L_y1)!{\displaystyle \underset{j=0}{\overset{[\frac{L_y}{2}]}{}}}{\displaystyle \frac{(L_yj)}{(j!)^2(L_y2j)!}}M_{L_y1}M_{L_y}\right]`$ (6.1.35)
and for the total number of (nonzero) terms,
$$N_{cg,sq,L_y,\lambda }=4(L_y1)!\underset{j=0}{\overset{[\frac{L_y}{2}]}{}}\frac{(L_yj)}{(j!)^2(L_y2j)!}M_{L_y1}M_{L_y}.$$
(6.1.36)
### 6.2 $`L_y=2`$ Cyclic Crossing-Subgraph Strip of the Square Lattice
For example, for $`L_y=2`$, using the Motzkin numbers $`M_1=1`$, $`M_2=2`$, the total number of terms is $`N_{cg,sq,2,\lambda }=5`$, with $`N_{cg,sq,2,up,\lambda }=3`$ and $`N_{cg,sq,2,opd,\lambda }=2`$. We have explicitly calculated the chromatic polynomial for this case using iterated deletion-contraction and coloring matrix methods and find (with the shorthand $`cgs2`$ for $`cg(sq,2\times L_x,cyc)`$)
$$P(cg(sq,2\times L_x,cyc))=\underset{j=1}{\overset{5}{}}c_{cgs2,j}(\lambda _{cgs2,j})^{L_x}$$
(6.2.1)
where
$$\lambda _{cgs2,j}=\pm 1,j=1,2$$
(6.2.2)
$$\lambda _{cgs2,3}=3q$$
(6.2.3)
$$\lambda _{cgs2,4}=q1$$
(6.2.4)
$$\lambda _{cgs2,5}=q^23q+3$$
(6.2.5)
with coefficients
$$c_{cgs2,1}=\frac{1}{2}(c^{(2)}c^{(0)})=\frac{1}{2}q(q3)$$
(6.2.6)
$$c_{cgs2,2}=\frac{1}{2}(c^{(2)}+c^{(0)})=\frac{1}{2}(q1)(q2)$$
(6.2.7)
$$c_{cgs2,j}=c^{(1)}=q1\mathrm{for}j=3,4$$
(6.2.8)
$$c_{cgs2,5}=c^{(0)}=1.$$
(6.2.9)
For even $`L_x`$, this chromatic polynomial for the width $`L_y=2`$ crossing graph strip of the square lattice reduces to the result for the regular cyclic strip of the square lattice with $`N_{sq,2,cyc,\lambda }=4`$, namely, with $`L_x=m`$,
$$P(sq,2\times L_x,cyc.)=(q^23q+1)+(q1)[(3q)^m+(1q)^m]+(q^23q+3)^m$$
(6.2.10)
while for odd $`L_x`$, the chromatic polynomial (6.2.1) reduces to the result for the $`L_y=2`$ Möbius strip with $`N_{sq,2,Mb,\lambda }=4`$, namely
$$P(sq,2\times L_x,Mb)=1+(q1)\left[(3q)^m(1q)^m\right]+(q^23q+3)^m.$$
(6.2.11)
### 6.3 $`L_y=3`$ Cyclic Crossing-Subgraph Strip of the Square Lattice
The chromatic polynomial for the $`L_y=3`$ cyclic crossing-subgraph strip of the square lattice (labelled $`cgs3`$) can be calculated directly or from the known results for the chromatic polynomials of the $`L_y=3`$ cyclic and Möbius strips of the square lattice. It is worthwhile to display the results here because of the unified understanding that they give concerning the structures of the chromatic polynomials for the cyclic and Möbius strips. From our general formulas above we have $`N_{cg,sq,3,\lambda }=14`$ with $`N_{cg,sq,3,up,\lambda }=6`$ and $`N_{cg,sq,3,opd,\lambda }=8`$. For the respective even and odd values of $`L_x`$ where the chromatic polynomial reduces to that for the $`L_y=3`$ cyclic and Möbius strips of the square lattice, we have $`N_{sq,3,cyc,\lambda }=N_{sq,3,Mb,\lambda }=14(1/2)8=10`$, in agreement with the previous calculations in . We find
$$P(cg(sq,3\times L_x,cyc))=\underset{j=1}{\overset{14}{}}c_{cgs3,j}(\lambda _{cgs3,j})^{L_x}$$
(6.3.1)
where
$$\lambda _{cgs3,j}=\pm 1,j=1,2$$
(6.3.2)
$$\lambda _{cgs3,j}=\pm (q1),j=3,4$$
(6.3.3)
$$\lambda _{cgs3,j}=\pm (q2),j=5,6$$
(6.3.4)
$$\lambda _{cgs3,j}=\pm (q4),j=7,8$$
(6.3.5)
$$\lambda _{cgs3,9}=(q2)^2$$
(6.3.6)
$$\lambda _{cgs3,10}=\lambda _{sq3,6}$$
(6.3.7)
$$\lambda _{cgs3,11}=\lambda _{sq3,7}$$
(6.3.8)
where $`\lambda _{sq3,j}`$ for $`j=6,7`$ were given in eq. (3.10) of , and
$$\lambda _{cgs3,j}=\lambda _{sq3,j4},12j14$$
(6.3.9)
where $`\lambda _{sq3,j}`$ for $`j=8,9,10`$ were defined by eq. (3.11) of .
The corresponding coefficients are
$$c_{cgs3,1}=\frac{1}{2}(c^{(3)}c^{(2)})=\frac{1}{2}(q2)(q^24q+1)$$
(6.3.10)
$$c_{cgs3,2}=\frac{1}{2}(c^{(3)}+c^{(2)})=\frac{1}{2}q(q1)(q3)$$
(6.3.11)
$$c_{cgs3,j}=\frac{1}{2}(c^{(2)}c^{(0)})=\frac{1}{2}q(q3)\mathrm{for}j=3,6,7$$
(6.3.12)
$$c_{cgs3,j}=\frac{1}{2}(c^{(2)}+c^{(0)})=\frac{1}{2}(q1)(q2)\mathrm{for}j=4,5,8$$
(6.3.13)
$$c_{cgs3,j}=c^{(1)}=q1\mathrm{for}j=9,12,13,14$$
(6.3.14)
$$c_{cgs3,j}=1,j=10,11.$$
(6.3.15)
### 6.4 $`L_y=4`$ Cyclic Crossing-Subgraph Strip of the Square Lattice
From our calculations of the chromatic polynomials for the width $`L_y=4`$ cyclic and Möbius strips of the square lattice , we have obtained the chromatic polynomial for the $`L_y=4`$ cyclic crossing-subgraph strip (labelled $`cgs4`$). In accord with our general formulas, we get $`N_{cg,sq,4,\lambda }=39`$ with $`N_{cg,sq,4,up,\lambda }=13`$ and $`N_{cg,sq,4,opd,\lambda }=26`$. For respective even and odd $`L_x`$, the reduction to the $`L_y=4`$ cyclic and Möbius strips has $`N_{sq,4,cyc,\lambda }=N_{sq,4,Mb,\lambda }=39(1/2)26=26`$. We omit listing the terms and their coefficients here since they are lengthy and our previous examples are sufficient to illustrate our general formulas.
In passing, we note that the cyclic crossing-subgraph strip of the triangular lattice does not yield either the cyclic or Möbius strip of this lattice for even or odd $`L_x`$. This provides a further understanding of our earlier findings that the coefficients for the $`L_y=2`$ and $`L_y=3`$ Möbius strips of the triangular lattice are not polynomials in $`q`$.
### 6.5 $`L_y=3`$ Toroidal Crossing-Subgraph Strip of the Square Lattice
For the toroidal crossing-subgraph strip of the square lattice with width $`L_y=3`$ (labelled $`cgst3`$) we find $`N_{cgst3,\lambda }=12`$ with $`N_{cgst3,up,\lambda }=4`$, $`N_{cgst3,opd,\lambda }=2`$, and $`N_{cgst3,ops,\lambda }=6`$ (see Table 4). Our result is
$$P(cg(sq,3\times L_x,PBC_y,PBC_x))=\underset{j=1}{\overset{12}{}}c_{cgst3,j}(\lambda _{cgst3,j})^{L_x}$$
(6.5.1)
with
$$\lambda _{cgst3,j}=\pm 1,j=1,2$$
(6.5.2)
$$\lambda _{cgst3,3}=1q$$
(6.5.3)
$$\lambda _{cgst3,j}=\pm (q2)\mathrm{for}j=4,5$$
(6.5.4)
$$\lambda _{cgst3,j}=\pm (q4)\mathrm{for}j=6,7$$
(6.5.5)
$$\lambda _{cgst3,8}=q5$$
(6.5.6)
$$\lambda _{cgst3,j}=\pm (q2)^2\mathrm{for}j=9,10$$
(6.5.7)
$$\lambda _{cgst3,11}=(q^27q+13)$$
(6.5.8)
$$\lambda _{cgst3,12}=q^36q^2+14q13$$
(6.5.9)
with corresponding coefficients
$$c_{cgst3,1}=\frac{1}{2}(q2)(q^24q+1)$$
(6.5.10)
$$c_{cgst3,2}=\frac{1}{2}q(q^26q+7)$$
(6.5.11)
$$c_{cgst3,j}=\frac{1}{2}(q1)(q2)\mathrm{for}j=3,6,7$$
(6.5.12)
$$c_{cgst3,j}=\frac{1}{2}q(q3)\mathrm{for}j=4,5,8$$
(6.5.13)
$$c_{cgst3,j}=q1\mathrm{for}9j11$$
(6.5.14)
$$c_{cgst3,12}=1.$$
(6.5.15)
Hence, in the even-$`L_x`$ case, the reduction to the chromatic polynomial for the $`L_y=3`$ toroidal strip of the square lattice (labelled $`st3`$) has, by eq. (6.1.17), $`N_{st3,\lambda }=1213=8`$ terms, while in the odd-$`L_x`$ case, the reduction to the chromatic polynomial for the $`L_y=3`$ Klein bottle strip of the square lattice (labelled $`sk3`$) has $`N_{sk3,\lambda }=1216=5`$. These are in agreement with the results that were obtained in . The resultant coefficients for the torus and Klein bottle strips can be computed from the analogues of formulas (6.1.8) and (6.1.11) and agree with those given in .
These correspondences shed further light on the coefficients that enter into the chromatic polynomials for strip graphs with torus and Klein bottle boundary conditions. We recall that these are not of the simple form expressible in terms of Chebyshev polynomials of the second kind that we found for strip graphs with cyclic and Möbius boundary conditions in . Among other differences, there is not a unique coefficient with a given degree $`d`$ in $`q`$; for example, for degree $`d=2`$, both $`(q1)(q2)/2`$ and $`q(q3)/2`$ as coefficients. However, the reduction of the chromatic polynomials for the toroidal crossing-subgraph strips to the respective chromatic polynomials for the torus and Klein bottle strips yields relations linking the coefficients in the latter to Chebyshev polynomials, such as (6.2.6) and (6.2.7). In the Appendix we list our results for the $`L_y=4`$ crossing-subgraph toroidal strip of the square lattice (labelled $`cgst4`$)
## 7 Toroidal Crossing-Subgraph Strips of the Triangular Lattice
We consider here a strip of the triangular lattice of fixed width $`L_y`$ and arbitrarily great length $`L_x`$ constructed as follows. As before, label the vertices of two successive transverse slices of the strip, starting at the top, as $`(1,2,..,L_y)`$ and $`(1^{},2^{},\mathrm{},L_y^{})`$. With periodic transverse boundary conditions, these transverse slices form circuit graphs, $`C_{L_y}`$. Connect these with edges linking vertices 1 to $`L_y^{}`$, 2 to $`(L_y1)^{}`$, and so forth. This forms the toroidal crossing-subgraph strip of the square lattice. Next, add diagonal edges as illustrated for the $`L_y=3`$ case in Fig. 4. Finally, impose periodic longitudinal boundary conditions. This yields the crossing-subgraph (cg) strip of the triangular (tri) lattice with toroidal boundary conditions. We shall label this strip as $`cg(tri,L_y\times L_x,PBC_y,PBC_x)`$.
We first observe that
$$cg(tri,L_y\times L_x,PBC_y,PBC_x)=tri(L_y\times L_x,PBC_y,PBC_x)\mathrm{if}L_x=0\mathrm{mod}2L_y$$
(7.1)
i.e., for $`L_x=0`$ mod $`2L_y`$, the toroidal crossing-subgraph strip of the triangular lattice reduces to the toroidal strip of this lattice. Associated with this result, there are terms in the chromatic polynomial for the resultant toroidal strip of the triangular lattice corresponding to opposite-sign pairs of $`\lambda _{cgtL_y,j}`$’s in the chromatic polynomial for the toroidal crossing-subgraph strip that appear with phase factors of the form $`e^{\pm \pi i/L_y}`$ and certain products thereof. For $`L_y=3`$ these are $`e^{\pm 2\pi i/3}`$, while for $`L_y=4`$ they are $`e^{\pm \pi \mathrm{}i/4}`$, $`\mathrm{}=1,2`$. It is thus necessary to distinguish between terms $`\lambda _{cgtL_y,j}`$ in the chromatic polynomial for the $`L_y`$ toroidal crossing-subgraph strip of the triangular lattice that correspond to terms in the toroidal strip of the triangular lattice that do, or do not, have the form of complex-conjugate pairs of terms $`\lambda _{ttL_y,j}`$ with complex prefactors such as $`e^{\pm 2\pi i/3}`$. (where the subscript $`ttL_y`$ refers to the toroidal strip of the triangular lattice with width $`L_y`$). We thus define $`N_{cg,tri,L_y,torus,ops,r,\lambda }`$ and $`N_{cg,tri,L_y,torus,ops,i,\lambda }`$ respectively as the number of $`\lambda _{cgtL_y,j}`$’s that comprise opposite-sign pairs such that the members of each pair have the same coefficient and correspond to a $`\lambda _{ttL_y,j}`$ that is real (is a member of a pair with complex prefactor) for real $`q`$. The resultant general formula relating these numbers of terms is
$$N_{tri,L_y,torus,\lambda }=N_{cg,tri,L_y,torus,\lambda }\frac{1}{2}N_{cg,tri,L_y,torus,opd,\lambda }\frac{1}{2}N_{cg,tri,L_y,torus,ops,r,\lambda }$$
(7.2)
If (i) $`L_y`$ is odd and $`L_x`$ is odd, or (ii) if $`L_y`$ is even and $`L_x=1`$ mod 4, then the toroidal crossing-subgraph strip of the triangular lattice reduces to the strip of this lattice with Klein bottle boundary conditions $`(PBC_y,TPBC_x)`$. In this case, the reduction of the number of terms is determined by eq. (6.1.18) with the obvious replacement of square by triangular lattice strip, i.e.,
$$N_{tri,L_y,Kb,\lambda }=N_{cg,tri,L_y,torus,\lambda }\frac{1}{2}N_{cg,tri,L_y,torus,opd,\lambda }N_{cg,tri,L_y,torus,ops,\lambda }.$$
(7.3)
Thus,
$$N_{tri,L_y,torus,\lambda }N_{tri,L_y,Kb,\lambda }=N_{cg,tri,L_y,torus,ops,\lambda }\frac{1}{2}N_{cg,tri,L_y,torus,ops,r,\lambda }.$$
(7.4)
In the limit $`L_x\mathrm{}`$, the width $`L_y`$ toroidal crossing-subgraph strip of the triangular lattice yields a $`W`$ function via (1.2) and hence a singular locus $``$. As before, we can take this limit using values of $`L_x`$ such that the crossing-subgraph strip reduces to the strip of the triangular lattice with either torus or Klein bottle boundary conditions. Given that the original limit exists, all three of these limits must be the same. This proves
Theorem 2: The $`W`$ function and singular locus $``$ are the same for the $`L_x\mathrm{}`$ limit of strip of the triangular lattice with width $`L_y`$ and length $`L_x`$ whether one imposes $`(PBC_y,PBC_x)`$ or $`(PBC_y,TPBC_x)`$, i.e. torus or Klein bottle boundary conditions.
Hence also
Corollary 2: Consider an $`L_y\times L_x`$ strip of the triangular lattice with torus or Klein bottle boundary conditions, and denote the set of nonzero eigenvalues that contribute to (1.4) for these two respective strips as $`\lambda _{ttL_y,j}`$, $`j=1,..,N_{ttL_y,\lambda }`$ and $`\lambda _{tkL_y,j}`$, $`j=1,..,N_{tkL_y,\lambda }`$. Focus on the set of eigenvalues $`\lambda _{ttL_y,j}`$ that do not occur among the set $`\lambda _{tkL_y,j}`$ (the number of these is given by (7.4)); none of these can be dominant eigenvalues.
For the $`L_y=3`$ crossing-subgraph toroidal strip of the triangular lattice (labelled $`cgt3`$) we find that $`N_{cgt3,\lambda }=12`$ and calculate
$$P(cg(tri,3\times L_x,torus))=\underset{j=1}{\overset{12}{}}c_{cgt3,j}(\lambda _{cgt3,j})^{L_x}$$
(7.5)
where
$$\lambda _{cgt3,j}=\pm 1\mathrm{for}j=1,2$$
(7.6)
$$\lambda _{cgt3,j}=\pm 2\mathrm{for}j=3,4$$
(7.7)
$$\lambda _{cgt3,j}=\pm (2q7)\mathrm{for}j=5,6$$
(7.8)
$$\lambda _{cgt3,7}=2q$$
(7.9)
$$\lambda _{cgt3,8}=3q14$$
(7.10)
$$\lambda _{cgt3,9}=2(q4)^2$$
(7.11)
$$\lambda _{cgt3,j}=\pm (q^25q+7)\mathrm{for}j=10,11$$
(7.12)
$$\lambda _{cgt3,12}=q^39q^2+29q32.$$
(7.13)
The corresponding coefficients are
$$c_{cgt3,j}=\frac{1}{2}q(q1)(2q7)\mathrm{for}j=1,2$$
(7.14)
$$c_{cgt3,3}=\frac{1}{6}(q1)(q2)(q3)$$
(7.15)
$$c_{cgt3,4}=\frac{1}{6}q(q1)(q5)$$
(7.16)
$$c_{cgt3,j}=\frac{1}{2}(q1)(q2)\mathrm{for}5j7$$
(7.17)
$$c_{cgt3,8}=\frac{1}{2}q(q3)$$
(7.18)
$$c_{cgt3,j}=q1\mathrm{for}9j11$$
(7.19)
$$c_{cgt3,12}=1.$$
(7.20)
Thus, $`N_{cgt3,up,\lambda }=4`$, $`N_{cgt3,opd,\lambda }=2`$, and $`N_{cgt3,ops,\lambda }=6`$. Hence, for $`L_x=0`$ mod 6, where the chromatic polynomial reduces to that for the $`L_y=3`$ toroidal strip of the triangular lattice, the number of $`\lambda _j`$’s is reduced, according to the general formula (7.2), to $`N_{tt3,\lambda }=121=11`$. These numbers agree with our previous result ,
$$P(tri,3\times L_x,torus)=\underset{j=1}{\overset{11}{}}c_{tt3,j}(\lambda _{tt3,j})^{L_x}$$
(7.21)
where
$$\lambda _{tt3,1}=2$$
(7.22)
$$\lambda _{tt3,2}=q2$$
(7.23)
$$\lambda _{tt3,3}=3q14$$
(7.24)
$$\lambda _{tt3,4}=2(q4)^2$$
(7.25)
$$\lambda _{tt3,5}=q^39q^2+29q32$$
(7.26)
$$\lambda _{tt3,j}=e^{\pm 2\pi i/3}\mathrm{for}j=6,7$$
(7.27)
$$\lambda _{tt3,j}=(q^25q+7)e^{\pm 2\pi i/3}\mathrm{for}j=8,9$$
(7.28)
$$\lambda _{tt3,j}=(2q7)e^{\pm 2\pi i/3}\mathrm{for}j=10,11$$
(7.29)
with coefficients
$$c_{tt3,1}=\frac{1}{3}(q1)(q^25q+3)$$
(7.30)
$$c_{tt3,j}=\frac{1}{2}(q1)(q2)\mathrm{for}j=2,10,11$$
(7.31)
$$c_{tt3,3}=\frac{1}{2}q(q3)$$
(7.32)
$$c_{tt3,j}=q1\mathrm{for}j=4,8,9$$
(7.33)
$$c_{tt3,5}=1$$
(7.34)
$$c_{tt3,6}=c_{tt3,7}=\frac{1}{6}q(q1)(2q7).$$
(7.35)
Similarly, for odd $`L_x`$ where the $`L_y=3`$ crossing-subgraph strip reduces to the $`L_y=3`$ Klein bottle strip of the triangular lattice, the number of nonzero terms is reduced, according to the general formula (7.3), to $`N_{tk3,\lambda }=1216=5`$, in agreement with our previous calculation
$$P(tri,3\times L_x,PBC_y,TPBC_x)=\underset{j=1}{\overset{5}{}}c_{tk3,j}(\lambda _{tk3,j})^{L_x}$$
(7.36)
where
$$\lambda _{tk3,1}=\lambda _{tt3,1}=2$$
(7.37)
$$\lambda _{tk3,2}=\lambda _{tt3,2}=q2$$
(7.38)
$$\lambda _{tk3,3}=\lambda _{tt3,3}=3q14$$
(7.39)
$$\lambda _{tk3,4}=\lambda _{tt3,4}=2(q4)^2$$
(7.40)
$$\lambda _{tk3,5}=\lambda _{tt3,5}=q^39q^2+29q32$$
(7.41)
with coefficients
$$c_{tk3,1}=(q1)$$
(7.42)
$$c_{tk3,2}=\frac{1}{2}(q1)(q2)$$
(7.43)
$$c_{tk3,3}=\frac{1}{2}q(q3)$$
(7.44)
$$c_{tk3,4}=q1$$
(7.45)
$$c_{tk3,5}=1.$$
(7.46)
In the Appendix we give our results for the $`L_y=4`$ crossing-subgraph toroidal strip of the triangular lattice.
## 8 Concluding Discussion
We comment further here on some features of our results.
1. Our exact calculations of the singular loci $``$ for $`L_y=4`$ strips of the square and triangular lattice with toroidal or Klein bottle conditions exhibit the following features, as did the earlier calculations for the $`L_y=3`$ strips of these lattices in : $``$ (i) passes through $`q=0`$, (ii) passes through $`q=2`$, (iii) passes through a maximal real point, thereby defining a $`q_c`$, and (iv) encloses one or more regions including the interval $`0<q<q_c`$ . As noted above, we also found that these four features hold for the ($`L_x\mathrm{}`$ limit of) strips with cyclic and Möbius boundary conditions, which leads to the inference that the key condition is the existence of periodic (or reversed-orientation periodic) longitudinal boundary conditions.
2. Previous exact calculations of $``$ for cyclic and Möbius strips of the square and triangular lattice of various widths are consistent with the inference that as $`L_y`$ increases, the outer envelope of $``$ moves outward, i.e. if $`L_y>L_y^{}`$, then the outer envelope of $``$ for $`L_y`$ encloses that for $`L_y^{}`$ . In particular, $`q_c`$ is a nondecreasing function of $`L_y`$. However, we have also shown that neither of these properties holds for strips with $`(PBC_y,FBC_x)`$, i.e., cylindrical, boundary conditions . Our present results show that for strips of the square and triangular lattices with toroidal or Klein bottle boundary conditions, $`(PBC_y,(T)PBC_x)`$, the outer envelope of $``$ does not, in general, move monotonically outward as one increases the width. This is illustrated in Figs. 5 and 6, which show the boundaries $``$ for the $`L_y=3`$ and $`L_y=4`$ strips of, respectively, the square and triangular lattices with torus or Klein bottle boundary conditions. See also Table 2. This monotonic (nonmonotonic) behavior of the outer envelope is reminiscent of the monotonic (nonmonotonic) behavior of the $`W`$ function for free (periodic) transverse boundary conditions discussed in (see also ).
3. As a special aspect of this outer envelope, $`q_c`$ decreases from 3 to approximately 2.78 for the (infinite-length limit of the) square-lattice strip with toroidal or Klein bottle boundary conditions when one increase the width from $`L_y=3`$ to $`L_y=4`$. In contrast, $`q_c`$ increases from about 3.72 to 4 for the (infinite-length limit of the) triangular-lattice strip with toroidal or Klein bottle boundary conditions when one increases $`L_y`$ from 3 to 4. Related to this, the calculation of $`P`$ and $`W`$ and $``$ for the $`L_y=3`$ strip of the square lattice with toroidal or Klein bottle boundary conditions in showed that $`q_c=3`$ for (the $`L_x\mathrm{}`$ limit of) that strip, and hence showed that $`q_c`$ for a finite-width, infinite-length strip of a given lattice can be the same as for the limit of infinite width, i.e. the full 2D infinite lattice, since $`q_c=3`$ for the square lattice (for general upper bounds, see ). So far, this was an isolated example. Our calculation in the present paper provides a second example of this phenomenon: $`q_c`$ for the infinite-length limit of the $`L_y=4`$ strip of the triangular lattice with toroidal or Klein bottle boundary conditions has the value $`q_c=4`$, which is equal to the value for the full 2D triangular lattice, i.e. the $`L_y\mathrm{}`$ limit of the strip. Parenthetically, we note that rigorous bounds on $`q_c`$ have been given in .
4. In all of the cases of strips of the square and triangular lattice with periodic or reversed-orientation periodic longitudinal boundary conditions for which we have performed exact calculations of the chromatic polynomials and have determined the respective singular loci $``$, we have found the following results for the coefficients corresponding to the dominant terms in various regions: (i) in region $`R_1`$, this coefficient has been proved to be unity ; (ii) in the region containing the interval $`0<q<2`$, the coefficient is $`c^{(1)}=q1`$ for the strips of the square and triangular lattice with cyclic and torus b.c., the square strips with Möbius and Klein bottle b.c. and the triangular strips with Klein bottle b.c. (the coefficients are not, in general, polynomials for Möbius strips of the triangular lattice); and (iii) for the observed complex-conjugate pairs of regions, the coefficients are also $`q1`$ for cyclic and torus b.c. and $`\pm (q1)`$ for Möbius ($`sq`$ case) and Klein bottle b.c. . A fourth finding is that (iv) for the torus/Klein bottle boundary conditions, in the cases that we have studied, we have found that the coefficient corresponding to the dominant $`\lambda _{G,j}`$ in the region containing the interval $`2<q<q_c`$, for each respective $`q_c`$, is $`q(q3)/2`$.
Acknowledgment: The research of R. S. was supported in part by the NSF grant PHY-9722101 and at Brookhaven by the DOE contract DE-AC02-98CH10886.<sup>3</sup><sup>3</sup>3Accordingly, the U.S. government retains a non-exclusive royalty-free license to publish or reproduce the published form of this contribution or to allow others to do so for U.S. government purposes.
## 9 Appendix: Further Results on Crossing-Subgraph Strips
### 9.1 $`L_y=4`$ Toroidal Crossing-Subgraph Strip of the Square Lattice
For the $`L_y=4`$ crossing-subgraph toroidal strip of the square lattice (labelled $`cgst4`$) we find that $`N_{cgst4,\lambda }=48`$ and calculate
$$P(cg(sq,4\times L_x,torus))=\underset{j=1}{\overset{48}{}}c_{cgst4,j}(\lambda _{cgst4,j})^{L_x}$$
(9.1.1)
where
$$\lambda _{cgst4,j}=\pm 1\mathrm{for}j=1,2$$
(9.1.2)
$$\lambda _{cgst4,j}=\pm (q1)\mathrm{for}j=3,4$$
(9.1.3)
$$\lambda _{cgst4,j}=\pm (q2)\mathrm{for}j=5,6$$
(9.1.4)
$$\lambda _{cgst4,j}=\pm (q3)\mathrm{for}j=7,8$$
(9.1.5)
$$\lambda _{cgst4,j}=\pm (q4)\mathrm{for}j=9,10$$
(9.1.6)
$$\lambda _{cgst4,j}=\pm (q5)\mathrm{for}j=11,12$$
(9.1.7)
$$\lambda _{cgst4,13}=(q^25q+5)$$
(9.1.8)
$$\lambda _{cgst4,14}=q^25q+7$$
(9.1.9)
$$\lambda _{cgst4,15}=(q1)(q3)$$
(9.1.10)
$$\lambda _{cgst4,j}=\lambda _{st4,j6}\mathrm{for}j=16,17$$
(9.1.11)
(where $`\lambda _{st4,j}`$ for $`j=10,11`$ were given in the text in eq. (2.16)),
$$\lambda _{cgst4,j}=\pm \lambda _{st4,12}\mathrm{for}j=18,19$$
(9.1.12)
$$\lambda _{cgst4,j}=\pm \lambda _{st4,13}\mathrm{for}j=20,21$$
(9.1.13)
$$\lambda _{cgst4,j}=\lambda _{st4,j8}\mathrm{for}j=22,23$$
(9.1.14)
$$\lambda _{cgst4,j}=\pm \lambda _{st4,16}\mathrm{for}j=24,25$$
(9.1.15)
$$\lambda _{cgst4,j}=\pm \lambda _{st4,17}\mathrm{for}j=26,27$$
(9.1.16)
$$\lambda _{cgst4,j}=\lambda _{st4,j10}\mathrm{for}j=28,29$$
(9.1.17)
$$\lambda _{cgst4,j}=\pm \lambda _{st4,20}\mathrm{for}j=30,31$$
(9.1.18)
$$\lambda _{cgst4,j}=\pm \lambda _{st4,21}\mathrm{for}j=32,33$$
(9.1.19)
$$\lambda _{cgst4,j}=\lambda _{st4,j12}\mathrm{for}34j36$$
(9.1.20)
$$\lambda _{cgst4,j}=\lambda _{st4,j12}\mathrm{for}37j42$$
(9.1.21)
$$\lambda _{cgst4,j}=\pm \lambda _{st4,31}\mathrm{for}j=43,44$$
(9.1.22)
$$\lambda _{cgst4,j}=\pm \lambda _{st4,32}\mathrm{for}j=45,46$$
(9.1.23)
$$\lambda _{cgst4,j}=\pm \lambda _{st4,33}\mathrm{for}j=47,48.$$
(9.1.24)
The corresponding coefficients are
$$c_{cgst4,1}=\frac{1}{2}(q1)(q^37q^2+13q2)$$
(9.1.25)
$$c_{cgst4,2}=\frac{1}{2}q(q3)(q^25q+5)$$
(9.1.26)
$$c_{cgst4,j}=\frac{1}{6}q(q1)(q5)\mathrm{for}j=3,12$$
(9.1.27)
$$c_{cgst4,4}=\frac{1}{6}(q1)(q2)(q3)$$
(9.1.28)
$$c_{cgst4,j}=\frac{1}{3}q(q2)(q4)\mathrm{for}j=5,9,10\mathrm{and}30j33$$
(9.1.29)
$$c_{cgst4,6}=\frac{1}{6}(q2)(q3)(2q+1)$$
(9.1.30)
$$c_{cgst4,j}=\frac{1}{3}(q1)(q^25q+3)\mathrm{for}j=7,8$$
(9.1.31)
$$c_{cgst4,11}=\frac{1}{6}(q1)(q2)(q3)$$
(9.1.32)
$$c_{cgst4,j}=\frac{1}{2}(q1)(q2)\mathrm{for}j=13,28,29\mathrm{and}43j48$$
(9.1.33)
$$c_{cgst4,j}=\frac{1}{2}q(q3)\mathrm{for}j=14\mathrm{and}22j27,40j42$$
(9.1.34)
$$c_{cgst4,j}=1\mathrm{for}j=15,16,17$$
(9.1.35)
$$c_{cgst4,j}=q1\mathrm{for}18j21\mathrm{and}34j39.$$
(9.1.36)
Thus, $`N_{cgst4,opd,\lambda }=8`$ and $`N_{cgst4,ops,\lambda }=22`$ and hence for the even and odd $`L_x`$ values where the chromatic polynomial for this crossing-subgraph strip reduces to the respective chromatic polynomial for the $`L_y=4`$ strip with torus and Klein bottle boundary conditions, we have, from eqs. (6.1.17) and (6.1.18), $`N_{st4,\lambda }=48(1/2)(8+22)=33`$ and $`N_{sk4,\lambda }=48(1/2)822=22`$, in agreement with our calculations in the text.
### 9.2 $`L_y=4`$ Toroidal Crossing-Subgraph Strip of the Triangular Lattice
For the $`L_y=4`$ crossing-subgraph toroidal strip of the triangular lattice (labelled $`cgt4`$) we find that there are $`N_{cgt4,\lambda }=40`$ different nonzero $`\lambda _{cgt4,j}`$ terms that enter into the chromatic polynomial and that the coloring matrix also has a zero eigenvalue, so that the total number of eigenvalues of the coloring matrix for this strip is $`N_{cgt4,\lambda ,tot}=41`$. We calculate
$$P(cg(tri,4\times L_x,torus))=\underset{j=1}{\overset{40}{}}c_{cgt4,j}(\lambda _{cgt4,j})^{L_x}$$
(9.2.1)
where
$$\lambda _{cgt4,j}=\pm 2\mathrm{for}j=1,2$$
(9.2.2)
$$\lambda _{cgt4,j}=\pm \sqrt{2}\mathrm{for}j=3,4$$
(9.2.3)
$$\lambda _{cgt4,5}=2(3q)$$
(9.2.4)
$$\lambda _{cgt4,j}=\pm (3q)\mathrm{for}j=6,7$$
(9.2.5)
$$\lambda _{cgt4,j}=\pm 2(2q9)\mathrm{for}j=8,9$$
(9.2.6)
$$\lambda _{cgt4,10}=2(q3)^2$$
(9.2.7)
$$\lambda _{cgt4,j}=\lambda _{tt4,j3}\mathrm{for}j=11,12$$
(9.2.8)
where $`\lambda _{tt4,j}`$ for $`j=8,9`$ were given above in eqs. (4.11),
$$\lambda _{cgt4,j}=\pm \sqrt{3}(q3)\mathrm{for}j=13,14$$
(9.2.9)
$$\lambda _{cgt4,j}=\pm (q2)\sqrt{2(q3)(q4)}\mathrm{for}j=15,16$$
(9.2.10)
$$\lambda _{cgt4,j}=\lambda _{tt4,j3}\mathrm{for}17j22.$$
(9.2.11)
The twelve terms $`\lambda _{cgt4,j}`$ for $`23j34`$ are related to the $`\lambda _{tt4,j}`$’s that are the roots of the quartic equations (4.24)-(4.34) as follows, where the $`a_{\mathrm{}}`$ and $`b_{\mathrm{}}`$ were defined in eqs. (4.42)-(4.47):
$$\lambda _{cgt4,j}=\pm \sqrt{2}a_1\mathrm{for}j=23,24$$
(9.2.12)
$$\lambda _{cgt4,j}=\pm \sqrt{2}b_1\mathrm{for}j=25,26$$
(9.2.13)
$$\lambda _{cgt4,j}=\pm \sqrt{2}a_2\mathrm{for}j=27,28$$
(9.2.14)
$$\lambda _{cgt4,j}=\pm \sqrt{2}b_2\mathrm{for}j=29,30$$
(9.2.15)
$$\lambda _{cgt4,j}=\pm \sqrt{2}a_3\mathrm{for}j=31,32$$
(9.2.16)
$$\lambda _{cgt4,j}=\pm \sqrt{2}b_3\mathrm{for}j=33,34.$$
(9.2.17)
The six terms $`\lambda _{cgt4,j}`$ for $`35j40`$ are related to the $`\lambda _{tt4,j}`$’s that are the roots of the sixth-degree equation (4.41) as follows, where the $`c_{\mathrm{}}`$ were defined in eqs. (4.48)-(4.50):
$$\lambda _{cgt4,j}=\pm \sqrt{2}c_1\mathrm{for}j=35,36$$
(9.2.18)
$$\lambda _{cgt4,j}=\pm \sqrt{2}c_2\mathrm{for}j=37,38$$
(9.2.19)
$$\lambda _{cgt4,j}=\pm \sqrt{2}c_3\mathrm{for}j=39,40.$$
(9.2.20)
The corresponding coefficients are
$$c_{cgt4,1}=\frac{1}{8}q(q1)(q2)(q3)$$
(9.2.21)
$$c_{cgt4,2}=\frac{1}{8}q(q2)(q3)(q5)$$
(9.2.22)
$$c_{cgt4,j}=\frac{1}{12}(q1)(q2)(3q^211q6)\mathrm{for}j=3,4$$
(9.2.23)
$$c_{cgt4,j}=\frac{1}{2}(q1)(q2)\mathrm{for}j=5,10\mathrm{and}35j40$$
(9.2.24)
$$c_{cgt4,j}=\frac{1}{3}q(q2)(q4)\mathrm{for}j=6,7,13,14\mathrm{and}31j34$$
(9.2.25)
$$c_{cgt4,8}=\frac{1}{6}(q1)(q2)(q3)$$
(9.2.26)
$$c_{cgt4,9}=\frac{1}{6}q(q1)(q5)$$
(9.2.27)
$$c_{cgt4,j}=1\mathrm{for}j=11,12$$
(9.2.28)
$$c_{cgt4,j}=\frac{1}{2}q(q3)\mathrm{for}j=15,16\mathrm{and}20j22,27j30$$
(9.2.29)
$$c_{cgt4,j}=q1\mathrm{for}17j19\mathrm{and}23j26.$$
(9.2.30)
Finally, the coloring matrix has a zero eigenvalue,
$$\lambda _{cgt4,41}=0$$
(9.2.31)
with multiplicity
$$c_{cgt4,41}=\frac{1}{12}q(q1)(3q^217q+40).$$
(9.2.32)
It follows from the general relation (7.1) that for $`L_x=0`$ mod 8, this $`L_y=4`$ crossing-subgraph strip of the triangular lattice reduces to the regular $`L_y=4`$ toroidal strip of the triangular lattice. This gives insight into the occurrence of the phase factors in several of the $`\lambda _{tt4,j}`$ terms. For this strip, we have $`N_{cgt4,up,\lambda }=2`$, $`N_{cgt4,opd,\lambda }=4`$, and $`N_{cgt4,ops,\lambda }=26`$. Further, $`N_{cgt4,ops,r,\lambda }=2`$ and $`N_{cgt4,ops,i,\lambda }=24`$. Hence, by the general relation (7.2), for $`L_x=0`$ mod 8, where the $`L_y=4`$ crossing-subgraph strip reduces to the toroidal strip of the triangular lattice, the number of nonzero terms is reduced to $`N_{tt4,\lambda }=4021=37`$. For $`L_x=1`$ mod 4 where the $`L_y=4`$ crossing-subgraph strip reduces to the Klein bottle strip of the triangular lattice, the number of nonzero terms is reduced, according to the general formula (7.3), to $`N_{tk4,\lambda }=40226=12`$. These numbers agree with our exact calculations presented in the text.
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# 1 Introduction
## 1 Introduction
As far as I know this is the first international conference devoted entirely to the relation of spin and statistics and to the investigation of possible small violations of statistics. My purpose in this talk is to give an overview of theoretical issues connected with violations of statistics. I have divided the talk into five parts: (1) general theoretical remarks, (2) types of experiments to detect violations, (3) attempts to violate statistics, (4) quons, the best formalism so far to describle small violations, and (5) summary and open questions.
## 2 General theoretical remarks
The general principles of quantum theory do not require that all particles be either bosons or fermions. This restriction requires an additional postulate which A.M.L. Messiah named the “symmetrization postulate,” which I quote as Messiah defined it: “The states of a system containing $`N`$ identical particles are necessaribly either all symmetrical or all antisymmetrical with respect to permutations of the $`N`$ particles.” The symmetrization postulate can be restated as: “All states of identical particles are in one-dimensional representations of the symmetric group.” With Messiah’s definition the spin-statistics connection, that integer spin particles are bosons and odd-half-integer spin particles are fermions, is a separate statement.
Messiah and I gave a detailed discussion of quantum mechanics without the symmetrization postulate. We emphasized that without the symmetrization postulate, a set of one-body measurements is never a maximal set; one needs additional measurements to fix the state of the system. Further there is a superselection rule separating states in inequivalent representations of the symmetric group. If identical particles can occur in states that violate the spin-statistics connection their transitions must occur in the same representation of the symmetric group. For example, in the experiment of Deilamian, et al which looked for anomalous helium atoms in which the two electrons violated the exclusion principle and were in the symmetric state, the search was for transitions among the symmetric states rather than between symmetric and antisymmetric states. This point was also made by R. Amado and H. Primakoff. In addition if one assumes that charged particles couple universally to the electromagnetic field, then the transitions among the anomalous states occur at the normal rate, so that isolated atoms will be in the lowest state of the anomalous system. Since the symmetrization postulate is not an intrinsic part of quantum theory, this postulate must be subjected both to theoretical study and to experimental tests. Quantitative tests require a theory in which the symmetrization postulate does not have to hold and in which the violation of the postulate is reflected in a parameter that departs from its standard value at which the symmetrization postulate and the spin-statistics connection do hold. It is certainly possible that violations of statistics are extremely small and require high-precision tests to be observed. This conference brings together leading workers in this search.
Because the notion that particles are identical requires that the Hamiltonian and, indeed, all observables must be symmetric in the dynamical variables associated with the identical particles, the observables can’t change the permutation symmetry type of the wave function. In particular one can’t introduce a small violation of statistics by assuming the Hamiltonian is the sum of a statistics-conserving and a small statistics-violating term,
$$H=H_S+ϵH_V$$
(1)
as one can for violations of parity, charge conjugation, etc. Violation of statistics has to be introduced in a more subtle way.
If charged particles couple universally to the electromagnetic field, then there can’t be two kinds of–say–electrons, “red” electrons and “blue” electrons, because then the lowest order pair production cross section,
$$\sigma (\gamma Xe^+e^{}X)$$
(2)
would double. A high-precision measurement is not needed to rule this out.
A convenient way to parametrize violations or bounds on violations of statistics uses the two-particle density matrix. For fermions,
$$\rho _2=(1v_F)\rho _a+v_F\rho _s;$$
(3)
for bosons,
$$\rho _2=(1v_B)\rho _s+v_B\rho _a;$$
(4)
in each case the violation parameter varies between zero if the statistics is not violated and one if the statistics is completely violated.
## 3 Types of experiments
There are three basic types of experiments to detect violations of statistics: (1) transitions among anomalous states–these can occur in solids, liquids or gases, (2) accumulation of particles in anomalous states, and (3) deviations from the usual statistical properties of the identical particles. Since, as mentioned earlier, a superselection rule prevents transitions between normal and anomalous states, experiments searching for such transitions do not provide a valid test of violation of statistics.
Transitions among anomalous states can provide a very sensitive test, since in some cases a single such transition can be observed. The prototype of this kind of test is the experiment of Maurice and Trudy Goldhaber. They asked the qualitative question, “Do the electrons from nuclear beta decay obey the same exclusion principle as electrons in atoms?” They knew that electrons from each source have the same charge, spin, and mass, etc., i.e. that the single-electron states in each case are identical, but there was no evidence that the many-electron states from each source are identical. They devised the following ingenious test: they let beta decay electrons from a nuclear source fall on a block of lead. They argued that if the many-electron states were not identical then the nuclear beta decay electrons would not obey the same exclusion principle as the electrons in the lead atoms. Then the beta decay electrons would not see the $`K`$ shells in the lead atoms as filled and could fall into the $`K`$ shells and would emit $`x`$-rays. A single such $`x`$-ray could be observed. They saw no such $`x`$-rays above background and thus answered their qualitative question in the affirmative. I estimate that their experiment gave the bound $`v_F5\times 10^2`$ for electrons.
E. Ramberg and G.A. Snow developed this experiment into one which yields a high-precision bound on violations of the exclusion principle. Their idea was to replace the natural $`\beta `$ source, which provides relatively few electrons, by an electric current, in which case Avogadro’s number is on their side. The possible violation of the exclusion principle is that a given collection of electrons can, with different probabilities, be in different permutation symmetry states. The probability to be in the “normal” totally antisymmetric state presumably would be close to one, the next largest probability would occur for the state with its Young tableau having one row with two boxes, etc. The idea of the experiment is that each collection of electrons has a possibility of being in an “abnormal” permutation state. If the density matrix for a conduction electron together with the electrons in an atom has a projection onto such an “abnormal” state, then the conduction electron will not see the $`K`$ shell of that atom as filled. Then a transition into the $`K`$ shell with $`x`$-ray emission is allowed. Each conduction electron which comes sufficiently close to a given atom has an independent chance to make such an $`x`$-ray-emitting transition, and thus the probability of seeing such an $`x`$-ray is proportional to the number of conduction electrons which traverse the sample and the number of atoms which the electrons visit, as well as the probability that a collection of electrons can be in the anomalous state. Ramberg and Snow chose to run 30 amperes through a thin copper strip for about a month. They surrounded the experiment with veto scintillators to remove background $`x`$-rays. They estimated the energy of the modified $`x`$-rays which would be emitted due to the transition to the $`K`$ shell. No excess of $`x`$-rays above background was found in this energy region. Ramberg and Snow set the limit
$$v_F1.7\times 10^{26}$$
(5)
for electrons. This is high precision indeed!
The Ramberg-Snow experiment may seem discouraging for the discovery of generalizations of bose and fermi statistics; however there are small numbers in physics which, if necessary, can occur in degree greater than one. For example the ratios
$$\frac{m_{proton}}{M_{Planck}}10^{19},\mathrm{and}\frac{G_Nm_e^2}{e^2}10^{43}$$
(6)
can provide numbers smaller than the Ramberg-Snow bound. In addition new physics effects such as violations of Lorentz invariance, spacetime discreteness, spacetime noncommutativity, etc. may provide small effects. Mohapatra and I gave an early survey of experimental bounds on violations of statistics.
Composite structure can mimic violations of statistics. This is not what I am considering here.
## 4 Attempts to violate statistics
### 4.1 Gentile’s “intermediate statistics”
The first attempt to go beyond bose and fermi statistics seems to have been made by G. Gentile who suggested an “intermediate statistics” in which at most $`n`$ identical particles could occupy a given quantum state. In intermediate statistics, fermi statistics is recovered for $`n=1`$ and bose statistics is recovered for $`n\mathrm{}`$; thus intermediate statistics interpolates between fermi and bose statistics. However Gentile’s statistics is not a proper quantum statistics, because the condition of having at most $`n`$ particles in a given quantum state is not invariant under change of basis. For example, for intermediate statistics with $`n=2`$, the state $`|\psi =|k,k,k`$ does not exist; however, the state $`|\chi =_{l_1,l_2,l_3}U_{k,l_1}U_{k,l_2}U_{k,l_3}|l_1,l_2,l_3`$, obtained from $`|\psi `$ by the unitary change of single-particle basis, $`|k^{}=_lU_{k,l}|l`$ does exist. By contrast, parafermi statistics of order $`n`$ which I discuss just below is invariant under change of basis. Parafermi statistics of order $`n`$ not only allows at most $`n`$ identical particles in the same state, but also allows at most $`n`$ identical particles in a symmetric state. In the example just described, neither $`|\psi `$ nor $`|\chi `$ exist for parafermi statistics of order two.
### 4.2 Parastatistics
H.S. Green proposed the first proper quantum statistical generalization of bose and fermi statistics. Green noticed that the commutator of the number operator with the annihilation and creation operators is the same for both bosons and fermions
$$[n_k,a_l^{}]_{}=\delta _{kl}a_l^{}.$$
(7)
The number operator can be written
$$n_k=(1/2)[a_k^{},a_k]_\pm +\mathrm{const}.,$$
(8)
where the anticommutator (commutator) is for the bose (fermi) case. If these expressions are inserted in the number operator-creation operator commutation relation, the resulting relation is trilinear in the annihilation and creation operators. Polarizing the number operator to get the transition operator $`n_{kl}`$ which annihilates a free particle in state $`l`$ and creates one in state $`k`$ leads to Green’s trilinear commutation relation for his parabose and parafermi statistics,
$$[[a_k^{},a_l]_\pm ,a_m^{}]_{}=2\delta _{lm}a_k^{}$$
(9)
Since these rules are trilinear, the usual vacuum condition,
$$a_k|0=0,$$
(10)
does not suffice to allow calculation of matrix elements of the $`a`$’s and $`a^{}`$’s; a condition on single-particle states must be added,
$$a_ka_l^{}|0=p\delta _{kl}|0.$$
(11)
Green found an infinite set of solutions of his commutation rules, one for each positive integer $`p`$, by giving an ansatz which he expressed in terms of bose and fermi operators. Let
$$a_k^{}=\underset{p=1}{\overset{n}{}}b_k^{(\alpha )},a_k=\underset{p=1}{\overset{n}{}}b_k^{(\alpha )},$$
(12)
and let the $`b_k^{(\alpha )}`$ and $`b_k^{(\beta )}`$ be bose (fermi) operators for $`\alpha =\beta `$ but anticommute (commute) for $`\alpha \beta `$ for the “parabose” (“parafermi”) cases. This ansatz clearly satisfies Green’s relation. The integer $`p`$ is the order of the parastatistics. The physical interpretation of $`p`$ is that, for parabosons, $`p`$ is the maximum number of particles that can occupy an antisymmetric state, while for parafermions, $`p`$ is the maximum number of particles that can occupy a symmetric state (in particular, the maximum number which can occupy the same state). The case $`p=1`$ corresponds to the usual bose or fermi statistics. Later, Messiah and I proved that Green’s ansatz gives all Fock-like solutions of Green’s commutation rules. Local observables have a form analogous to the usual ones; for example, the local current for a spin-1/2 theory is $`j_\mu =(1/2)[\overline{\psi }(x),\psi (x)]_{}`$. From Green’s ansatz, it is clear that the squares of all norms of states are positive, since sums of bose or fermi operators give positive norms. Thus parastatistics gives a set of orthodox theories.
This is all well and good; however, the violations of statistics provided by parastatistics are gross. Parafermi statistics of order two has up to two particles in each quantum state. High-precision experiments are not necessary to rule this out for all particles we think are fermions.
### 4.3 The Ignatiev-Kuzmin model
Interest in possible small violations of the exclusion principle was revived by a paper of Ignatiev and Kuzmin in 1987. They constructed a model of one oscillator with three possible states: a vacuum state, a one-particle state and, with small amplitude $`\beta `$, a two-particle state. They gave trilinear commutation relations for their oscillator. Mohapatra and I noticed that the Ignatiev-Kuzmin oscillator could be represented by a modified form of the order-two Green ansatz. We suspected that a field theory generalization of this model having an infinite number of oscillators would not have local observables and set about trying to prove this. To our surprize, we found that we could construct local observables and gave trilinear relations which guarantee the locality of the current.
### 4.4 Parons
Following Ignatiev and Kuzmin we introduced a parameter $`\beta `$ that gives the deformation of the Green trilinear commutation relations. For $`\beta 1`$ the relations reduce to those of the $`p=2`$ parafermi field; for $`\beta 0`$ the double occupancy is completely suppressed and the theory is equivalent to a fermi theory. A random state of two paronic electrons has the violation parameter $`\beta ^2/2`$. Mohapatra and I checked that the norms are positive for states of up to three particles. At this stage, we were carried away with enthusiasm, named these particles “parons” since their algebra is a deformation of the parastatistics algebra, and thought we had found a local theory with small violation of the exclusion principle. Unknown to us Govorkov, using a detailed algebraic argument, already had shown in generality that any deformation of the Green commutation relations necessarily has states with negative squared norms in the Fock-like representation. For our model, the first such negative-probability state occurs for four particles in the representation of $`𝒮_4`$ with three boxes in the first row and one in the second. We were able to understand Govorkov’s result qualitatively as follows: Since parastatistics of order $`p`$ is related by a Klein transformation to a model with exact $`SO(p)`$ or $`SU(p)`$ internal symmetry, a deformation of parastatistics which interpolates between Fermi and parafermi statistics of order two would be equivalent to interpolating between the trivial group whose only element is the identity and a theory with $`SO(2)`$ or $`SU(2)`$ internal symmetry. This is impossible, since there is no such interpolating group.
### 4.5 The Doplicher-Haag-Roberts analysis
S. Doplicher, R. Haag and J. Roberts made a general study of identical particle statistics using the algebraic field theory methods pioneered by Haag. They found parabose and parafermi statistics of positive integer orders which as mentioned above were introduced by Green. They also found another case which they called infinite statistics. Young patterns label the inequivalent irreducible representations of the symmetric group. In parabose (parafermi) statistics of order $`p`$ the Young patterns have at most $`p`$ rows (columns) corresponding to having at most $`p`$ particles in an antisymmetric (a symmetric) state. In infinite statistics all irreducibles of the symmetric group occur. Doplicher, et al, did not give an operator realization of infinite statistics.
### 4.6 Infinite statistics
In 1989 I gave an evening lecture at Wake Forest University. My talk was attended by physicists, philosophers, and among people in other disciplines, chemists. In my talk I mentioned the bose and fermi commutation relations. After the talk Roger Hegstrom, a chemist, asked “Why not average the bose and fermi commutation relations and consider the relation
$$a(k)a^{}(l)=\delta (k,l)\mathrm{?}\mathrm{"}$$
(13)
I was surprized to find that such a simple case had not been considered. Later I found out that it had, in the mathematical literature, by J. Cuntz. With Hegstrom’s permission I developed this case, which turned out to be the first operator example of infinite statistics. In order to select the Fock-like representation, one must add the vacuum condition
$$a(k)|0=0.$$
(14)
We can calculate all vacuum matrix elements of products of $`a`$’s and $`a^{}`$’s using the commutation relation and the vacuum condition. There is no commutation relation involving two $`a`$’s or two $`a^{}`$’s. There are $`n!`$ linearly independent $`n`$-particle states in Hilbert space if all quantum numbers are distinct; these states differ only by permutations of the order of the creation operators. (Later we will see that there are not that many independent density matrices or other observables.) The matrix of scalar products of these states is the identity matrix,
$$M_{P,Q}^n(q)=(Pa^{}(k_1)a^{}(k_2)\mathrm{}a^{}(k_n)|0,Qa^{}(l_1)a^{}(l_2)\mathrm{}a^{}(l_n)|0=_{i=1}^n\delta (k_i,l_i)\delta (P,Q)$$
(15)
where $`P`$ and $`Q`$ are permutations from $`S_n`$; that is, the scalar product is zero unless there are the same number of creation operators on each side of the scalar product and they have the same quantum numbers in the same order. This algebra can be viewed as a deformation of either the bose or the fermi algebras. As is typical for deformed algebras, there is an element that is infinite degree in the generators of the algebra. In this case the number operator that obeys
$$[n(k),a^{}(l)]_{}=\delta (k,l)a^{}(l)$$
(16)
is the operator of infinite degree; in terms of the $`a`$’s and the $`a^{}`$’s,
$$n(k)=a^{}(k)a(k)+\underset{t}{}a^{}(t)a^{}(k)a(k)a(t)+\underset{t_1,t_2}{}a^{}(t_2)a^{}(t_1)a^{}(k)a(k)a(t_1)a(t_2)+\mathrm{}.$$
(17)
There is an analogous formula for the transition operator, $`n(k,l)`$, that obeys
$$[n(k,l),a^{}(m)]_{}=\delta (l,m)a^{}(k).$$
(18)
## 5 Quons
### 5.1 The quon algebra
The quon algebra is the best attempt so far to violate statistics by a small amount. The infinite statistics algebra just discussed is the average of the bose and fermi algebras. The quon algebra can be obtained as the convex sum of these two algebras,
$$\frac{1+q}{2}[a(k),a^{}(l)]_{}+\frac{1q}{2}[a(k),a^{}(l)]_+=\delta (k,l),$$
(19)
or
$$a(k)a^{}(l)qa^{}(l)a(k)=\delta (k,l).$$
(20)
As usual the Fock-like representation is selected by the vacuum condition
$$a(k)|0=0.$$
(21)
Convexity requires $`0q1`$; for this range the states have positive squared norms. Outside this range the squared norms become negative. Using the algebra (20) and the vacuum condition (21) all vacuum matrix elements of polynomials in the $`a`$’s and $`a^{}`$’s can be calculated; for example,
$$(a^{}(k_1)a^{}(k_2)|0,a^{}(l_1)a^{}(l_2)|0)=$$
$$\delta (k_1,l_1)\delta (k_2,l_2)+q\delta (k_1,l_2)\delta (k_2,l_1)=$$
$$\frac{1+q}{2}[\delta ((k_1,l_1)\delta (k_2,l_2)+\delta (k_1,l_2)\delta (k_2,l_1)]+\frac{1q}{2}[\delta ((k_1,l_1)\delta (k_2,l_2)\delta (k_1,l_2)\delta (k_2,l_1)].$$
(22)
The first proof of the positivity of the norms was given by D. Zagier, who gave a tour-de-force calculation of the determinant of the $`n!\times n!`$ matrix of scalar products (15) for arbitrary $`n`$,
$$detM_{P,Q}^n(q)=\underset{k=1}{\overset{n1}{}}(1q^{k(k+1)})^{\frac{nk)n!}{k(k+1)}}.$$
(23)
As shown above, at $`q=0`$ the norms are positive and the determinant is one. In order for a norm to become negative the determinant has to change sign. From Zagier’s formula this happens only when $`q^{k(k+1)}=1`$, i.e., on the unit circle. This proves that the norms remain positive between negative one and one.
Speicher gave an ingenious proof of the positivity of the norms using an ansatz for the Fock-like representation of quons analogous to Green’s ansatz for parastatistics. Speicher represented the quon annihilation operator as the weak operator limit,
$$a_k=\mathrm{lim}_N\mathrm{}N^{1/2}\underset{\alpha =1}{\overset{N}{}}b_k^{(\alpha )},$$
(24)
where the $`b_k^{(\alpha )}`$ are bose oscillators for each $`\alpha `$, but with relative commutation relations given by
$$b_k^{(\alpha )}b_l^{(\beta )}=s^{(\alpha ,\beta )}b_l^{(\beta )}b_k^{(\alpha )},\alpha \beta ,\mathrm{where}s^{(\alpha ,\beta )}=\pm 1.$$
(25)
This limit is taken as the limit, $`N\mathrm{}`$, in the vacuum expectation state of the Fock space representation of the $`b_k^{(\alpha )}`$. In this respect Speicher’s ansatz differs from Green’s, which is an operator identity. To get the Fock-like representation of the quon algebra, Speicher chose a probabilistic condition for the signs $`s^{(\alpha ,\beta )}`$,
$$\mathrm{prob}(s^{(\alpha ,\beta )}=1)=(1+q)/2,$$
(26)
$$\mathrm{prob}(s^{(\alpha ,\beta )}=1)=(1q)/2.$$
(27)
Speicher’s rules reproduce the quon algebra. The norms are positive since the sums of bose or fermi operators have positive norms. The constraint on $`q`$ follows because the probabilities have to lie between zero and one.
The number and transition operators for general $`q`$ have infinite degree expansions analogous to but more complicated than those for the $`q=0`$ case,
$$n(k,l)=a^{}(k)a(l)+\frac{1}{1q^2}\underset{t}{}(a^{}(t)a^{}(k)qa^{}(k)a^{}(t))(a(l)a(t)qa(t)a(l))+\mathrm{}.$$
(28)
The general formula for the number operator was given by S. Stanciu.
At $`q=\pm 1`$ only the symmetric (antisymmetric) representation of $`𝒮_n`$ occurs. The quon operators interpolate smoothly between fermi and bose statistics in the sense that as $`q`$ departs from $`\pm 1`$ the vectors formed by polynomials in the creation operators, which are superpositions of vectors in different irreducible representations of the symmetric group, have higher weights in the more symmetric (antisymmetric) representations, and as $`q1`$ the antisymmetric (symmetric) representations smoothly become more heavily weighted.
### 5.2 Observables in quon theory
It is important to note that although there are $`n!`$ linearly independent vectors in Fock space associated with a degree $`n`$ monomial in creation operators that carry disjoint quantum numbers acting on the vacuum, there are fewer than $`n!`$ observables associated with such vectors. For example, for two identical quons $`1`$ and $`2`$ in orthogonal quantum states the two vectors $`a^{}(1)a^{}(2)|0`$ and $`a^{}(2)a^{}(2)|0`$ are orthogonal and each is normalized to one. Let
$$|\varphi _{s,a}=N_{s,a}(a^{}(1)a^{}(2)\pm a^{}(2)a^{}(1))|0$$
(29)
be normalized states that are symmetric or antisymmetric under transposition of $`1`$ and $`2`$. The quon algebra gives
$$N_{s,a}=\frac{1}{\sqrt{2(1\pm q)}}.$$
(30)
One can then calculate the expansion
$$a^{}(1)a^{}(2)|0=\alpha |\varphi _s+\beta |\varphi _a,$$
(31)
either using
$$a^{}(1)a^{}(2)|0=(1/2)[(a^{}(1)a^{}(2)+a^{}(2)a^{}(1))+(a^{}(1)a^{}(2)a^{}(2)a^{}(1))]|0$$
(32)
or using
$$a^{}(1)a^{}(2)|0=\varphi _s|a_1^{}a_2^{}|0|\varphi _s+\varphi _a|a_1^{}a_2^{}|0|\varphi _a.$$
(33)
Either way gives
$$\alpha =\sqrt{(1+q)/2},\beta =\sqrt{(1q)/2}$$
(34)
so that
$$a_1^{}a_2^{}|0=\sqrt{\frac{1+q}{2}}\varphi _s+\sqrt{\frac{1q}{2}}\varphi _a$$
(35)
and
$$a_2^{}a_1^{}|0=\sqrt{\frac{1+q}{2}}\varphi _s\sqrt{\frac{1q}{2}}\varphi _a.$$
(36)
Then, dropping the cross terms that are excluded by the superselection rule separating symmetric and antisymmetric states of identical particles (and, indeed, states of identical particles in different representations of the symmetric group generally),
$$a_1^{}a_2^{}|00|a_2a_1=a_2^{}a_1^{}|00|a_1a_2=\frac{1+q}{2}|\varphi _s\varphi _s|+\frac{1q}{2}|\varphi _a\varphi _a|.$$
(37)
This shows that since particles 1 and 2 are identical the same observable results follow when the labels 1 and 2 are transposed. That also shows that the relative phase in Eq.(35) and Eq.(36) is not observable. Equation (37) states that the density matrices for $`a_1^{}a_2^{}|0`$ and $`a_2^{}a_1^{}|0`$ are identical, which means that these two “states” correspond to exactly the same physical situation. We put quotation marks around the word “states” to indicate that these should really be represented by density matrices. Note that the sum of the coefficients of the two terms in the two-particle density matrix is one, as it should be. The general observable is a linear combination of projectors on the irreducibles of the symmetric group.
The parameters $`v_F`$ and $`v_B`$ that represent small violations of statistics can be written in terms of the $`q`$ parameters; the result is,
$$q_F=2v_F1\mathrm{or}v_F=\frac{1}{2}(1+q_F);q_B=12v_B\mathrm{or}v_B=\frac{1}{2}(1q_B).$$
(38)
### 5.3 Properties of quon theory
Surprizingly several properties of relativistic theories that I would expect to fail for the quon theory (made relativistic kinematically) actually hold. These include Wick’s theorem, cluster decomposition theorems and the $`CPT`$ theorem. We are familiar with Wick’s theorem for bosons which states that the vacuum matrix element of a product of free fields is the sum of all possible products of two-point functions, with each product occuring with factor one. For fermions the factors are plus or minus one, depending on the parity of the permutation between the order in the vacuum matrix element and the order in the product of two point functions. In Wick’s theorem for quons the corresponding factors are $`q`$ raised to the inversion number of the permutation between the order in the vacuum matrix element and the order in the product of two point functions. This result reduces to the usual Wick’s theorem when $`q\pm 1`$. The inversion number can be found conveniently by drawing lines above the vacuum matrix element to indicate the pairs that are contracted into two-point functions. The minumum number of crossings of these lines is the inversion number. Since both cluster decomposition and the $`CPT`$ theorem for vacuum matrix elements of free fields depend on the properties of two-point functions, these theorems hold for quon fields. Note that quon fields, which clearly violate the spin-statistics theorem, obey the $`CPT`$ theorem; this emphasises the point made by R. Jost that the $`CPT`$ theorem requires only very weak assumptions.
If all the usual properties of relativistic field theory hold, then the spin-statistics theorem holds; thus some property must fail for quons. The property that does not hold is locality in the sense of the commutativity of observables at spacelike separation. Jost showed that if locality holds in an open spacelike region then analyticity arguments prove that it holds everywhere outside the lightcone. This result does not hold if the violation of locality decreases–say–exponentially away from the lightcone. The experimental bounds on such a violation are not clear. Note that the nonrelativistic form of locality
$$[\rho (𝐱),\psi ^{}(𝐲)]_{}=\delta (𝐱𝐲)\psi ^{}(𝐲),$$
(39)
where $`\rho `$ is the charge density, does hold.
### 5.4 Conservation of statistics rules for quon theory
For the energies of systems that are widely (spacelike) separated to be additive, all terms in the Hamiltonian must be effective bose operators in the sense that
$$[H(x),\varphi (y)]_{}0,\mathrm{as}xy\mathrm{}\mathrm{spacelike}$$
(40)
for all fields $`\varphi `$. This condition imposes “conservation of statistics” rules, the simplest of which is that only an even number of fermi fields can appear in any term of the Hamiltonian. For parafields, which have local observables, Messiah and I showed that parafields must occur in even degree, except that, for $`p`$ odd, $`p`$ parafields can occur. I defined paragrassmann and quongrassmann numbers which must be used in coupling para and quon operators to external sources and showed that for external parasources there are analogous restrictions. Using these results, R.C. Hilborn and I gave an heuristic argument relating the $`q`$ parameter for electrons to that for photons. The result,
$$q_e^2=q_\gamma ,$$
(41)
allows the very accurate bound from the Ramberg-Snow experiment to be carried over to photons with comparable accuracy. Similar arguments work for any particles that are coupled to electrons through any chain of reactions.
### 5.5 Bound states of quons
The classical result about bound states of bosons and fermion due to E.P. Wigner and to P. Ehrenfest and J.R. Oppenheimer states that a bound state of bosons and fermions is a boson unless it has an odd number of fermions, in which case it is a fermion. Hilborn and I showed that this result generalizes for quons. A bound state of $`n`$ identical quons with parameter $`q_{constituent}`$ has parameter $`q_{bound}=q_{constituent}^{n^2}`$. This implies that if $`q_{nucleus}`$ is bounded within $`ϵ`$ of $`\pm 1`$ for a nucleus with $`A`$ nucleons, then $`q_{nucleon}`$ is bounded within $`ϵ/A^2`$. Thus the bound on the nucleons is stronger than the bound on the nucleus. Analogously the bound on quarks is improved by $`1/9`$ over the bound on nucleons. Michael Berry pointed out that in the context of the quon theory this result on $`q`$ for bound states implies that either the layers of compositeness stop or all particles are bosons or fermions.
This result reduces to the usual one for bosons and fermions since $`q`$ and $`q^2`$ are even or odd together.
## 6 Summary and open questions
Like any other physical property, statistics should be subjected to high-precision experimental tests. In order to interpret such tests we need a theory in which statistics can be violated and a parameter that gives a quantitative measure of the validity of statistics. A theory that allows violations of statistics cannot have all the properties we might like. So far quons are the best theory that allows small violations.
In summary the positive properties of quons as a field theory are (a) norms are positive, (b) a simple modification of Wick’s theorem holds, (c) cluster decomposition theorems hold, (d) the $`CPT`$ theorem holds, and (e) free fields can have relativistic kinematics. The negative properties are (a) spacelike commutativity of observables fails, and because of this (b) interacting relativistic field theory is in doubt. I do not have a concrete suggestion for the possible origin of small violations of the exclusion principle. One could turn this issue around and observe that the constraints of bose and fermi statistics are grafted onto the general structure of quantum theory in an ad hoc way and ask why these constraints are realized in nature. Study of the situation in which these constraints are violated may shed light on why they hold for the known particles. In any case a fundamental issue such as statistics should be subjected to experimental tests and to theoretical study, just as is being done for Lorentz and $`CPT`$ invariance.
What we lack is an “external” motivation for violation of statistics—that is a connection of violations with some other physical property. We also don’t have any insight into the level at which we can expect violations if they do occur. Possible external motivations for violation of statistics include (a) violation of $`CPT`$, (b) violation of locality, (c) violation of Lorentz invariance, (d) extra space dimensions, (e) discrete space and/or time and (f) noncommutative spacetime. Of these, (a) seems unlikely because the quon theory which obeys $`CPT`$ allows violations, (b) seems likely because if locality is satisfied we can prove the spin-statistics connection and there will be no violations, (c), (d), (e) and (f) seem possible.
At the conference the question was raised whether the stability of matter which depends on the exclusion principle might set stringent bounds on possible violations of statistics. This certainly deserves careful study.
Hopefully either violations will be found experimentally or our theoretical efforts will lead to understanding of why only bose and fermi statistics occur in Nature.
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# Three-party entanglement from positronium
## I Introduction
Entanglement or quantum correlations between many space-separated subsystems has been recognized as one of the most intrinsic properties of quantum mechanics and provides the basis for many genuine applications of quantum information theory. It is, then, quite natural to look for physical situations in which quantum entangled states are obtained. Most of the theoretical and experimental effort has so far been devoted to unveil physical realizations of quantum states describing two quantum correlated subsystems. The search for physical systems displaying clean three-party entanglement is not simple. In this paper, we shall analyze decays of particles as a natural scenario for fulfilling such a goal. More precisely, we shall show that the decay of ortho-positronium into three photons corresponds to a highly entangled state. Let us now review what entanglement can be used for and why it is interesting to look for quantum correlation between more than two particles.
In 1935 Einstein, Podolsky and Rosen , starting from three reasonable assumptions of locality, reality and completeness that every physical theory must satisfy, argued that quantum mechanics (QM) is an incomplete theory. They did not question quantum mechanics predictions but rather quantum mechanics interpretation . Their argument was based on some inconsistencies between quantum mechanics and their local-realistic premises (LR) which appear for quantum states of bipartite systems, $`|\psi _{d_1}_{d_2}`$. It was in 1964 when Bell showed that any theory compatible with LR assumptions can not reproduce some of the statistical predictions of QM, using a gedankenexperiment proposed in with two quantum correlated spin-$`\frac{1}{2}`$ particles in the singlet state
$$|s=\frac{1}{\sqrt{2}}\left(|01|10\right).$$
(1)
In his derivation, as it is well-known, quantum correlations or entanglement have a crucial role. Actually, the singlet state is known to be the maximally entangled state between two particles. The conflict between LR and QM arises since the latter violates some experimentally verifiable inequalities, called Bell inequalities, that any theory according to the local-realistic assumptions ought to satisfy. It is then possible to design real experiments testing QM against LR (for a detailed discussion see ). Correlations of linear polarizations of pair of photons were measured in 1982 showing strong agreement with quantum mechanichs predictions and violating Bell inequalities . Nowadays, Bell inequalities have been tested thoroughly in favor of QM .
More recently, it has been pointed out that some predictions for quantum systems having quantum correlations between more than two particles give a much stronger conflict between LR and QM than any entangled state of two particles. The maximally entangled state between three spin-$`\frac{1}{2}`$ particles, the so-called GHZ (Greenberger, Horne and Zeilinger) state
$$|\mathrm{GHZ}=\frac{1}{\sqrt{2}}(|000+|111),$$
(2)
shows some perfect correlations incompatible with any LR model (see and also for more details). It is then of obvious relevance to obtain these GHZ-like correlations. Producing experimentally a GHZ state has turned out to be a real challenge yet a controlled instance has been produced in a quantum optics experiment .
Entanglement is then important for our basic understanding of quantum mechanics. Recent developments on quantum information have furthermore shown that it is also a powerful resource for quantum information applications. For instance, teleportation uses entanglement in order to obtain surprising results which are impossible in a classical context. A lot of work has been performed trying to know how entanglement can be quantified and manipulated. Our aim in this paper consists on looking for GHZ-like correlations, which are truly three-party pure state entanglement, in the decay of ortho-positronium to three photons. The choice of this physical system has been motivated mainly by several reasons. First, decay of particles seems a very natural source of entangled particles. Indeed, positronium decay to two photons was one of the physical systems proposed long time ago as a source of two entangled space-separated particles . On a different line of thought, some experiments for testing quantum mechanics have been recently proposed using correlated neutral kaons coming from the decay of a $`\varphi `$-meson . In the case of positronium, three entangled photons are obtained in the final state, so it offers the opportunity of analyzing a quantum state showing three-party correlations similar to other experiments in quantum optics.
The structure of the paper goes as follows. We first review the quantum states emerging in both para- and ortho-positronium decays. Then, we focus on their entanglement properties and proceed to a modern analysis of the three-photon decay state of ortho-positronium. Using techniques developed in the context of quantum information theory, we show that this state allows in principle for an experimental test of QM finer than the ones based on the use of the singlet state. We have tried to make the paper self-contained and easy to read for both particle physicists and quantum information physicists. The first ones can found a translation of some of the quantum information ideas to a well-known situation, that is, the positronium decay to photons, while the second ones can see an application of the very recent techniques obtained for three-party entangled states, which allow to design a QM vs LR test for a three-particle system in a situation different from the GHZ state.
## II Positronium decays
### A Positronium properties
Let us start reminding some basic facts about positronium. Positronium corresponds to a $`e^+e^{}`$ bound state. These two spin-$`\frac{1}{2}`$ particles can form a state with total spin equal to zero, para-positronium (p-Ps), or equal to one, ortho-positronium (o-Ps). Depending on the value of its angular momentum, it can decay to an even or an odd number of photons as we shall see shortly.
Positronium binding energy comes from the Coulomb attraction between the electron and the positron. In the non-relativistic limit, its wave function is
$$\mathrm{\Psi }(r)=\frac{1}{\sqrt{\pi a^3}}e^{\frac{r}{a}}=\frac{d^3p}{(2\pi )^{3/2}}e^{i\stackrel{}{p}\stackrel{}{r}}\stackrel{~}{\mathrm{\Psi }}(\stackrel{}{p})=\frac{d^3p}{(2\pi )^{3/2}}e^{i\stackrel{}{p}\stackrel{}{r}}\frac{\sqrt{8a^3}}{\pi (1+a^2p^2)^2},$$
(3)
where $`a=\frac{2}{m\alpha }`$, i.e. twice the Bohr radius of atomic hydrogen, and $`m`$ is the electron mass. Note that the wave function takes significant values only for three-momenta such that $`p\frac{1}{a}m`$, which is consistent with the fact that the system is essentially non-relativistic.
The parity and charge conjugation operators are equal to
$$U_P=(1)^{L+1}U_C=(1)^{L+S},$$
(4)
where $`L`$ and $`S`$ are the orbital and spin angular momentum. Positronium states are then classified according to these quantum numbers so that the ground states are $`{}_{}{}^{1}S_{0}^{}`$, with $`J^{PC}=0^+`$, for the p-Ps and $`{}_{}{}^{3}S_{1}^{}+^3D_1`$, having $`J^{PC}=1^{}`$, for the o-Ps.
Positronium is an unstable bound state that can decay to photons. Since a $`n`$-photon state transforms as $`U_C|n\gamma =(1)^n|n\gamma `$ under charge conjugation, which is an exact discrete symmetry for any QED process such as the decay of positronium, we have that the ground state of p-Ps (o-Ps) decays to an even (odd) number of photons . The analysis of the decay of positronium to photons can be found in a standard QED textbook . Para-positronium lifetime is about 0.125 ns, while for the case of ortho-positronium the lifetime is equal to approximately 0.14 $`\mu `$s .
The computation of positronium decays is greatly simplified due to the following argument. The scale which controls the structure of positronium is of the order of $`|\stackrel{}{p}|\alpha m`$. On the other hand, the scale for postrinomium annihilation is of the order of $`m`$. Therefore, it is easy to prove that positronium decays are only sensitive to the value of the wave function at the origin. As a consequence, it is possible to factor out the value of the wave function from the tree-level QED final state computation . A simple computation of Feymann diagrams will be enough to write the precise structure of momenta and polarizations which describe the positronium decays. Furthermore, only tree-level amplitudes need to be computed since higher corrections are suppressed by one power of $`\alpha `$. Let us now proceed to analyze the decays of p-Ps and o-Ps in turn.
### B Para-positronium decay
Para-positronium ground state decays into two photons. Because of the argument mentioned above, the determination of the two-photon state coming from the p-Ps decay is simply given by the lowest order Feynmann diagram of $`e^+e^{}\gamma \gamma `$. Since positronium is a non-relativistic particle to a very good approximation, the three-momenta of $`e^+`$ and $`e^{}`$ are taken equal to zero, and the corresponding spinors are replaced by a two-component spin. This implies that the tree-level calculation of the annihilation of p-Ps into two photons is equal to, up to constants,
$$(e^+e^{}\gamma \gamma )\chi _+^cM_2\chi _{},$$
(5)
where (see for more details) $`\chi _\pm `$ is the two-component spinor describing the fermions, $`\chi ^c\chi ^Ti\sigma _2`$, and $`M_2`$ gives
$$M_2=\underset{perm}{}(\stackrel{}{ϵ}_1^{}\times \stackrel{}{ϵ}_2^{})\widehat{k}I_{2\times 2}A(\widehat{k}_1,\lambda _1;\widehat{k}_2,\lambda _2)I_{2\times 2},$$
(6)
where $`\stackrel{}{ϵ}_i^{}\stackrel{}{ϵ}^{}(\widehat{k}_i,\lambda _i)`$ stands for the circular polarization vector associated to the outgoing photon $`i`$ and $`I_{2\times 2}`$ is the $`2\times 2`$ identity matrix. More precisely, for a photon having the three-momentum vector $`\stackrel{}{k}=|\stackrel{}{k}|\widehat{k}=|\stackrel{}{k}|(\mathrm{sin}\theta \mathrm{cos}\varphi ,\mathrm{sin}\theta \mathrm{sin}\varphi ,\mathrm{cos}\theta )`$, the polarization vectors can be chosen
$$\stackrel{}{ϵ}(\widehat{k},\lambda )=\frac{\lambda }{\sqrt{2}}(\mathrm{cos}\theta \mathrm{cos}\varphi i\lambda \mathrm{sin}\varphi ,\mathrm{cos}\theta \mathrm{sin}\varphi +i\lambda \mathrm{cos}\varphi ,\mathrm{sin}\theta ),$$
(7)
where $`\lambda =\pm 1`$ and they obey
$`\widehat{k}\stackrel{}{ϵ}(\widehat{k},\lambda )=0\widehat{k}\times \stackrel{}{ϵ}(\widehat{k},\lambda )=i\lambda \stackrel{}{ϵ}(\widehat{k},\lambda )`$ (8)
$`\stackrel{}{ϵ}(\widehat{k}_i,\lambda _i)\stackrel{}{ϵ}(\widehat{k}_j,\lambda _j)={\displaystyle \frac{1}{2}}\left(1\lambda _i\lambda _j\widehat{k}_i\widehat{k}_j\right).`$ (9)
From the expressions of the polarizaton vectors and the three-momentum and energy conservation, it follows that the scalar term $`A`$ is
$$A(\widehat{k},\lambda _1;\widehat{k},\lambda _2)=\frac{i}{2}(\lambda _1+\lambda _2),$$
(10)
and it verifies
$`A(\widehat{k},+1;\widehat{k},+1)=A(\widehat{k},1;\widehat{k},1)`$ (11)
$`A(\widehat{k},+1;\widehat{k},1)=A(\widehat{k},+1;\widehat{k},1)=0.`$ (12)
The two fermions in the para-positonium ground state are in the singlet state, $`|S`$$`=`$$`0,S_z`$$`=`$$`0=\frac{1}{\sqrt{2}}(|\frac{1}{2},\frac{1}{2}|\frac{1}{2},\frac{1}{2})`$, and then, using the previous relations for $`A`$ and (5), the two-photon state resulting of the p-Ps desintegration is
$$|\psi _p=\frac{1}{\sqrt{2}}(|++|).$$
(13)
The two-photon state resulting from p-Ps decay is thus equivalent to a maximally entangled state of two spin-$`\frac{1}{2}`$ particles. This is a well-known result and was, actually, one of the physical system first proposed as a source of particles having the quantum correlations needed to test QM vs LR .
### C Ortho-positronium decay
The ground state of ortho-positronium has $`J^{PC}=1^{}`$ and, due to the fact that charge conjugation is conserved, decays to three photons. Repeating the treatment performed for the p-Ps annihilation, the determination of the three-photon state resulting from the o-Ps decay requires the simple calculation of the tree-level Feynmann diagrams corresponding to $`e^+e^{}\gamma \gamma \gamma `$. Its tree-level computation gives, up to constants,
$$(e^+e^{}\gamma \gamma \gamma )\chi _+^cM_3\chi _{},$$
(14)
and the $`2\times 2`$ matrix $`M_3`$ is equal to
$$M_3=\underset{cyclicperm.}{}\left(\left(\stackrel{}{ϵ}_2^{}\stackrel{}{ϵ}_3^{}\stackrel{}{\delta }_2\stackrel{}{\delta }_3\right)\stackrel{}{ϵ}_1^{}+\left(\stackrel{}{ϵ}_2^{}\stackrel{}{\delta }_3+\stackrel{}{ϵ}_3^{}\stackrel{}{\delta }_2\right)\stackrel{}{\delta }_1\right)\stackrel{}{\sigma },$$
(15)
where
$$\stackrel{}{\delta }_i=\stackrel{}{k}_i\times \stackrel{}{ϵ}_i^{}.$$
(16)
Using (8) we can rewrite $`M_3`$ in the following way
$$M_3\stackrel{}{\sigma }\stackrel{}{V}(\widehat{k}_1,\lambda _1;\widehat{k}_2,\lambda _2;\widehat{k}_3,\lambda _3),$$
(17)
where
$`\stackrel{}{V}=`$ $`((\lambda _1\lambda _2)(\lambda _2+\lambda _3)\stackrel{}{ϵ}^{}(\widehat{k}_1,\lambda _1)(\stackrel{}{ϵ}^{}(\widehat{k}_2,\lambda _2)\stackrel{}{ϵ}^{}(\widehat{k}_3,\lambda _3))`$ (20)
$`+(\lambda _2\lambda _3)(\lambda _3+\lambda _1)\stackrel{}{ϵ}^{}(\widehat{k}_2,\lambda _2)\left(\stackrel{}{ϵ}^{}(\widehat{k}_3,\lambda _3)\stackrel{}{ϵ}^{}(\widehat{k}_1,\lambda _1)\right)`$
$`+(\lambda _3\lambda _1)(\lambda _1+\lambda _2)\stackrel{}{ϵ}^{}(\widehat{k}_3,\lambda _3)(\stackrel{}{ϵ}^{}(\widehat{k}_1,\lambda _1)\stackrel{}{ϵ}^{}(\widehat{k}_2,\lambda _2))).`$
Notice that the helicity coefficient $`(\lambda _i\lambda _j)(\lambda _j+\lambda _k)`$ for the cyclic permutations of $`ijk`$ explicitly enforces the vanishing of the $`(+++)`$ and $`()`$ polarizations,
$$\stackrel{}{V}(\widehat{k}_1,+;\widehat{k}_2,+;\widehat{k}_3,+)=\stackrel{}{V}(\widehat{k}_1,;\widehat{k}_2,;\widehat{k}_3,)=0.$$
(21)
On the other hand, the rest of structures are different from zero
$`\stackrel{}{V}(\widehat{k}_1,;\widehat{k}_2,+;\widehat{k}_3,+)=2\stackrel{}{ϵ}^{}(\widehat{k}_1,)(1\widehat{k}_2\widehat{k}_3)`$ (22)
$`\stackrel{}{V}(\widehat{k}_1,+;\widehat{k}_2,;\widehat{k}_3,)=2\stackrel{}{ϵ}^{}(\widehat{k}_1,+)(1\widehat{k}_2\widehat{k}_3),`$ (23)
and similar expressions for the other cyclic terms.
The original $`e^+e^{}`$ in the ortho-positronium could be in any of the three triplet states. It can be shown, using (14) and (17), that when the initial positronium state is $`|S`$$`=`$$`1,S_z`$$`=`$$`1=|\frac{1}{2},\frac{1}{2}`$, the decay amplitude is proportional to $`V_1+iV_2`$, while the same argument gives $`V_1+iV_2`$ for $`|S`$$`=`$$`1,S_z`$$`=`$$`1=|\frac{1}{2},\frac{1}{2}`$ and $`\sqrt{2}V_3`$ for $`|S`$$`=`$$`1,S_z`$$`=`$$`0=\frac{1}{\sqrt{2}}(|\frac{1}{2},\frac{1}{2}+|\frac{1}{2},\frac{1}{2})`$. Now, considering the explicit expressions of the polarization vectors (7), with $`\theta =\frac{\pi }{2}`$ without loss of generality, and (22), it is easy to see that the three-photon state coming from the o-Ps decay is, up to normalization,
$`|\psi _0(\widehat{k}_1,\widehat{k}_2,\widehat{k}_3)=`$ $`(1\widehat{k}_1\widehat{k}_2)(|+++|+)`$ (24)
$`+`$ $`(1\widehat{k}_1\widehat{k}_3)(|+++|+)`$ (25)
$`+`$ $`(1\widehat{k}_2\widehat{k}_3)(|+++|+),`$ (26)
when the third component of the ortho-positronium spin, $`S_z`$, is equal to zero, and
$`|\psi _1(\widehat{k}_1,\widehat{k}_2,\widehat{k}_3)=`$ $`(1\widehat{k}_1\widehat{k}_2)(|++|+)`$ (27)
$`+`$ $`(1\widehat{k}_1\widehat{k}_3)(|++|+)`$ (28)
$`+`$ $`(1\widehat{k}_2\widehat{k}_3)(|++|+),`$ (29)
when $`S_z=\pm 1`$.
The final state of the o-Ps decay is, thus, an entangled state of three photons, whose quantum correlations depend on the angles among the momenta of the outgoing three photons. For the rest of the paper we will consider the first family of states ($`S_z=0`$) although equivalent conclusions are valid for the second one. In the next sections we will analyze the entanglement properties of the states $`|\psi _0(\widehat{k}_1,\widehat{k}_2,\widehat{k}_3)`$, using some of the quantum information techniques and comparing them to the well-known cases of the singlet and GHZ state.
## III Entanglement properties
The quantum correlations of the three-photon entangled state obtained from the o-Ps annihilation depend on the position of the photon detectors, i.e. on the photon directions we are going to measure. Our next aim will be to choose from the family of states given by (24) the one that, in some sense, has the maximum amount of GHZ-like correlations. In order to do this, we first need to introduce some recent results on the study of three-party entanglement.
The set of states $`|\psi _0(\widehat{k}_1,\widehat{k}_2,\widehat{k}_3)`$ form a six-parameter dependent family in the Hilbert space $`_2_2_2`$, so that each of its components is equivalent to a state describing three spin-$`\frac{1}{2}`$ particles or three qubits (a qubit, or quantum bit, is the quantum version of the classical bit and corresponds to a spin-$`\frac{1}{2}`$ particle). Two pure states belonging to a generic composite system $`_d^N`$, i.e. $`N`$ parties each having a $`d`$-dimensional Hilbert space, are equivalent as far as their entanglement properties go when they can be transformed one into another by local unitary transformations. This argument gives a lower bound for the entanglement parameters a generic state $`|\varphi _2^N`$ depends on. Since the number of real parameters for describing it is $`2^{N+1}`$, and the action of an element of the group of local unitary transformations $`U(2)^N`$ is equivalent to the action of $`U(1)\times SU(2)^N`$, which depends on $`3N+1`$ real parameters, the number of entanglement parameters is bounded by $`2^{N+1}(3N+1)`$. For our case this counting of entanglement parameters gives six, since we have $`N=3`$, and it can be proved that this is indeed the number of nonlocal parameters describing a state in $`_2_2_2`$ .
The above arguments imply that six independent quantities invariant under the action of the group of local unitary transformations will be enough, up to some discrete symmetry, to describe the entanglement properties of any three-qubit pure state. Given a generic state $`|\varphi _2^3`$
$$|\varphi =\underset{i,j,k}{}t_{ijk}|ijki,j,k=1,2,$$
(30)
where $`|i,|j,|k`$ are the elements of a basis in each subsystem, A, B and C, the application of three local unitary transformations $`U^A`$, $`U^B`$ and $`U^C`$ transforms the coefficients $`t_{ijk}`$ into
$$t_{ijk}^{}=U_{i\alpha }^AU_{j\beta }^BU_{k\gamma }^Ct_{\alpha \beta \gamma }.$$
(31)
From this expression it is not difficult to build polynomial combinations of the coefficient $`t_{ijk}`$ which are invariant under local unitary transformations . These quantities are good candidates for being an entanglement parameter. For example, one of these invariants is
$$t_{i_1j_1k_1}t_{i_1j_2k_2}^{}t_{i_2j_2k_2}t_{i_2j_1k_1}^{}=\mathrm{tr}(\rho _A^2),$$
(32)
where $`\rho _A=\mathrm{tr}_{BC}(|\varphi \varphi |)`$ is the density matrix describing the local quantum state of A (and the same happens for B and C). In the six linearly independent polynomial invariants of minor degree were found (a trivial one is the norm) and a slightly modified version of these quantities was also proposed in . In the rest of the paper we will not consider the norm, so the space of entanglement parameters of the normalized states belonging to $`_2_2_2`$ has dimension equal to five.
A particularly relevant polynomial invariant is the so-called tangle, $`\tau `$, introduced in . There is strong evidence that somehow it is a measure of the amount of “GHZ-ness” of a state . It corresponds to the modulus of the hyperdeterminant of the hypermatrix given by the coefficients $`t_{ijk}`$ , which from (30) corresponds to
$$\tau (|\varphi )=|\mathrm{Hdet}(t_{ijk})|=\left|ϵ_{i_1i_2}ϵ_{i_3i_4}ϵ_{j_1j_2}ϵ_{j_3j_4}ϵ_{k_1k_3}ϵ_{k_2k_4}t_{i_1j_1k_1}t_{i_2j_2k_2}t_{i_3j_3k_3}t_{i_4j_4k_4}\right|,$$
(33)
where $`ϵ_{00}=ϵ_{11}=0`$ and $`ϵ_{01}=ϵ_{10}=1`$. This quantity can be shown to be symmetric under permutation of the indices $`i,j,k`$.
Because of the interpretation of the tangle as a measure of the GHZ-like correlations, we will choose the position of the photon detectors, from the set of states (24), the ones that are associated to a maximum tangle. In the figure 1 it is shown the variation of the tangle with the position of the detectors. It is not difficult to see that the state of (24) with maximum tangle corresponds to the case $`\widehat{k}_1\widehat{k}_2=\widehat{k}_1\widehat{k}_3=\widehat{k}_2\widehat{k}_3=\frac{1}{2}`$, i.e. the most symmetric configuration, that we shall call “Mercedes-star” geometry. The normalized state obtained from (24) for this geometry is
$$|\psi =\frac{1}{\sqrt{6}}(|+++|++|+++|++|+++|+).$$
(34)
Note that the GHZ state has tangle equal to $`\frac{1}{4}`$, while the value of the tangle of (34) is lower,
$$\tau (|\psi )=\frac{1}{12}.$$
(35)
It is arguable that a “Mercedes-star” geometry was naturally expected to produce a maximum tangle state. Indeed, GHZ-like quantum correlations do not singularize any particular qubit.
Let us also mention that the state we have singled out has some nice properties from the point of view of group theory. It does correspond to the sum of two of the elements of the coupled basis resulting from the tensor product of three spin-$`\frac{1}{2}`$ particles, $`\frac{1}{2}\frac{1}{2}\frac{1}{2}`$,
$$|\psi =\frac{1}{\sqrt{2}}\left(|\frac{3}{2},+\frac{1}{2}+|\frac{3}{2},\frac{1}{2}\right),$$
(36)
where
$`|{\displaystyle \frac{3}{2}},+{\displaystyle \frac{1}{2}}={\displaystyle \frac{1}{\sqrt{3}}}(|+++|+++|++)`$ (37)
$`|{\displaystyle \frac{3}{2}},{\displaystyle \frac{1}{2}}={\displaystyle \frac{1}{\sqrt{3}}}(|++|++|+).`$ (38)
The quantum correlations of (34) will be now analyzed.
## IV Useful decompositions
In this section, the state (34) will be rewritten in some different forms that will help us to understand better its nonlocal properties. First, let us mention that for any generic three-qubit pure state and by performing change of local bases, it is possible to make zero at least three of the coefficients $`t_{ijk}`$ of (30) . A simple counting of parameters shows that this is in fact the expected number of zeros. This means that by a right choice of the local bases, any state can be written with the minimum number of coefficients $`t_{ijk}`$, i.e. we are left with all the non-local features of the state, having removed all the “superfluous” information due to local unitary tranformations. For the case of the state (34) it is easy to prove that it can be expressed as
$$|\psi =\frac{1}{2\sqrt{3}}\left(|001+|010+|100\right)+\frac{\sqrt{3}}{2}|111,$$
(39)
which is the minimum decomposition in terms of product states built from local bases (four of the coefficients $`t_{ijk}`$ are made equal to zero).
An alternative decomposition, that will prove to be fruitful for the rest of the paper, consists of writing the state as a sum of two product states. This decomposition is somewhat reminiscent of the form of the GHZ state, which is a sum of just two product states, and is only possible when the tangle is different from zero as it happens for our state (see 35). The state then can be written as
$`|\psi `$ $`={\displaystyle \frac{2}{3}}\left(\left(\begin{array}{c}1\\ 0\end{array}\right)\left(\begin{array}{c}1\\ 0\end{array}\right)\left(\begin{array}{c}1\\ 0\end{array}\right)+\left(\begin{array}{c}\frac{1}{2}\\ \frac{\sqrt{3}}{2}\end{array}\right)\left(\begin{array}{c}\frac{1}{2}\\ \frac{\sqrt{3}}{2}\end{array}\right)\left(\begin{array}{c}\frac{1}{2}\\ \frac{\sqrt{3}}{2}\end{array}\right)\right)`$ (41)
$`\alpha (|000+|aaa),`$
where $`|0\left(\begin{array}{c}1\\ 0\end{array}\right)`$ and $`a\left(\begin{array}{c}\frac{1}{2}\\ \frac{\sqrt{3}}{2}\end{array}\right)`$. We omit the details for the explicit computation of this expression since they can be found in . It is worth noticing that o-Ps decay is hereby identified to belonging to an interesting type of states already classified in quantum information theory .
The above decomposition allows for an alternative interpretation of the initial state as an equally weighted sum of two symmetric product states. Note that the Bloch vector, $`\widehat{n}=(\mathrm{sin}\theta \mathrm{cos}\varphi ,\mathrm{sin}\theta \mathrm{sin}\varphi ,\mathrm{cos}\theta )`$, representing the first local spinor appearing in (41) is pointing to the $`z`$ axis, i.e. $`\widehat{n}_1=(0,0,1)`$, while the second is located in the $`XZ`$ plane with an angle of 120 with the $`z`$ axis, i.e. $`\widehat{n}_2=(\frac{\sqrt{3}}{2},0,\frac{1}{2})`$. By performing a new unitary transformation, (41) can be written as
$$|\psi =\frac{2}{3}\left(\left(\begin{array}{c}c\\ s\end{array}\right)\left(\begin{array}{c}c\\ s\end{array}\right)\left(\begin{array}{c}c\\ s\end{array}\right)+\left(\begin{array}{c}s\\ c\end{array}\right)\left(\begin{array}{c}s\\ c\end{array}\right)\left(\begin{array}{c}s\\ c\end{array}\right)\right),$$
(42)
where $`c=\mathrm{cos}15^{},s=\mathrm{sin}15^{}`$. Now, the two Bloch vectors are in the $`XZ`$ plane, pointing to the $`\theta =30^{}`$ and $`\theta =150^{}`$ directions. The GHZ state corresponds to the particular case $`c=1`$ and $`s=0`$.
## V Quantum mechanics vs local realism
The quantum correlations present in some three-qubit pure states show, as it was mentioned in the introduction, a much stronger disagreement with the predictions of a local-realistic model than any two-qubit entangled state. In fact, contrary to the case of the singlet state, no LR model is able to reproduce all the perfect correlations predicted for the maximally entangled state of three qubits . The state (34) emerging from o-Ps decay is not a GHZ state, although it has been chosen the one with the maximum tangle in order to maximize GHZ-like correlations. In this section we will show how to use it for testing quantum mechanics against local-realistic models, and then we will compare its performance against existing tests for the maximally entangled states of two and three spin-$`\frac{1}{2}`$ particles. We start reviewing some of the consequences derived from the arguments proposed in .
### A QM vs LR conflict
Given a generic quantum state of a composite system shared by $`N`$ parties, there should be an alternative LR theory which reproduces all its statistical predictions. In this LR model, a state denoted by $`\lambda `$ will be assigned to the system specifying all its elements of physical reality. In particular, the result of a measurement depending on a set of parameters $`\{n\}`$ performed locally by one of the parties, say A, will be specified by a function $`a_\lambda (\{n\})`$. The same will happen for each of the space-separated parties and, since there is no causal influence among them, the result measured on A can not modify the measurement on B. For example, if the measurement is of the Stern-Gerlach type, the parameters labeling the measurement are given by a normalized vector $`\widehat{n}`$ and $`a_\lambda (\widehat{n})a`$ are the LR functions describing the outcome.
The LR model can be very general provided that some conditions must be satisfied. Consider a generic pure state belonging to $`_2_2_2`$ shared by three observers A, B and C, which are able to perform Stern-Gerlach measurements in any direction. Since the outcomes of a Stern-Gerlach measurement are only $`\pm 1`$, it is easy to check that for any pair of measurements on each subsystem, described by the LR functions $`a`$ and $`a^{}`$, $`b`$ and $`b^{}`$, $`c`$ and $`c^{}`$, and for all their possible values, it is always verified
$$a^{}bc+ab^{}c+abc^{}a^{}b^{}c^{}=\pm 2.$$
(43)
It follows from this relation that
$$2a^{}bc+ab^{}c+abc^{}a^{}b^{}c^{}2.$$
(44)
This constraint is known as Mermin inequality and has to be satisfied by any LR model describing three space-separated systems.
Let us now take the GHZ state (2). It is quite simple to see that if the observables $`a`$ and $`a^{}`$ are equal to $`\sigma _y`$ and $`\sigma _x`$ (the same for parties B and C), the value of (44) is $`4`$, so an experimental condition is found that allows to test quantum mechanics against local realism. Note that this is the maximal violation of inequality (44). Moreover, the GHZ state also satisfies that $`a^{}bc=ab^{}c=abc^{}=a^{}b^{}c^{}=1`$ and no LR model is able to take into account this perfect correlation result because of (43) . This is a new feature that does not appear for the case of a two maximally entangled state of two spin-$`\frac{1}{2}`$ particles. In this sense it is often said that a most dramatic contrast between QM and LR emerges for entanglement between three subsystems.
Let us go back to the state given by the ortho-positronium decay (34). Our aim is to design an experimental situation where a conflict between QM and LR appears, so we will look for the observables that give a maximal violation of (44). Such observables will extremize that expression. Using the decomposition (42), the expectation value of three local observables is
$`abc`$ $`=`$ $`\psi |(\widehat{n}_a\stackrel{}{\sigma })(\widehat{n}_b\stackrel{}{\sigma })(\widehat{n}_c\stackrel{}{\sigma })|\psi `$ (45)
$`=`$ $`{\displaystyle \frac{4}{9}}({\displaystyle \underset{i=a,b,c}{}}(\stackrel{~}{c}\mathrm{cos}\theta _i+\stackrel{~}{s}\mathrm{sin}\theta _i\mathrm{cos}\varphi _i)+{\displaystyle \underset{i=a,b,c}{}}(\stackrel{~}{c}\mathrm{cos}\theta _i+\stackrel{~}{s}\mathrm{sin}\theta _i\mathrm{cos}\varphi _i)`$ (46)
$`+`$ $`{\displaystyle \underset{i=a,b,c}{}}\mathrm{sin}\theta _i(c^2e^{i\varphi _i}+s^2e^{i\varphi _i})+{\displaystyle \underset{i=a,b,c}{}}\mathrm{sin}\theta _i(c^2e^{i\varphi _i}+s^2e^{i\varphi _i})),`$ (47)
where $`\stackrel{~}{c}c^2s^2`$ and $`\stackrel{~}{s}2sc`$. Because of the symmetry of the state under permutation of parties, the Stern-Gerlach directions are taken satisfying $`\widehat{n}_a=\widehat{n}_b=\widehat{n}_c=(\mathrm{sin}\theta \mathrm{cos}\varphi ,\mathrm{sin}\theta \mathrm{sin}\varphi ,\mathrm{cos}\theta )`$ and $`\widehat{n}_a^{}=\widehat{n}_b^{}=\widehat{n}_c^{}=(\mathrm{sin}\theta ^{}\mathrm{cos}\varphi ^{},\mathrm{sin}\theta ^{}\mathrm{sin}\varphi ^{},\mathrm{cos}\theta ^{})`$. Substituting this expression in (44), we get the explicit function $`f(\theta ,\varphi ,\theta ^{},\varphi ^{})`$ to be extremized. For the case of the GHZ state described above, the extreme values were obtained using two observables with $`\theta =\theta ^{}=\frac{\pi }{2}`$, i.e. in the $`XY`$ plane. Since (42) is the GHZ-like decomposition of the initial state, we take $`\theta =\theta ^{}=\frac{\pi }{2}`$ and it is easy to check that in this case $`\frac{f}{\theta }|_{\theta =\theta ^{}=\frac{\pi }{2}}=\frac{f}{\theta ^{}}|_{\theta =\theta ^{}=\frac{\pi }{2}}=0,\varphi ,\varphi ^{}`$. Mantaining the parallelism with the GHZ case, it can be seen that all the partial derivatives vanish when it is also imposed $`\varphi =\frac{\pi }{2}`$ and $`\varphi ^{}=0`$. In our case the calculation of (44) gives $`3`$, so a conflict between local-realistic models and quantum mechanics again appears, and then the three-photon state coming from the ortho-positronium decay can be used, in principle, to test QM vs LR with the set of observables given by the normalized vectors
$$\widehat{n}_a=\widehat{n}_b=\widehat{n}_c=(0,1,0)\widehat{n}_a^{}=\widehat{n}_b^{}=\widehat{n}_c^{}=(1,0,0).$$
(48)
There is an alternative set of angles $`\varphi `$ and $`\varphi ^{}`$ that makes zero all the partial derivatives of $`f`$: the combination of local observables (44) is equal to $`3.046`$ for
$$\varphi ^{}=\mathrm{arctan}\left(\frac{\sqrt{17+27\sqrt{41}}}{10}\right)126^{}\varphi =\frac{1}{2}\mathrm{arctan}\left(\frac{2\sqrt{17+27\sqrt{41}}}{25}\right)24^{}.$$
(49)
This second set of parameters will be seen to produce in the end a weaker dismissal of LR.
Our next step will be to carry over the comparison of this QM vs LR test against the existent ones for the maximally entangled states of three and two spin-$`\frac{1}{2}`$ particles, i.e. the GHZ and singlet state. It is quite evident that the described test should be worse than the obtained for the GHZ state. It is less obvious how this new situation will compare with the singlet case.
### B Comparison with the maximally entangled states of two and three spin-$`\frac{1}{2}`$ particles
We will now estimate the “strength” of the QM vs LR test proposed above, being this “strength” measured by the number of trials needed to rule out local-realism at a given confidence level, as Peres did in . A reasoning anologous to the one given in will be done here for the state (34) and the observables (48).
Imagine a local-realistic physicist who does not believe in quantum mechanics. He assigns prior subjective probabilities to the validity of LR and QM, $`p_r`$ and $`p_q`$, expressing his personal belief. Take for instance $`\frac{p_r}{p_q}=100`$. His LR theory is not able to reproduce exactly all the QM statistical results of some quantum states. Consider the expectation value of some observable $`𝒪`$ with two outcomes $`\pm 1`$ such that $`𝒪=E_q`$ is predicted for some quantum state, while LR gives $`𝒪=E_rE_q`$. Since the value of the two possible outcomes are $`\pm 1`$, the probablity of having $`𝒪=+1`$ is $`q=\frac{1+E_q}{2}`$ for QM and $`r=\frac{1+E_r}{2}`$ for LR. An experimental test of the observable $`𝒪`$ now is performed $`n`$ times yielding $`m`$ times the result $`+1`$. The prior probabilities $`p_q`$ and $`p_r`$ are modified according to the Bayes theorem and their ratio has changed to
$$\frac{p_r^{}}{p_q^{}}=\frac{p_r}{p_q}\frac{p(m|_{LR})}{p(m|_{QM})},$$
(50)
where
$$p(m|_{LR})=\left(\begin{array}{c}n\\ m\end{array}\right)r^m(1r)^{nm},$$
(51)
is the LR probability of having $`m`$ times the outcome $`+1`$, and we have the same for $`p(m|_{QM})`$, being $`r`$ replaced by $`q`$. Following Peres , the confidence depressing factor is defined
$$D\frac{p(m|_{QM})}{p(m|_{LR})}=\left(\frac{q}{r}\right)^m\left(\frac{1q}{1r}\right)^{nm},$$
(52)
which accounts for the change in the ratio of the probabilities of the two theories, i.e. it reflects how the LR belief changes with the experimental results. Like in a game, our aim is to destroy as fast as we can the LR faith of our friend by choosing an adequate experimental situation. It can be said, for example, that he will give up when, for example, $`D=10^4`$. Since the world is quantum, $`m=qn`$, and the number of experimental tests needed to obtain $`D=10^4`$ is equal to
$$n_D(q,r)\frac{4}{q\mathrm{log}_{10}\left(\frac{q}{r}\right)+(1q)\mathrm{log}_{10}\left(\frac{1q}{1r}\right)}=\frac{4}{K(q,r)},$$
(53)
being $`K(q,r)`$ the information distance between the QM and LR binomial distribution for the outcome $`+1`$. The more separate the two probability distributions are, measured in terms of the information distance, the fewer the number of experiments $`n_D`$ is.
Let us come back to the three-party entangled state coming from the ortho-positronium decay (34) under the local measurements described by (48). As it has been shown above, a contradiction with any LR model appears for the combination of the observables given by the Mermin inequality. In our case quantum mechanics gives the following predictions
$$a^{}bc=ab^{}c=abc^{}=\frac{2}{3}a^{}b^{}c^{}=+1,$$
(54)
and this implies that $`q_1=prob(a^{}bc=+1)=prob(ab^{}c=+1)=prob(abc^{}=+1)=\frac{1}{6}`$ and $`q_2=prob(a^{}b^{}c^{}=+1)=1`$. This is the QM data that our LR friend has to reproduce as well as possible. Because of the symmetry of the state he will assign the same probability $`r_1`$ to the events $`a^{}bc=+1`$, $`ab^{}c=+1`$ and $`abc^{}=+1`$ and $`r_2`$ to $`a^{}b^{}c^{}=+1`$. However, his model has to satisfy the constraint given by (44), so the best he can do is to saturate the bound and then
$$3r_1=r_2\mathrm{\hspace{0.17em}0}r_1\frac{1}{3}.$$
(55)
Now, according to the probabilities $`r_1`$ and $`r_2`$ his LR model predicts, we choose the experimental test that minimizes (53), i.e. we consider the event $`a^{}bc=+1`$ ($`a^{}b^{}c^{}=+1`$) when $`n_D(q_1,r_1)<n_D(q_2,r_2)`$ ($`n_D(q_1,r_1)>n_D(q_2,r_2)`$), and the experimental results will destroy his LR belief after $`n_D(q_1,r_1)`$ ($`n_D(q_2,r_2)`$) trials. The best value our LR friend can assign to $`r_1`$ is the solution to
$$n_D(q_1,r_1)=n_D(q_2,r_2),$$
(56)
with the constraint (55), and this condition means that $`r_10.315`$ and $`n_D161`$ trials are needed to have a depressing factor equal to $`10^4`$. Repeating the same calculation for the observables giving by (49), the number of trials slightly increases, $`n_D166`$, despite of the fact that the violation of the inequality is greater than the obtained for (48).
In ref. the same reasoning was applied to the maximally entangled state of two and three spin-$`\frac{1}{2}`$ particles, showing that $`n_D200`$ in the first case, and $`n_D32`$ for the latter (see table I). Our result then implies that the three-photon entangled state produced in the ortho-positronium decay has, in some sense, more quantum correlations than any entangled state of two spin-$`\frac{1}{2}`$ particles.
### C Generalization of the results
It is easy to generalize some of the results obtained for the entangled state resulting from the o-Ps decay. As it has been mentioned, this state can be understood as an equally weighted sum of two symmetric product states, since it can be written as (42). The Bloch vectors of the two local states appearing in this decomposition form an angle of $`120^{}`$. It is clear that the conclusions seen above depend on the angle between these vectors, i.e. with their degree of non-orthogonality. The family of states to be analyzed can be parametrized in the following way
$$|\psi (\delta )=\alpha _\delta \left(\left(\begin{array}{c}c_\delta \\ s_\delta \end{array}\right)\left(\begin{array}{c}c_\delta \\ s_\delta \end{array}\right)\left(\begin{array}{c}c_\delta \\ s_\delta \end{array}\right)+\left(\begin{array}{c}s_\delta \\ c_\delta \end{array}\right)\left(\begin{array}{c}s_\delta \\ c_\delta \end{array}\right)\left(\begin{array}{c}s_\delta \\ c_\delta \end{array}\right)\right),$$
(57)
where $`\delta `$ is the angle between the two local Bloch vectors, $`c_\delta \mathrm{cos}\left(\frac{\pi \delta }{4}\right)`$ and $`s_\delta \mathrm{sin}\left(\frac{\pi \delta }{4}\right)`$ and $`\alpha _\delta `$ is a positive number given by the normalization of the state. An alternative parametrization of this family is, using (39) and defining $`\delta ^{}\frac{\delta }{4}`$,
$$|\psi (\delta )=2\alpha _\delta \left(\mathrm{sin}^2\delta ^{}\mathrm{cos}\delta ^{}\left(|001+|010+|100\right)+\mathrm{cos}^3\delta ^{}|111\right).$$
(58)
The expectation value of three local observables for this set of states follows trivially from (45). Using this expression it is easy to see that the combination of the expectation values of (44) has all the partial derivatives equal to zero for the set of observables given in (48) independently of $`\delta `$. For these observables, the dependence of expression (44) with the degree of orthogonality between the two product states is given in figure 2. There is no violation of the Mermin inequality for the case in which $`\delta 85^{}`$. In this situation one can always found a LR model able to reproduce the QM statistical prediction given by (44) and the observables (48). We can now repeat all the steps made in order to determine the number of trials needed to rule out local realism as a function of the angle $`\delta `$. In figure 3 we have summaryzed the results. We have shown only the cases where the number of trials is minor than two hundred, since this is the value obtained for the singlet. Note that the case $`\delta =120^{}`$, which corresponds to (34), is very close to the region where there is no improvement compared to the maximally entangled state of two qubits.
All these results can be understood in the following way: the smaller the angle between the two local states, $`\delta `$, the higher the overlap of the state $`|\psi (\delta )`$ with the product state having each local Bloch vector pointing in the direction of the $`x`$ axis, which corresponds to the state $`|111`$ in (58). This means that the quantum state we are handling is too close to a product state , and thus, no violation of the Mermin inequality can be observed.
## VI Concluding remarks
In this work we have analyzed the three-particle quantum correlations of a physical system given by the decay of the ortho-positronium into a three-photon pure state. After obtaining the state describing the polarization of the three photons (34), some of the recent techniques developed for the study of three-party entanglement have been applied. The particular case where the three photons emerge in a symmetric, Mercedes-star-like configuration, corresponds to the state with the maximum tangle. We have shown that this state allows a priori for a QM vs LR test which is stronger than any of the existing ones that use the singlet state. In this sense, ortho-positronium decays into a state which carries stronger quantum correlations than any entangled state of two spin-$`\frac{1}{2}`$ particles.
Bose symmetrization has played a somewhat negative role in reducing the amount GHZ-ness of the o-Ps decay state. Indeed, the natural GHZ combination $`|+++|+`$ emerging from the computation of Feynmann diagrams has been symmetrized due to the absence of photon tagging to our state $`|+++|+++|+++|++|++|+`$, inducing a loss of tangle. The quantum optics realization of the GHZ state does avoid symmetrization through a geometric tagging . It is, thus, reasonable to look for pure GHZ states in decays to distinct particles, so that tagging would be carried by other quantum numbers, as e.g. charge. It is, on the other hand, peculiar to note that symmetrization in the $`K^0\overline{K}^0`$ system is responsible for its entanglement ($`|++|+`$) .
Finally, let us briefly discuss the experimental requirements needed for testing quantum mechanics as it has been described in this paper. In order to do this, the circular polarizations of the three photons resulting from an ortho-positronium decay have to be measured. The positions of the three detectors are given by the “Mercedes-star” geometry and their clicks have to detect the coincidence of the three photons. The energy of these photons is of the order of 1 Mev. Polarization analyzers with a good efficiency would allow us to acquire statistical data showing quantum correlations which would violate the Mermin inequality discussed above. Unfortunately, as far as we know, no such analizers exist for this range of energies. A possible way-out might be to use Compton scattering to measure the photon polarizations . However, Compton effect just gives a statistical pattern depending on the photon and electron polarizations which is not a direct measurement of the polarizations. Further work is needed to modify our analysis of QM vs LR to accommodate for such indirect measurements.
## Acknowledgments
We acknowledge J. Bernabeu for suggesting positronium as a source of three entangled particles and reading carefully the paper. We also thank A. Czarnecki, D. W. Gidley, M. A. Skalsey and V. L. Telegdi for comments about the measurement of the photon polarizations in the ortho-positronium decay. We acknowledge financial support by CICYT project AEN 98-0431, CIRIT project 1998SGR-00026 and CEC project IST-1999-11053, A. A. by a grant from MEC (AP98). Financial support from the ESF is also acknowledged. This work was concluded during the 2000 session of the Benasque Center for Science, Spain.
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# Finite-Size Scaling in the transverse Ising Model on a Square Lattice
## I Introduction
Recent advances in computer technology have allowed the exact diagonalization of Ising-type quantum spin systems up to 36 sites in size. Schulz, Ziman and Poilblanc (1996), for example, studied the J1-J2 XXZ Heisenberg spin model on square lattices up to 6x6 sites. Our aim in this paper is to carry out an exact diagonalization study of the transverse Ising model on the square lattice, in order to estimate its critical parameters and study its finite-size scaling behaviour.
The transverse Ising model in (2+1)D is well-known to be the quantum Hamiltonian corresponding to the classical 3D Ising model (Suzuki 1976, Fradkin and Susskind 1978), and exhibits a quantum phase transition in the same universality class as the classical 3D Ising thermal transition. It was first studied by series expansion methods by Pfeuty and Elliott (1971), and there have been several further series expansion calculations since then, both ‘low-temperature’ (Marland 1981, Yanase et al 1976, Oitmaa et al 1991) and ‘high-temperature’ (Hamer and Irving 1984, Hamer and Guttmann 1989, He et al 1990). Exact finite-lattice calculations have also been carried out previously (Roomany and Wyld 1980, Hamer 1983, Henkel 1984, 1987) for square lattices of up to 5x5 sites, and similar calculations have also been done for the triangular lattice (Hamer and Johnson 1986, Henkel 1990, Price et al 1993). Here we extend these calculations for the first time to the 6x6 lattice, and use finite-size scaling theory to obtain improved estimates of the critical point and critical index $`\nu `$.
The finite-size scaling amplitudes at the critical point are also of interest. In (1+1)D, it is well-known that the theory of conformal invariance relates the scaling amplitudes to fundamental parameters of the underlying effective field theory at the critical point, such as the conformal anomaly and scaling indices. In higher dimensions, a similar scenario is known to hold at ‘first-order’ transitions, where a continuous symmetry is spontaneously broken, giving rise to Goldstone bosons (Hasenfratz and Leutwyler 1990): the finite-size scaling amplitudes are related to parameters of the Goldstone bosons such as the spin-wave stiffness and spin-wave velocity. Does something similar apply at second-order transitions in higher dimensions? Apart from a discussion by Cardy (1985), little has been done in this area. One peculiar result was obtained by Henkel (1986, 1987) and Weston (1990), and confirmed recently by Weigel and Janke (1999): the scaling amplitudes of the spin-spin and energy-energy correlation lengths on antiperiodic lattices have a universal ratio:
$$\frac{A_\sigma }{A_ϵ}=\frac{x_\sigma }{x_ϵ}$$
(1)
where the $`x_i`$ are the scaling indices in the respective sectors. This phenomenon appears to have no good theoretical explanation at the present time.
Our exact diagonalization methods are outlined briefly in Section 2, and the numerical results are presented in Section 3. The critical parameters so obtained are compared with other estimates in Section 4, and the critical amplitudes are discussed.
## II Method
The transverse Ising model on the square lattice has the Hamiltonian
$`H={\displaystyle \underset{i}{}}(1\sigma _3(i))x{\displaystyle \underset{<ij>}{}}\sigma _1(i)\sigma _1(j)h{\displaystyle \underset{i}{}}\sigma _1(i)`$ (2)
where the sum $`<ij>`$ runs over nearest neighbour pairs on the lattice, and the $`\sigma `$ matrices are the usual Pauli spin operators acting on a 2-state spin-variable at each site. The coupling x is analogous to an inverse ‘temperature’, and h represents an external ‘magnetic field’. We shall employ a representation in which the $`\sigma _3(i)`$ are diagonal. Periodic boundary conditions are assumed.
The unperturbed ground-state of the model at $`x=0`$ has all spins ‘up’, i.e. $`\sigma _3(i)=+1`$, all $`i`$. The interaction term will induce an admixture of states with ‘flipped’ spins. The Hilbert space of the model consists of two sectors, containing an odd and even number of flipped spins respectively.
Exact diagonalizations have been carried out for LxL lattices, L = 1,..,6. The methods employed are fairly standard, for the most part, and will not be described in detail here. First, a list of allowed basis states in the given sector was prepared, using the ‘sub-lattice coding’ technique of Lin (1990). This efficient technique produces a sorted list of states, requiring only one integer word of storage per state. Since only the zero-momentum states are considered here, the states were ‘symmetrized’: that is, all copies of a given state under translations, reflections and rotations were represented by a single state. Thus for the 6x6 lattice in the even sector, the total number of ‘unsymmetrized’ states is approximately $`2^{35}`$, whereas under symmetrization this is reduced by a factor of approximately 288, down to 119,539,680.
Next, the Hamiltonian matrix elements are generated, by applying the interaction operators of equation (2) to each initial state, symmetrizing the resulting final state, and looking it up in the master file. The elements were grouped into blocks, each of which acts between small sub-sets of the initial and final state vectors, to avoid ‘thrashing’ during the matrix multiplications. Within each sub-set, the initial and final addresses can each be fitted into a half-integer, so that the matrix elements occupied 35 Gbyte of storage over all.
Finally, the lowest eigenvalue and eigenvector of the Hamiltonian were found in each sector, using the conjugate gradient method. Nightingale et al (1993) showed that the conjugate gradient method converges faster than the Lanczos method for large problems such as this. We found that the eigenvalue converged to an accuracy of 1 part in $`10^{10}`$ in 20-25 iterations for the 6x6 lattice in the neighbourhood of the critical point.
Having determined the quantities of interest for each finite lattice, it is then necessary to make an extrapolation to $`L\mathrm{}`$, to estimate the bulk behaviour of the system. In the vicinity of the critical point, the finite-lattice sequence will typically behave as
$$f_L=f_{\mathrm{}}+a_1L^{\omega _1}+a_2L^{\omega _2}+\mathrm{}$$
(3)
where the $`\omega _i`$ are non-integer exponents, in general (Barber 1983). The problem of extrapolating such a sequence has been discussed in several reviews (Smith and Ford 1982, Barber and Hamer 1982, Guttmann 1989). We have employed a number of different algorithms, including:
* the Neville table (Guttmann 1989), which is best suited to a simple polynomial sequence, with integer exponents $`\omega _i`$;
* the alternating VBS algorithm (van den Broeck and Schwartz 1979, Barber and Hamer 1982), which can give good convergence for sequences of type (3), but needs at least two iterations to work well;
* the Lubkin algorithm (1952), which is more suitable for short sequences;
* the Bulirsch-Stoer algorithm (1964), which has been applied in this context by Henkel and Patkos (1987), and Henkel and Schütz (1988). This algorithm involves an explicit parameter $`\omega `$ which can be optimized to match the leading power-law correction. It has been claimed by Henkel and Schütz that the algorithm is more robust and more accurate than the VBS algorithm, especially for short sequences.
## III Results
### A Finite-lattice data
The pseudo-critical point at lattice size L can be defined according to finite-size scaling theory (Barber 1983) as the coupling $`x_L`$ such that
$$R_L(x_L)=1$$
(4)
where $`R_L(x)`$ is the scaled energy-gap ratio
$$R_L(x)=\frac{LF_L(x)}{(L1)F_{L1}(x)}$$
(5)
and $`F_L(x)`$ is the energy gap for lattice size L. This point is found by calculating the energy eigenvalues at a cluster of 5 equally spaced points in the neighbourhood of $`x_L`$, and then finding $`x_L`$ by interpolation between them. The spacing between the points was chosen as $`\mathrm{\Delta }x=0.001`$, estimated to balance the truncation and round-off errors in the calculation. The values of all other observables can then be estimated at $`x_L`$ by the same finite-difference interpolation procedures. Tables 1 and 2 list the pseudo-critical points $`x_L`$, and the values of the calculated observables at coupling $`x_L`$ for each pair of lattice sizes L and (L-1), for L = 2,3,4,5 and 6. The values of $`x_L`$ for $`L=2`$ to 5 listed in Table 1 agree through six figures with those calculated previously (Hamer 1983).
Table 1 lists values for the ground-state energy per site $`ϵ_{0,L}`$ for lattice size L, and its derivatives $`ϵ_{0,L}^{}`$ and $`ϵ_{o,L}^{\prime \prime }`$, where the prime denotes differentiation with respect to x. The values are expected to be accurate to the figures quoted (or better) as regards round-off error. The truncation error in the 5-point interpolation process is harder to estimate, since it involves unknown higher derivatives of $`ϵ_0`$, but we estimate it should be no more than about 1 part in $`10^{12}`$ for $`ϵ_0`$ and 1 part in $`10^6`$ for $`ϵ_0^{\prime \prime }`$.
We have also listed values in Table 1 for the magnetic susceptibility, defined by
$$\chi _L=\frac{1}{L^2}\frac{^2E_{0,L}(x,h)}{h^2}|_{h=0}$$
(6)
This derivative was also estimated by a finite difference method, using a cluster of 5 data points around $`h=0`$, with a spacing $`\mathrm{\Delta }h=0.0003`$, giving an estimated truncation error of no more than 1 part in $`10^6`$ in the susceptibility. This calculation was a little too large to carry through for $`L=6`$, with the facilities available.
Table 2 lists the energy gap $`F_L`$ between the odd and even sectors, and its derivatives $`F_L^{}`$ and $`F_L^{\prime \prime }`$, at each $`x_L`$. Values are also listed here for the quantity $`M_L`$ defined by
$$M_L=0|\sigma _1(1)|1$$
(7)
where $`|0`$, $`|1`$ are the lowest-lying energy eigenvectors in the even and odd sectors, respectively. It can be shown (Yang 1952, Uzelac 1980, Hamer 1982) that this quantity converges towards the spontaneous magnetization in the bulk limit. Unfortunately, for technical reasons we were again unable to calculate this quantity for L=6. Since the accuracy of the wavefunction is only the square root of that of the eigenvalue, the round-off error in these values is expected to be about 1 part in $`10^6`$.
### B Critical Point
The sequence of pseudo-critical points $`x_L`$ converges rapidly, as can be seen in Fig. 1, where $`x_L`$ is plotted against $`1/L^4`$. To estimate the bulk limit ($`L\mathrm{}`$) of this sequence, we have employed various algorithms discussed above, as well as a simple polynomial fit in $`1/L^4`$ and higher powers.
Our final estimate of the critical point is
$$x_c=0.32841(2)$$
(8)
This is consistent with our earlier finite-size estimate of $`x_c=0.3289(10)`$ (Hamer 1983), but nearly two orders of magnitude more accurate. Henkel (1987) obtained an improved estimate $`x_c=0.3282(1)`$ from lattices up to 5x5 sites.
### C Critical Indices
Finite-size scaling theory (Barber 1983) also tells us how to estimate the critical indices for the model. The finite-lattice susceptibility $`\chi _L`$, for instance, is predicted to scale at the critical point like
$$\chi _L(x_c)L^{\gamma /\nu },L\mathrm{}$$
(9)
and hence one finds that
$$L(1\frac{\chi _L(x_L)}{\chi _{L1}(x_L)})\frac{\gamma }{\nu },L\mathrm{}$$
(10)
Similarly, ratios of the finite-lattice ‘magnetizations’ (equation 7) give estimates of $`\beta /\nu `$. Finally, estimates of the index $`1/\nu `$ can be obtained from the Callan-Symanzik ‘beta function’ (Barber 1983),
$$\beta _L(x)/g=\frac{F_L(x)}{(F_L(x)2xF_L^{}(x))}$$
(11)
via
$$L(1\frac{\beta _L(x_L)}{\beta _{L1}(x_L)})\frac{1}{\nu },L\mathrm{}$$
(12)
One would expect to obtain estimates of the ratio $`\alpha /\nu `$ in a similar fashion from the ‘specific heat’,
$$C_L(x)=\frac{x^2}{L^2}\frac{^2ϵ_0}{x^2},$$
(13)
but it is known (Hamer 1983) that these estimates are very poor, too high by a factor of nearly 2. The reason is easily found: the ground-state energy or specific heat contains a ‘regular’ or analytic piece as well as the singular term (Privman and Fisher 1984). Henkel (1987) has cleverly sidestepped this problem, using a transition amplitude to find $`\alpha /\nu `$, in analogy to equation (7). Here, we eliminate the regular term by subtracting:
$$ϵ_{0,L}^{\prime \prime }ϵ_{0,L1}^{\prime \prime }L^{\alpha /\nu 1},L\mathrm{},$$
(14)
and using successive ratios of these differences to estimate $`1\alpha /\nu `$ The estimates so obtained for the critical index ratios are listed in Table 3.
Alternatively, ‘logarithmic’ estimates of the critical indices may be obtained as follows:
$$\frac{\mathrm{ln}[\chi _L(x_L)/\chi _{L1}(x_L)]}{\mathrm{ln}[L/(L1)]}\frac{\gamma }{\nu },L\mathrm{}$$
(15)
These alternative estimates are listed in Table 4. The finite-size corrections are generally smaller for the logarithmic estimates.
These finite-size estimates of the critical indices agree closely with the previous calculation of Hamer (1983) up to $`L=5`$; and remarkably enough, most of them agree to within 4 significant figures with the equivalent results obtained for the triangular lattice (Hamer and Johnson 1986, Price et al 1993).
The same algorithms mentioned above have been employed to extrapolate these sequences to their bulk limit. The sequences are very short, and may have slight irregularities, so that the tabular algorithms are generally no more accurate than simple graphical methods or polynomial fits in the extrapolation. The resulting estimates are listed at the foot of Tables 3 and 4. The errors in these estimates are inevitably rather subjective, but the variation between different algorithms gives some indication of the likely error.
Figure 2 graphs the estimates of $`1/\nu `$ from Tables 3 and 4 as a function of $`1/L`$: it can be seen that the behaviour is almost precisely linear for the estimates from Table 3. Correspondingly, the Neville tables and polynomial fits give stable results, while the Lubkin and Bulirsch-Stoer algorithms give less stable results, possibly a little higher. We conclude that
$$\frac{1}{\nu }=1.591(2)$$
(16)
The estimates for $`\alpha /\nu `$ are not quite so well-behaved, but our final estimate is
$$\frac{\alpha }{\nu }=0.16(1)$$
(17)
This is a much better result than can be obtained directly from the specific heat, equation (13).
The estimates for the other indices $`\beta /\nu `$ and $`\gamma /\nu `$ are rapidly convergent, but we only have data up to $`L=5`$, which were known previously. We find
$$\frac{\beta }{\nu }=0.522(2)$$
(18)
$$\frac{\gamma }{\nu }=1.96(1)$$
(19)
### D Energy Amplitudes
The finite-size behaviour of the energy gap at the critical point is
$$F_L(x_L)\frac{A_1}{L},L\mathrm{}$$
(20)
so the amplitude $`A_1`$ can be estimated by
$$LF_L(x_L)A_1,L\mathrm{}$$
(21)
The sequence of estimates for $`A_1`$ is shown in Figure 3. It extrapolates to a value
$$A_1=1.39(1)$$
(22)
A value of 1.42 was previously estimated by Henkel (1987).
The finite-size scaling behaviour of the ground-state energy per site $`ϵ_0`$ at the pseudo-critical point is shown in Figure 4. The finite-size scaling corrections appear to decrease like $`1/L^3`$, in accordance with the Privman-Fisher scaling hypothesis (1984), which states that the singular part of the free energy density of a system of finite size L should scale as $`L^d`$ (here d = 3). A polynomial fit on this assumption gives
$$ϵ_{0,L}(x_L)ϵ_0^{}\frac{A_0}{L^3},L\mathrm{}$$
(23)
with
$$ϵ_0^{}=0.624(1)$$
(24)
and
$$A_0=0.38(5)$$
(25)
Further evidence for this power-law behaviour can be obtained as follows. Suppose
$$ϵ_{0,L}(x_L)ϵ_0^{}A_0/L^p,L\mathrm{}$$
(26)
then
$$L[1\frac{(ϵ_{0,L}(x_L)ϵ_{0,L1}(x_L))}{(ϵ_{0,L1}(x_{L1})ϵ_{0,L2}(x_{L1}))}]p_Lp\text{as }L\mathrm{}$$
(27)
and
$$\mathrm{ln}[\frac{ϵ_{0,L}(x_L)ϵ_{0,L1}(x_L)}{ϵ_{0,L1}(x_{L1})ϵ_{0,L2}(x_{L1})}]/\mathrm{ln}[L/L1]p\text{as }L\mathrm{}$$
(28)
The sequences of finite lattice estimates for p are shown in Figure 5. It can be seen that the ‘linear’ sequence comes down towards 3 from above, whilst the ‘logarithmic’ sequence comes up towards 3 from below. The sequences are a little irregular, however, and the best estimate we can obtain for the bulk limit is
$$p=2.8(2)$$
(29)
a little lower than, but still consistent with 3.
Assuming that $`p=3`$, the scaling amplitude $`A_0`$ for the ground-state energy can be found by
$$\frac{L^4}{3}(ϵ_{0,L}(x_L)ϵ_{0,L1}(x_L))A_0,L\mathrm{}$$
(30)
The sequence of estimates for $`A_0`$ is graphed in Figure 6, and extrapolates to a value
$$A_0=0.35(2),$$
(31)
which is in reasonable agreement with equation (25). Henkel (1987) previously obtained an estimate of 0.39 for this quantity (allowing for the different normalization of his Hamiltonian). Mon (1985) obtained a Monte Carlo estimate of the corresponding free energy amplitude in the 3D classical model.
In order to ‘calibrate’ this result, we need to know the “speed of light’ v, or in other words the scale factor needed in this model to make the long-range correlations isotropic in space and time at the critical point. We have attempted to estimate this using the dispersion relation for the lowest excited state at the critical point, expected to be of the form
$$E(k)=vk$$
(32)
in the bulk system. We have calculated the finite-lattice eigenvalues for low-lying excited states with non-zero momntum for lattice sizes L=2 to 5, and set
$$v_L=\frac{L}{2\pi }(F_L(x_L,\frac{2\pi }{L})F_L(x_L,0))v,L\mathrm{}$$
(33)
where $`F_L(x,k)`$ is the energy at coupling x for momentum k. Figure 7 shows the sequence of finite-lattice estimates for v as a function of $`1/L`$. They extrapolate to a bulk value
$$v=0.99(3)$$
(34)
The ratio $`A_0/v`$ should be a universal number, independent of the normalization of the Hamiltonian. From the results above, we find
$$\frac{A_0}{v}=0.35(2),$$
(35)
to be compared with values of 0.719 expected according to effective field theory for a single free boson (Hasenfratz and Niedermayer 1993), or 0.211 for a single free fermion degree of freedom (Appendix). The result (3.31) matches neither of these values. This is not surprising, since the effective field theory at the critical point is expected to be a non-trivial interacting theory. It might be possible to estimate this quantity via the $`ϵ`$-expansion, using a Landau-Ginzburg effective field theory. This has not yet been done, as far as we are aware.
## IV Conclusions
We have calculated the lowest-lying energy eigenvalues of the transverse Ising model on the square lattice with periodic boundary conditions for lattice sizes up to 6x6 sites, using the conjugate gradient method. Finite-size scaling theory has been employed to estimate the critical parameters, which are compared with previous estimates in Table 5.
It can be seen that our present estimates agree well with earlier finite-size scaling results. We have achieved a substantial increase in accuracy for the critical point, but only a more modest increase for the critical index $`\nu `$. The results appear very compatible with previous series analyses, and also with recent estimates for the classical 3D Ising model, and field theory. This provides further confirmation of the universality between these transitions. Finally, it can be seen that the accuracy of the exponents for the quantum model is now not very far behind that for the classical model.
We have also estimated the finite-size scaling amplitudes for the energy eigenvalues at the critical point. For the spin gap we find
$$A_1=1.39(1)$$
(36)
to be compared with a previous estimate by Henkel (1987) of $`A_1=1.42`$.
For the ground-state energy, we have shown evidence that
$$ϵ_{0,L}(x_L)ϵ_0^{}\frac{A_0}{L^3},L\mathrm{}$$
(37)
and estimated
$$A_0=0.35(2),$$
(38)
to be compared with a previous estimate of 0.39 by Henkel (1987). It should be possible to predict this amplitude from Landau-Ginzburg effective field theory.
An extension to 7x7 sites of these exact diagonalization calculations is hardly feasible at the present time, but there are some very precise approximate methods now available, such as the density matrix renormalization group (White 1992) and path integral Monte Carlo techniques (Sandvik 1992). These might well be able to extend the results to larger lattice sizes, and allow much improved finite-size scaling estimates of the critical parameters. They could also confirm whether or not the Casimir energy scales as in equation (37). Our exact diagonalization results should provide a useful calibration for such studies. We look forward to seeing such calculations in the future.
We have chosen here to work on the square lattice rather than the triangular one, because the Hamiltonian matrix is somewhat smaller, and the leading finite-size corrections are expected to be much the same for both lattices. There are, however, some hints of irregularity or alternating behaviour in some of the square lattice sequences. It might well be that the triangular lattice results are smoother.
###### Acknowledgements.
I would like to thank Dr. P.F. Price and Prof. I. Affleck for useful discussions. Part of this work was carried out while on study leave at the Institute for Theoretical Physics, University of California at Santa Barbara, and at the Centre for Nonlinear Studies, Los Alamos National Laboratory. I would like to thank Prof. R. Singh and the organizers of the Workshop on ‘Magnetic Materials in Novel Materials and Geometries’ for their hospitality in Santa Barbara, and Dr. J. Gubernatis for his kind hospitality in Los Alamos. The calculations were performed using facilities at the New South Wales Centre for Parallel Computing, and at the Centre for Nonlinear Studies, Los Alamos. I am very grateful to Prof. R. Standish and Mr. D. Neal for their assistance in this regard. This research was supported in part by the National Science Foundation under grant no. PHY94-07194, and also by a grant from the Australian Research Council. Appendix The finite-size scaling amplitude for the ground-state energy (“Casimir amplitude”) can be calculated for free fields as follows.
### 1 Free boson case
The zero-point energy of a free boson field is given in $`d`$ space dimensions by
$$E_0=\frac{1}{2}\underset{k}{}\omega _k$$
(39)
i.e. $`\omega _k/2`$ for each momentum mode. On a lattice, the free particle Hamiltonian can be written in a finite-difference form
$$H=\frac{1}{2}\underset{𝐧}{}[\dot{\varphi }^2(𝐧)+\underset{i=1}{\overset{d}{}}(\frac{\varphi (𝐧+𝐢)\varphi (𝐧𝐢)}{2})^2]$$
(40)
where the lattice spacing has been set to 1. The eigenmodes are plane waves
$$\varphi (𝐧,t)=\frac{1}{N^{1/2}}\underset{𝐤}{}[a_ke^{i(𝐤.𝐧\omega _kt)}+a_k^{}e^{i(𝐤.𝐧\omega _kt)}]$$
(41)
whence
$$\omega _k=[2\underset{i=1}{\overset{d}{}}(1\mathrm{cos}k_i)]^{1/2}$$
(42)
and for periodic boundary conditions the allowed momenta are
$$k_i=\frac{2\pi }{L}l_i,l_i=0,1,2,\mathrm{}$$
(43)
for a lattice of $`N=L^d`$ sites. Hence
$$E_0=\frac{1}{2}\underset{𝐤}{}[2\underset{i=1}{\overset{d}{}}(1\mathrm{cos}k_i)]^{1/2}$$
(44)
i.e.
$$ϵ_0=\frac{1}{L^d}\underset{\{l_i\}=0}{\overset{L1}{}}[\underset{i=1}{\overset{d}{}}\mathrm{sin}^2(\frac{\pi l_i}{L})]^{1/2}$$
(45)
Now the leading finite-size correction to this sum arises from the infrared (small momentum) behaviour of the lattice sum (Hasenfratz and Leutwyler 1990), and does not depend on the cutoff or regularization at large momentum. Thus we may approximate for our purposes
$$ϵ_0\frac{1}{L^d}\underset{\{l_i\}=\mathrm{}}{\overset{\mathrm{}}{}}[\underset{i=1}{\overset{d}{}}(\frac{\pi l_i}{L})^2]^{1/2}$$
(46)
Now use the Poisson resummation formula
$$\underset{𝐦=\mathrm{}}{\overset{+\mathrm{}}{}}f(𝐦L)=\frac{1}{L^d}\underset{𝐧=\mathrm{}}{\overset{+\mathrm{}}{}}g(\frac{2\pi 𝐧}{L})$$
(47)
with
$$g(𝐤)=_{\mathrm{}}^+\mathrm{}e^{i𝐤.𝐱}f(𝐱)d^dx$$
(48)
to show
$$ϵ_0^{}=\frac{1}{4\pi }\underset{\{m_i\}^{}=\mathrm{}}{\overset{+\mathrm{}}{}}_0^{\mathrm{}}k^d𝑑k\frac{J_{d/21}(kx)}{(2\pi kx)^{d/21}}(x=L|𝐦|)$$
(49)
$$=\frac{\mathrm{\Gamma }(\frac{d+1}{2})}{2\pi ^{\frac{d+1}{2}}L^{d+1}}\underset{\{m_i\}^{}=\mathrm{}}{\overset{+\mathrm{}}{}}\frac{1}{|𝐦|^{d+1}}$$
(50)
The dash here implies removal of the term $`𝐦=0`$, which corresponds to the (infinite, non-universal) bulk ground-state energy per site, which we simply drop.
The sum involved here is a generalization of the Riemann zeta function. It gives
$$ϵ_0^{}=\frac{A_0}{L^{d+1}}$$
(51)
where for $`d=1`$, $`A_0`$ is easily evaluated
$$A_0=\frac{\pi }{6}=0.5236$$
(52)
the result familiar from conformal field theory. For higher dimensions, we have evaluated the sum numerically
$$d=2:A_0=0.7189$$
(53)
$$d=3:A_0=0.8375$$
(54)
The result for $`d=2`$ was given previously by Hasenfratz and Niedermayer (1993) .
### 2 Free fermion case
A similar naive argument can be given for the case of a single species of free Weyl (spinless) fermions. The filled Dirac sea has energy
$$E_0=\underset{k}{}\omega _k$$
(55)
where again
$$\omega _k=[2\underset{i=1}{\overset{d}{}}(1\mathrm{cos}k_i)]^{1/2}$$
(56)
and we assume antiperiodic boundary conditions for the fermions
$$k_i=\frac{\pi }{L}(2l_i+1),l_i=0,1,2,\mathrm{}$$
(57)
Hence
$$ϵ_0\frac{1}{L^d}\underset{\{l_i\}=\mathrm{}}{\overset{\mathrm{}}{}}[\underset{i=1}{\overset{d}{}}(\frac{\pi (2l_i+1)}{2L})^2]^{1/2}$$
(58)
Use the Poisson resummation formula again to find
$$ϵ_0^{}=\frac{1}{4\pi }\underset{\{m_i\}^{}=\mathrm{}}{\overset{+\mathrm{}}{}}(1)^{_im_i}_0^{\mathrm{}}k^d𝑑k\frac{J_{d/21}(kx)}{(2\pi kx)^{d/21}}(x=L|𝐦|)$$
(59)
$$=\frac{\mathrm{\Gamma }(\frac{d+1}{2})}{2\pi ^{\frac{d+1}{2}}L^{d+1}}\underset{\{m_i\}^{}=\mathrm{}}{\overset{+\mathrm{}}{}}\frac{(1)^{_im_i}}{|𝐦|^{d+1}}$$
(60)
For $`d=1`$ $`A_0`$ is easily evaluated to give
$$A_0=\frac{\pi }{12}=0.2618$$
(61)
also familiar from conformal field theory; while for higher dimensions, we find numerically
$$d=2:A_0=0.2106$$
(62)
$$d=3:A_0=0.1957$$
(63)
These numbers have not been obtained previously, as far as we are aware.
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# Untitled Document
New Non-Abelian Zeta Functions for Curves over Finite Fields
Lin WENG
Graduate School of Mathematics, Nagoya University, Chikusa-ku, Nagoya 464-8602, Japan
In this paper, we introduce and study two new types of non-abelian zeta functions for curves over finite fields, which are defined by using (moduli spaces of) semi-stable vector bundles and non-stable bundles. A Riemann-Weil type hypothesis is formulated for zeta functions associated to semi-stable bundles, which we think is more canonical than the other one. All this is motivated by (and hence explains in a certain sense) our work on non-abelian zeta functions for number fields.
§1. Restricted Non-Abelian Zeta Functions
Let $`C`$ be a regular, irreducible, projective curve of genus $`g`$ defined over an algebraically closed field, and $`L`$ a line bundle over $`C`$. Then, by a result of Mumford, we know that moduli space $`_r(L)`$ of semi-stable vector bundles of rank $`r`$ over $`C`$ with $`L`$ as determinants is projective (\[Mu\]). And, if $`C`$ and $`L`$ are defined over the finite field $`𝐅_q`$ with $`q`$ elements, then by taking a finite base field extension if necessary, we may assume that $`_r(L)`$ is also defined over $`𝐅_q`$. Moreover, by a result of Harder and Narasimhan \[HN\], the $`𝐅_q`$-rational points of $`_r(L)`$ are exactly rank $`r`$ semi-stable vector bundles of $`C`$ defined over $`𝐅_q`$ with $`L`$ as determinants.
Clearly, the canonical line bundle $`K_C`$ of $`C`$ is defined over $`𝐅_q`$ as well. Thus, for all $`n𝐙`$, we obtain the following natural isomorphisms defined over $`𝐅_q`$ too:
$$\begin{array}{ccccccc}_r(L)& & _r(LK_C^{nr});& & _r(L)& & _r(L^1K_C^{nr})\\ E& & EK_C^n;& & E& & E^{}K_C^n,\end{array}$$
where $`E^{}`$ denotes the dual of $`E`$.
Define the so-called Harder-Narasimhan number $`M_{C,r,L}`$ by setting
$$M_{C,r,L}:=\underset{E_r(L)(𝐅_q)}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}.$$
Here $`\mathrm{Aut}(E)`$ denotes the automorphism group of $`E`$. (See e.g., \[HN\], \[Se\].) Easily, we have, for all $`n𝐙`$,
$$M_{C,r,L}=M_{C,r,LK_C^{nr}}=M_{C,r,L^1K_C^{nr}}.$$
The Harder-Narasimhan numbers may be calculated by counting contributions of non-stable bundles, with the help of a famuous result of Siegel on quadratic forms and Harder-Narasimhan’s filtration for vector bundles. (For details, see e.g., \[HN\], \[DR\] and \[AB\].) But in this section, we relate them with a new type of (restricted) non-abelian zeta functions.
Denote by $`h^i`$ the dimension of cohomology groups $`H^i`$, $`i=0,1`$, by $`\chi (C,E):=h^0(C,E)h^1(C,E)`$ the Euler-Poincaré characteristic, and by $`d(E)`$ the degree of $`E`$. Set
$$_{r;L}(𝐅_q):=_{n𝐙}\left(_r(LK_C^{nr})(𝐅_q)_r(L^1K_C^{nr})(𝐅_q)\right).$$
Definition. Let $`C`$ be a regular, geometrically irreducible, projective curve defined over the finite field $`𝐅_q`$ with $`q`$ elements. With respect to any fixed positive integer $`r𝐙_{>0}`$, and any line bundle $`L`$ of $`C`$ defined over $`𝐅_q`$, define a weight $`r`$ and level $`L`$ restricted non-abelian zeta function $`\xi _{r,L}(s)`$ of $`C`$ by
$$\xi _{r,L}(s):=\underset{E_{r,L}(𝐅_q)}{}\frac{q^{h^0(C,E)}1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s,\mathrm{Re}(s)>1.$$
Clearly, $`\xi _{r,L}(s)=\xi _{r,L^{\pm 1}K_C^{nr}}(s).`$ Hence, from now on, we assume $`0d(L)r(g1)`$.
Theorem 1. For $`\mathrm{Re}(s)>1`$, we have
$$\begin{array}{cc}\hfill \xi _{r,L}(s)=& \frac{1}{2}\underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(\left(q^{\chi (C,E)}\right)^s+\left(q^{\chi (C,E)}\right)^{s1}\right)\hfill \\ & +[\left(q^{d(L)r(g1)}\right)^{1s}\frac{1}{q^{r(2g2)(s1)}1}+\left(q^{d(L)r(g1)}\right)^s\frac{1}{q^{sr(2g2)}1}\hfill \\ & +\left(q^{d(L)r(g1)}\right)^{s1}\frac{1}{q^{r(2g2)(s1)}1}+\left(q^{d(L)r(g1))}\right)^s\frac{1}{q^{sr(2g2)}1}]M_{C,r,L}.\hfill \end{array}$$
$`()`$
Hence, in particular, we have
(a) $`\xi _{r,L}(s)`$ can be meromorphically extended to the whole complex $`t=q^s`$-plane;
(b) the extension, denoted also by $`\xi _{r,L}(s)`$, has simple poles at $`s=0`$ and $`s=1`$ with the Harder-Narasimhan number $`M_{C,r,L}`$ as residues;
(c) $`\xi _{r,L}(s)`$ satisfies the functional equation $`\xi _{r,L}(s)=\xi _{r,L}(1s).`$ Proof of the Theorem. It suffices to prove (\*). For this, we proceed as follows. Recall that for $`E`$ semi-stable, if $`d(E)>r(2g2)`$, then $`h^1(C,E)=0`$; while if $`d(E)<0`$, then $`h^0(C,E)=0`$. Thus,
$$\begin{array}{cc}& \xi _{r,L}(s)=\underset{E_{r,L}(𝐅_q);d(E)0}{}\frac{q^{h^0(C,E)}1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\hfill \\ \hfill =& \underset{E_{r,L}(𝐅_q);d(E)0}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\underset{E_{r,L}(𝐅_q);d(E)0}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\hfill \\ \hfill =& \underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\hfill \\ & +\underset{E_{r,L}(𝐅_q);d(E)>r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\underset{E_{r,L}(𝐅_q);d(E)0}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s.\hfill \end{array}$$
So $`\xi _{r,L}(s)=S_{r,L}^{0d(E)r(2g2)}(s)+T_{r,L}(s)`$ if we set
$$S_{r,L}^{0d(E)r(2g2)}(s):=\underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s,$$
and
$$\begin{array}{cc}\hfill T_{r,L}(s):=& \underset{E_{r,L}(𝐅_q);d(E)>r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\underset{E_{r,L}(𝐅_q);d(E)0}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s.\hfill \end{array}$$
Lemma. The function $`t^{r(g1)}S_{r,L}^{0d(E)r(2g2)}(s)`$ is holomorphic in $`t=q^s`$. Moreover
$$S_{r,L}^{0d(E)r(2g2)}(s)=\frac{1}{2}\underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(\left(q^{\chi (C,E)}\right)^s+\left(q^{\chi (C,E)}\right)^{s1}\right)$$
and hence satisfies the functional equation
$$S_{r,L}^{0d(E)r(2g2)}(s)=S_{r,L}^{0d(E)r(2g2)}(1s).$$
Proof of the lemma. Holomorphicity is clear, as only finitely many terms are involved. The functional equation comes from the Riemann-Roch theorem. Indeed,
$$\begin{array}{cc}\hfill S_{r,L}^{0d(E)r(2g2)}(s)=& \frac{1}{2}(\underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\hfill \\ & +\underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s)\hfill \\ \hfill =& \frac{1}{2}(\underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (E)}\right)^s\hfill \\ & +\underset{E^{}K_C_{r,L}(𝐅_q);0d(E^{}K_C)r(2g2)}{}\frac{q^{h^0(C,E^{}K_C)}}{\mathrm{\#}\mathrm{Aut}(E^{}K_C)}\left(q^{\chi (C,E^{}K_C)}\right)^s)\hfill \\ \hfill =& \frac{1}{2}(\underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\hfill \\ & +\underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^0(C,E^{}K_C)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E^{}K_C)}\right)^s)\hfill \\ \hfill =& \frac{1}{2}(\underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\hfill \\ & +\underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^1(C,E^{}K_C)+d(E^{}K_C)r(g1)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E^{}K_C)}\right)^s)\hfill \\ \hfill =& \frac{1}{2}(\underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\hfill \\ & +\underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^0(C,E)d(E)+r(g1)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)+r(g1)}\right)^s)\hfill \\ \hfill =& \frac{1}{2}\underset{E_{r,L}(𝐅_q);0d(E)r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(\left(q^{\chi (C,E)}\right)^s+\left(q^{\chi (C,E)}\right)^{s1}\right).\hfill \end{array}$$
This completes the proof of the lemma. As for $`T_{r,L}(s)`$, clearly, by the vanishing recalled at the beginning and the Riemann-Roch, we get
$$\begin{array}{cc}\hfill T_{r,L}(s)=& \underset{E_{r,L}(𝐅_q);d(E)>r(2g2)}{}\frac{q^{d(E)r(g1)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\underset{E_{r,L}(𝐅_q);d(E)0}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\hfill \\ \hfill =& \underset{E_{r,L}(𝐅_q);d(E)>r(2g2)}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^{1s}\underset{E_{r,L}(𝐅_q);d(E)0}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\hfill \\ \hfill =& \left(\underset{E_r(LK_C^{rn})(𝐅_q);d(E)>r(2g2)}{}+\underset{E_r(L^1K_C^{rn})(𝐅_q);d(E)>r(2g2)}{}\right)\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^{1s}\hfill \\ & \left(\underset{E_r(LK_C^{rn})(𝐅_q);d(E)0}{}+\underset{E_r(L^1K_C^{rn})(𝐅_q);d(E)0}{}\right)\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s\hfill \end{array}$$
$$\begin{array}{cc}\hfill =& (\underset{E_r(LK_C^{rn})(𝐅_q);d(E)>r(2g2)}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^{1s}\hfill \\ & \underset{E_r(L^1K_C^{rn})(𝐅_q);d(E)0}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s)\hfill \\ & +(\underset{E_r(L^1K_C^{rn})(𝐅_q);d(E)>r(2g2)}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^{1s}\hfill \\ & \underset{E_r(LK_C^{rn})(𝐅_q);d(E)0}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{\chi (C,E)}\right)^s)\hfill \\ \hfill =& [(\underset{n;d(E)=d(L)+nr(2g2)>r(2g2)}{}\left(q^{\chi (C,E)}\right)^{1s}\underset{n;d(E)=d(L)+nr(2g2)0}{}\left(q^{\chi (C,E)}\right)^s)\hfill \\ & +(\underset{n;d(E)=d(L)+nr(2g2)>r(2g2)}{}\left(q^{\chi (C,E)}\right)^{1s}\underset{n;d(E)=d(L)+nr(2g2)0}{}\left(q^{\chi (C,E)}\right)^s)]M_{C,r,L},\hfill \end{array}$$
in which we omit stating that $`E`$’s are semi-stable with determinants $`L^\pm K_C^{nr}`$. Therefore,
$$\begin{array}{cc}\hfill T_{r,L}(s)=& [(\underset{n=1}{\overset{\mathrm{}}{}}\left(q^{d(L)+nr(2g2)r(g1)}\right)^{1s}\underset{n=1}{\overset{\mathrm{}}{}}\left(q^{d(L)+nr(2g2)r(g1)}\right)^s)\hfill \\ & +(\underset{n=2}{\overset{\mathrm{}}{}}\left(q^{d(L)+nr(2g2)r(g1)}\right)^{1s}\underset{n=0}{\overset{\mathrm{}}{}}\left(q^{d(L)+nr(2g2)r(g1)}\right)^s)]M_{C,r,L}\hfill \\ \hfill =& [(\underset{n=1}{\overset{\mathrm{}}{}}\left(q^{\chi (E_0)+nr(2g2)}\right)^{1s}\underset{n=1}{\overset{\mathrm{}}{}}\left(q^{\chi (C,E_0)+nr(2g2)}\right)^s)\hfill \\ & +(\underset{n=2}{\overset{\mathrm{}}{}}\left(q^{\chi (C,E_0)+(n1)r(2g2)}\right)^{1s}\underset{n=0}{\overset{\mathrm{}}{}}\left(q^{\chi (C,E_0)+nr(2g2)}\right)^s)]M_{C,r,L},\hfill \end{array}$$
for any $`E_0_r(L)(𝐅_q)`$. Thus, if $`q^{r(2g2)(1s)}<1`$ and $`q^{sr(2g2)}<1`$, by $`0d(L)r(g1)`$,
$$\begin{array}{cc}\hfill T_{r,L}(s)=& [(\underset{n=1}{\overset{\mathrm{}}{}}\left(q^{\chi (C,E_0)+nr(2g2)}\right)^{1s}\underset{n=1}{\overset{\mathrm{}}{}}\left(q^{\chi (C,E_0)+(n1)r(2g2)}\right)^s)\hfill \\ & +(\underset{n=2}{\overset{\mathrm{}}{}}\left(q^{\chi (C,E_0)+(n1)r(2g2)}\right)^{1s}\underset{n=0}{\overset{\mathrm{}}{}}\left(q^{\chi (C,E_0)+nr(2g2)}\right)^s)]M_{C,r,L}\hfill \\ \hfill =& [\left(q^{\chi (C,E_0)}\right)^{1s}\underset{n=1}{\overset{\mathrm{}}{}}\left(q^{r(2g2)(1s)}\right)^n\left(q^{\chi (C,E_0)}\right)^s\underset{n=0}{\overset{\mathrm{}}{}}\left(q^{sr(2g2)}\right)^n\hfill \\ & +\left(q^{\chi (C,E_0)}\right)^{s1}\underset{n=1}{\overset{\mathrm{}}{}}\left(q^{r(2g2)(1s)}\right)^n\left(q^{\chi (C,E_0)}\right)^s\underset{n=0}{\overset{\mathrm{}}{}}\left(q^{sr(2g2)}\right)^nBig]M_{C,r,L}\hfill \\ \hfill =& [\left(q^{\chi (C,E_0)}\right)^{1s}\frac{q^{r(2g2)(1s)}}{1q^{r(2g2)(1s)}}+\left(q^{\chi (C,E_0)}\right)^s\frac{1}{q^{sr(2g2)}1}\hfill \\ & +\left(q^{\chi (C,E_0)}\right)^{s1}\frac{q^{r(2g2)(1s)}}{1q^{r(2g2)(1s)}}+\left(q^{\chi (C,E_0)}\right)^s\frac{1}{q^{sr(2g2)}1}]M_{C,r,L}\hfill \\ \hfill =& [\left(q^{\chi (C,E_0)}\right)^{1s}\frac{1}{q^{r(2g2)(s1)}1}+\left(q^{\chi (C,E_0)}\right)^s\frac{1}{q^{sr(2g2)}1}\hfill \\ & +\left(q^{\chi (C,E_0)}\right)^{s1}\frac{1}{q^{r(2g2)(s1)}1}+\left(q^{\chi (C,E_0)}\right)^s\frac{1}{q^{sr(2g2)}1}]M_{C,r,L}.\hfill \end{array}$$
This then proves the existence of meromorphic extension, the statement for simple poles and their residues, and the functional equation for $`T_{r,L}(s)`$. Thus by the Lemma, we complete the proof of the Theorem.
§2. New Non-Abelian Zeta Functions for Curves over Finite Fields
From definition, clearly,
$$\xi _{1,L}(s)=q^{(g1)s}\underset{E=L^\pm K_C^n;n𝐙}{}\frac{q^{h^0(C,E)}1}{q1}(q^{d(E)})^s.$$
So, by taking a finite sum over suitable $`L`$, we could arrive at
$$q^{(g1)s}\underset{M\mathrm{Pic}(C)(𝐅_q)}{}\frac{q^{h^0(C,E)}1}{q1}(q^{d(E)})^s,$$
which is nothing but
$$q^{(g1)s}\underset{D0}{}N(D)^s,$$
i.e., the standard abelian zeta function $`\zeta _C(s)`$, or Artin zeta function (\[A\]), for $`C`$ times with $`q^{(g1)s}`$. This then suggests the the following discussion.
Denote by $`_{r,d}(C)`$ the moduli space of degree $`d`$ semi-stable vector bundles of rank $`r`$ on $`C`$. Fixed a degree 1 line bundle $`A`$ on $`C`$ which is $`𝐅_q`$-rational. (One may give a proof of this fact by using properties of Artin zeta function $`\zeta _C(s)`$.) Clearly, we then have the following isomorphisms defined over $`𝐅_q`$:
$$\begin{array}{ccc}_r(L)& & _r(LA^{nr})\\ E& & EA^n.\end{array}$$
Thus, as before, we may assume that, if necessary, by taking a finite fields extension, $`_{r,d}(C)`$ are defined over $`𝐅_q`$ as well. Moreover, we know that all $`𝐅_q`$-rational semi-stable vector bundles are indeed an element in
$$_r(C)(𝐅_q):=_d_{r,d}(C)(𝐅_q).$$
Definition. For $`r𝐙_{>0}`$, define a weight $`r`$ non-abelain zeta function $`\zeta _{C,r}(s)`$ of $`C`$ by
$$\zeta _{C,r}(s):=\underset{E_r(C)(𝐅_q)}{}\frac{q^{h^0(C,E)}1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s,\mathrm{Re}(s)>1.$$
Set also $`\xi _{C,r}(s):=q^{r(g1)s}\zeta _{C,r}(s)`$, $`t=q^s`$ and $`Z_{C,r}(t)=\zeta _{C,r}(s)`$. Theorem 2. (a) $`\zeta _{C,r}(s)`$ is absolutely convergent for $`\mathrm{Re}(s)>1`$, can be meromorphically extended to the whole complex $`t=q^s`$ plane which has simple poles at $`s=0`$ and $`s=1`$ with the same residue, and satisfies the functional equation $`\xi _{r,L}(s)=\xi _{r,L}(1s).`$
(b) $`Z_{r,C}(t)={\displaystyle \frac{P(t)}{(1t^r)(1(qt)^r)}}`$ where $`P(t)𝐐[t]`$ is a degree $`2rg`$ polynomial with rational coefficients.
(c) Let $`P(t)=P(0){\displaystyle \underset{i=1}{\overset{2rg}{}}}(1\omega _it),`$ then after a suitable arrangement, we have
$$\omega _i\omega _{2rgi}=q,i=1,\mathrm{},rg.$$
From the proof given below, one can give the precise values of the residues in terms of a certain combination of Harder-Narasimhan numbers. We leave this to the reader. Moreover, (c) suggests that Weil-Riemann Hypothesis (\[W\]) holds also for our new zeta functions. This then leads to the following Riemann-Weil Hypothesis. The reciprocal roots $`\omega _i`$, $`i=1,\mathrm{},2rg`$, of the weight $`r`$ non-abelain zeta functions of curves defined over finite fields satisfy
$$|\omega _i|=q^{\frac{1}{2}},i=1,\mathrm{},2rg.$$
Proof of the Theorem. By a finite field extension, we may assume that all $`_r(L)`$ are defined over $`𝐅_q`$ if $`L`$ is defined over $`𝐅_q`$. With this, (a) is a direct consequence of Theorem 1 in §1 by the fact that
$$_{r,d}(C)(𝐅_q)=_{d=0}^{r1}_{L\mathrm{Pic}^d(C)(𝐅_q)}_{k=1}^{g1}_{r,LA^{kr}}(𝐅_q).$$
Hence we only need to prove (b) and (c).
Let us assume (b), then by the functional equation for $`t=q^s`$, which changes $`t`$ to $`\frac{1}{qt}`$, $`_{i=1}^{2rg}(1\omega _it)`$ changes to $`_{i=1}^{2rg}(1\frac{\omega _i}{qt})`$ accordingly. Thus
$$\underset{i=1}{\overset{2rg}{}}(1\omega _it)=\underset{k=1}{\overset{2rg}{}}(1\frac{1}{\omega _kq}t).$$
This then implies (c).
So it suffices to prove (b), the rationality of our non-abelian zeta function. For this, let us first assume that
$$Z_{r,C}(t)=\frac{P(t)}{(1t^r)(1(qt)^r)}$$
where $`P(t)𝐐[t]`$. Then by the functional equation for $`\zeta _{r,C}(s)`$, we have
$$\frac{P(t)}{(1t^r)(1(qt)^r)}t^{r(g1)}=\frac{P(\frac{1}{qt})}{(1\frac{1}{(qt)^r})(1\frac{1}{t^r})}(qt)^{r(g1)}.$$
Thus, $`\mathrm{deg}(P)2rr(g1)=r(g1)`$ by comparing degrees of $`t`$ of rational functions on both sides. That is to say, we have $`\mathrm{deg}(P)=2rg`$. In this way, we are lead to prove the following Theorem. With the same notation as above, there exists a polynomial $`P(t)𝐐[t]`$ such that
$$Z_{r,C}(t)=\frac{P(t)}{(1t^r)(1(qt)^r)}.$$
Proof. Note that we have the following vanishing: for $`E`$ semi-stable, if $`d(E)>r(2g2)`$, then $`h^1(C,E)=0`$; while if $`d(E)<0`$, then $`h^0(C,E)=0`$. Thus, by definition,
$$\begin{array}{cc}& Z_{r,C}(t)\hfill \\ \hfill =& \underset{E_r(C)(𝐅_q),d(E)0}{}\frac{q^{h^0(C,E)}1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ \hfill =& \underset{E_r(C)(𝐅_q),d(E)0}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\underset{E_r(C)(𝐅_q),d(E)0}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ \hfill =& \underset{E_r(C)(𝐅_q),r(2g2)d(E)0}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s+\underset{E_r(C)(𝐅_q),d(E)>r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ & \underset{E_r(C)(𝐅_q),d(E)0}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s.\hfill \end{array}$$
Set
$$\begin{array}{cc}\hfill I:=& \underset{E_r(C)(𝐅_q),r(2g2)d(E)0}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ \hfill II:=& \underset{E_r(C)(𝐅_q),d(E)>r(2g2)}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ \hfill III:=& \underset{E_r(C)(𝐅_q),d(E)0}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \end{array}$$
Then
$$Z_{r,C}(t)=I(t)+II(t)+III(t).$$
Thus it suffices to prove the following Lemma. With the same notation as above,
$$I(t),\left(1(qt)^r\right)II(t),\mathrm{and}\left(1t^r\right)III(t)𝐐[t].$$
Proof of the Lemma. By definition, $`I(t)𝐐[t]`$. Hence we should prove the assertions for $`II(t)`$ and $`III(t)`$.
We first deal with $`II(t)`$. By the vanishing, we have $`h^0(C,E)=dr(g1)`$. Hence
$$II(t)=q^{r(g1)}\underset{E_r(C)(𝐅_q),d(E)>r(2g2)}{}\frac{(qt)^{d(E)}}{\mathrm{\#}\mathrm{Aut}(E)}.$$
Let $`_{r,k}:=_{r,k}(C)(𝐅_q)`$, then by tensoring with $`A^{}n`$, $`n𝐙_0`$, we have isomorphisms
$$\begin{array}{cc}& _{r,0}_{r,r}\mathrm{}_{r,r(2g2)}_{r,(r+1)(2g2)}\mathrm{}\hfill \\ & _{r,1}_{r,r+1}\mathrm{}_{r,r(2g2)+1}_{r,(r+1)(2g2)+1}\mathrm{}\hfill \\ & \mathrm{}\mathrm{}\mathrm{}\hfill \\ & _{r,r1}_{r,r+(r1)}\mathrm{}_{r,r(2g2)+(r1)}_{r,(r+1)(2g2)+(r1)}\mathrm{}\hfill \end{array}$$
Thus, if we further set
$$M_r:=_{r,r(2g2)+1}\mathrm{}_{r,(r+1)(2g2)+(r1)}_{r,(r+1)(2g2)},$$
then
$$\begin{array}{cc}\hfill II(t)=& q^{r(g1)}\underset{EM_r}{}\underset{n=0}{\overset{\mathrm{}}{}}\frac{(qt)^{d(EA^n)}}{\mathrm{\#}\mathrm{Aut}(E)}\hfill \\ \hfill =& q^{r(g1)}\underset{EM_r}{}\underset{n=0}{\overset{\mathrm{}}{}}\frac{(qt)^{d(E)+nr}}{\mathrm{\#}\mathrm{Aut}(E)}=q^{r(g1)}\underset{EM_r}{}\frac{(qt)^{d(E)}}{\mathrm{\#}\mathrm{Aut}(E)}\underset{n=0}{\overset{\mathrm{}}{}}((qt)^r)^n\hfill \\ \hfill =& q^{r(g1)}\underset{EM_r}{}\frac{(qt)^{d(E)}}{\mathrm{\#}\mathrm{Aut}(E)}\frac{1}{1(qt)^r}.\hfill \end{array}$$
This proves the assertion for $`II(t)`$.
Similarly, for $`III(t)`$, set
$$M_r^{}:=_{r,0}_{r,1}\mathrm{}_{r,r1}.$$
Then
$$\begin{array}{cc}\hfill III(t)=& \underset{EM_r^{}}{}\underset{n=0}{\overset{\mathrm{}}{}}\frac{t^{d(EA^n)}}{\mathrm{\#}\mathrm{Aut}(E)}\hfill \\ \hfill =& \underset{EM_r^{}}{}\underset{n=0}{\overset{\mathrm{}}{}}\frac{t^{d(E)+nr}}{\mathrm{\#}\mathrm{Aut}(E)}=\underset{EM_r}{}\frac{t^{d(E)}}{\mathrm{\#}\mathrm{Aut}(E)}\underset{n=0}{\overset{\mathrm{}}{}}(t^r)^n\hfill \\ \hfill =& \underset{EM_r^{}}{}\frac{t^{d(E)}}{\mathrm{\#}\mathrm{Aut}(E)}\frac{1}{1t^r}.\hfill \end{array}$$
This proves the assertion for $`III(t)`$. All in all, we have proved the lemma, hence the theorem and the theorem itself. As a direct consequence of the rationality of our new zeta functions, i.e., (b) of the Theorem, we have the following
Corollary. (a) For each $`m1`$ there exists suitable number $`N_m`$ such that
$$Z_{r,C}(t)=\mathrm{exp}\left(\underset{m=1}{\overset{\mathrm{}}{}}N_m\frac{t^m}{n}\right).$$
$`()`$
Moreover,
$$N_m=\{\begin{array}{cc}r(1+q^m)_{i=1}^{2rg}\omega _i^m,\hfill & \text{if }r|m\text{;}\hfill \\ _{i=1}^{2rg}\omega _i^m,\hfill & \text{if }rm\text{.}\hfill \end{array}$$
(b) For any positive integer $`a`$ such that $`(a,r)=1`$, if $`\zeta _{i,a},i=1,\mathrm{},a`$ denote all $`a`$-th roots of unity, then
$$\underset{i=1}{\overset{a}{}}Z_{C,r}(\zeta _{i,a}t)=\mathrm{exp}\left(\underset{m=1}{\overset{\mathrm{}}{}}N_{ma}\frac{T^m}{m}\right)$$
$`()`$
where $`T=t^a`$.
Proof. In fact, note that $`\mathrm{log}(1x)=_{m=1}^{\mathrm{}}\frac{x^m}{m}`$, from the rationality of $`Z_{C,r}(t)`$, we get (\**) directly. As for the precise formula of $`N_m`$, we use the following fact which also implies (b) directly:
$$\underset{i=1}{\overset{a}{}}\zeta _{i,a}^m=\{\begin{array}{cc}a,\hfill & \text{if }a|m\text{,}\hfill \\ 0,\hfill & \text{if }am\text{.}\hfill \end{array}$$
Clearly, (\***) suggests that
$$\underset{i=1}{\overset{a}{}}Z_{C,r}(\zeta _{i,a}t)=Z_{C_{ar},r}(T)$$
with $`C_{ar}`$ is obtained from $`C`$ by taking simply extension of constant fields from $`𝐅_q`$ to $`𝐅_{q^{ar}}`$. (Question: Why $`ar`$ not simply $`a`$?) As a matter of fact, $`N_1,\mathrm{},N_r`$ do have Diophantine interpretations. Moreover, there are many important relations in this enumerative aspect of moduli spaces of semi-stable vector bundles. As all are associated to the above Riemann-Weil hypothesis, which we cannot verify now, we leave them to some other occasions.
§3. Non-Stable Contributions
In this section, we justify our new non-abelian zeta functions by looking at non-stable contributions. But for simplicity, we only study the case when $`r=2`$.
So let $`W_d(C)`$ be the collection of isomorphism classes of rank two degree $`d`$ non-stable vector bundles $`E`$ (over $`C`$) defined over $`𝐅_q`$. Thus clearly, we have then bijections $`W_0(C)W_{2n}(C)`$ and $`W_1(C)(W_{2n+1}(C)`$ for all $`n𝐙`$ as for semi-stable vector bundles.
Naturally, one may try to define a new zeta function for $`C`$ by considering a formal summation
$$\underset{EW_d(C),d𝐙}{}\frac{q^{h^0(C,E)}1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s,\mathrm{Re}(s)>1.$$
But this time, we have some troubles. Recall that in semi-stable case, despite that the formal summation is for all $`d𝐙`$, finally, only terms with $`d0`$ contribute, since $`q^{h^0(C,E)}1=0`$ when $`d<0`$. Clearly, for non-stable bundles, this does not hold. Hence, we should modify our definition as follows.
Definition. Let $`C`$ be a regular, geometrically irreducible, projective curve defined over the finite field $`𝐅_q`$ with $`q`$ elements. Define a new zeta functions $`\zeta _{C,\mathrm{ns}}(s)`$ by setting
$$\zeta _{C,\mathrm{ns}}(s):=\underset{EW_d(C),d0}{}\frac{q^{h^0(C,E)}1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s,\mathrm{Re}(s)>1.$$
The main result of this section will be the following
Theorem 3. $`\zeta _{C,\mathrm{ns}}(s)`$ is well-defined and admits a meromorphic extension to the whole complex $`t=q^s`$-plane, which is indeed rational too, i.e., can be written as a quotient of two polynomials with rational coefficients.
Proof. Call any element $`EW_0(C)`$ an $`E_0`$ and $`E_{2n}=E_0A^n`$ with $`A`$ as in the previous section. Similarly, call $`EW_1(C)`$ an $`E_1`$ and $`E_{2n+1}=E_1A^n`$. Then clearly $`\mathrm{\#}\mathrm{Aut}(E_{})=\mathrm{\#}\mathrm{Aut}(E_{2n+})`$. Moreover, for $`E_{}W_{}(C)`$, $`=0,1`$, by using Harder-Narasimhan filtration, over $`𝐅_q`$, we have the following short exact sequences:
$$0L_{+,}E_{}L_,0$$
$`()`$
such that
($`a_0`$) if $`E_0W_0(C)`$, $`d_{+,0}:=d(L_{+,0})=1,2,3,\mathrm{}`$ and $`d_{,0}:=d(L_{,0})=d(L_{+,0})=1,2,3,\mathrm{}`$;
($`a_1`$) if $`E_1W_1(C)`$, $`d_{+,1}:=d(L_{+,1})=1,2,3,\mathrm{}`$ and $`d_{,1}:=d(L_{,1})=1d(L_{+,0})=0,1,2,\mathrm{}`$.
Thus if we set $`L_{+,2n+}=L_{+,}A^n`$ with $`=0,1`$, then we have also
($`a_{2n}`$) if $`E_{2n}W_{2n}(C)`$, $`d_{+,2n}:=d(L_{+,2n})=n+1,n+2,n+3,\mathrm{}`$ and $`d_{,2n}:=d(L_{,2n})=n1,n2,n3,\mathrm{}`$;
($`a_{2n+1}`$) if $`E_{2n+1}W_{2n+1}(C)`$, $`d_{+,2n+1}:=d(L_{+,2n+1})=n+1,n+2,n+3,\mathrm{}`$ and $`d_{,2n+1}:=d(L_{,2n+1})=n,n1,n2,\mathrm{}`$.
Next, we compute $`\mathrm{\#}\mathrm{Aut}(E)`$ according to whether (\*) is trivial or not.
(1) If $`E_m=L_{+,m}L_{,m}`$, the automorphisms consist of $`𝐅_q^{}\times 𝐅_q^{}`$ together with the unipotents of the form $`1+\varphi `$ with
$$\varphi \mathrm{Hom}(L_{,m},L_{+,m})=H^0(C,L_{,m}^1L_{+,m})=\{\begin{array}{cc}H^0(C,L_{,0}^1L_{+,0}),\hfill & \text{if }m=2n\text{;}\hfill \\ H^0(C,L_{,1}^1L_{+,1}),\hfill & \text{if }m=2n+1\text{.}\hfill \end{array}$$
Hence
$$\mathrm{\#}\mathrm{Aut}(E_m)=(q1)^2\{\begin{array}{cc}q^{h^0(C,L_{,0}^1L_{+,0})},\hfill & \text{if }m=2n\text{;}\hfill \\ q^{h^0(C,L_{,1}^1L_{+,1})},\hfill & \text{if }m=2n+1\text{.}\hfill \end{array}$$
(2) For non-trivial extensions, we have only one copy of $`𝐅_q^{}`$ and hence
$$\mathrm{\#}\mathrm{Aut}(E_m)=(q1)\{\begin{array}{cc}q^{h^0(C,L_{,0}^1L_{+,0})},\hfill & \text{if }m=2n\text{;}\hfill \\ q^{h^0(C,L_{,1}^1L_{+,1})},\hfill & \text{if }m=2n+1\text{.}\hfill \end{array}$$
But non-trivial extensions $`E_m`$ correspond to non-zero elements of $`H^1(C,L_{,m}^1L_{+,m})`$ and proportional vectors give isomorphic bundles. Hence the number of isomorphism classes of bundles $`E_m`$ for which $`E`$ is non-trivial is
$$\frac{q^{h^1(C,C,L_{,m}^1L_{+,m})}1}{q1}.$$
Thus in particular, in the summation
$$\underset{EW_d(C),d0}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s,$$
the contributions arising from given $`L_{+,2n+},=0,1`$ are
$$\begin{array}{cc}& \left(\frac{1}{(q1)^2q^{h^0(C,L_,^1L_{+,})}}+\frac{q^{h^1(C,L_,^1L_{+,})}1}{(q1)^2q^{h^0(C,L_,^1L_{+,})}}\right)(q^s)^{2n+}\hfill \\ \hfill =& \frac{q^{h^1(C,L_,^1L_{+,})}}{(q1)^2q^{h^0(C,L_,^1L_{+,})}}(q^s)^{2n+}\hfill \\ \hfill =& q^{d_{+,}+d_,+(g1)}(q^s)^{2n+}\hfill \\ \hfill =& \{\begin{array}{cc}q^{2d_{+,0}+(g1)}(q^s)^{2n},\hfill & \text{if }m=2n\text{;}\hfill \\ q^{2d_{+,0}+g}(q^s)^{2n+1},\hfill & \text{if }m=2n+1\text{.}\hfill \end{array}\hfill \end{array}$$
$`()`$
(This trick is first used by Atiyah and Bott in \[AB\] p.595.)
Therefore, if $`J_0(C)`$ denotes degree zero Jacobian of $`C`$, then
$$\begin{array}{cc}& \underset{EW_d(C),d0}{}\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ \hfill =& \mathrm{\#}J_0(C)(𝐅_q)\underset{d_{+,}=1}{\overset{\mathrm{}}{}}\underset{n0}{\overset{\mathrm{}}{}}\frac{1}{(q1)^2}\left(q^{2d_{+,0}+(g1)}(q^s)^{2n}+q^{2d_{+,0}+g}(q^s)^{2n+1}\right)\hfill \\ \hfill =& \mathrm{\#}J_0(C)(𝐅_q)\frac{1}{(q1)^2}\frac{q^2}{1q^2}q^{g1}\left(1+q^{2s}+q^{4s}+\mathrm{}+q\left(q^s+q^{3s}+q^{5s}+\mathrm{}\right)\right)\hfill \\ \hfill =& \frac{q^{g1}}{(q1)^2(q^21)}\frac{1}{1q^{2s}}(1+q^{1s})\mathrm{\#}J_0(C)(𝐅_q),\hfill \end{array}$$
$`()`$
provided that $`\mathrm{Re}(s)>1`$.
Before going further, we reminder the reader that the summation here involves a double infinite summations, and hence is quite different in nature comparing with semi-stable cases.
Next, consider the part
$$\underset{EW_d(C),d0}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s.$$
Similarly as above, we want to group non-trivial extensions and trivial extensions together so that automorphisms may be calculated easily. Clearly for this purpose, then we should assume that $`h^0(C,E_m)=h^0(C,L_{+,m})+h^0(C,L_{,m})`$, which is clearly the case if $`h^1(C,L_{+,m})=0`$, or even better, $`d_{+,m}>2g2`$. By the complete list for degrees at the beginning of this proof, we see that there are only finitely many $`EW_d(C),d0`$ such that $`d_{+,m}2g2`$. Hence, by definition, up to a polynomial, we are essentially dealing with non-stable bundles such that $`h^0(C,E_m)=h^0(C,L_{+,m})+h^0(C,L_{,m})`$. Moreover, if we can calculate $`h^0(C,L_{+,m})`$ and $`h^0(C,L_{,m})`$ precisely in terms of degrees $`d_{+,m}`$ and $`d_{,m}`$, say in cases when we have vanishing of $`h^0`$ or $`h^1`$, proceeding similarly as in (\**), we can prove the theorem.
With this in mind, we write
$$\begin{array}{cc}& \underset{EW_d(C),d0}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ \hfill =& \underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g1}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ & +\underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g}{}\frac{q^{h^0(C,L_{+,2n+})+h^0(C,L_{,2n+})}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ \hfill =& \underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g1}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ & +\underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g,0d_{,m}2g2}{}\frac{q^{h^0(C,L_{+,2n+})+h^0(C,L_{,2n+})}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ \hfill +& \underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g,d_{,m}>2g2}{}\frac{q^{\chi (C,E_m)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ & +\underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g,d_{,m}<0}{}\frac{q^{\chi (C,L_{+,m})}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s,\hfill \end{array}$$
by a simple observation on cohomology groups. Set now
$$\begin{array}{cc}\hfill I:=& \underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g,d_{,m}>2g2}{}\frac{q^{\chi (C,E_m)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s,\hfill \\ \hfill II:=& \underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g,d_{,m}<0}{}\frac{q^{\chi (C,L_{+,m})}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s,\hfill \\ \hfill III:=& \underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g,0d_{,m}2g2}{}\frac{q^{h^0(C,L_{+,2n+})+h^0(C,L_{,2n+})}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s,\hfill \\ \hfill IV:=& \underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g1}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s.\hfill \end{array}$$
Then clearly,
(i) $`{\displaystyle \underset{EW_d(C),d0}{}}{\displaystyle \frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}}\left(q^{d(E)}\right)^s=I+II+III+IV.`$
(ii) By a similar calculation as in (\**) for the case $`_E\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s`$, we see that $`I`$ and $`II`$ are convergent provided $`\mathrm{Re}(s)>1`$ and are rational.
Therefore, it suffices to deal with $`III`$ and $`IV`$. We study $`IV`$ first. For this, note that $`d_{+,m}+d_{,m}=m`$ or $`m+1`$, we have then
$$\begin{array}{cc}\hfill IV=& \underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g1,d_{+,m}2g2}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ & +\underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g1,d_{+,m}>2g2}{}\frac{q^{h^0(C,L_{+,m})+h^0(C,L_{,m})}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ \hfill =& \underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g1,d_{+,m}2g2}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ & +\underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g1,d_{+,m}>2g2,0d_{,m}2g1}{}\frac{q^{h^0(C,L_{+,m})+h^0(C,L_{,m})}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ & +\underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g1,d_{+,m}>2g2,d_{,m}<0}{}\frac{q^{\chi (C,L_{+,m})}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s.\hfill \end{array}$$
Now set
$$\begin{array}{cc}\hfill IV_a:=& \underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g1,d_{+,m}2g2}{}\frac{q^{h^0(C,E)}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ \hfill IV_b:=& \underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g1,d_{+,m}>2g2,0d_{,m}2g1}{}\frac{q^{h^0(C,L_{+,m})+h^0(C,L_{,m})}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s\hfill \\ \hfill IV_c:=& \underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g1,d_{+,m}>2g2,d_{,m}<0}{}\frac{q^{\chi (C,L_{+,m})}}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s.\hfill \end{array}$$
Then
(a) $`IV=IV_a+IV_b+IV_c`$;
(b) $`IV_a`$ and $`IV_b`$ consist of only finitely many terms and hence are polynomials;
(c) $`IV_c`$ may be calculated similarly as in the case $`_E\frac{1}{\mathrm{\#}\mathrm{Aut}(E)}\left(q^{d(E)}\right)^s`$.
Hence we get the convergence provided $`\mathrm{Re}(s)>1`$ and the rationality for $`IV`$.
Finally, let us deal with $`III`$. This seems to be a rather delicate work: the summation is an infinite one; while previous trick cannot be applied since we cannot calculate $`h^0(C,L_{,m})`$ precisely. However, we are quite lucky: In any case, now we have $`h^0(C,E_m)=h^0(C,L_{+,m})+h^0(C,L_{,m})`$ and $`h^0(C,L_{+,m})=\chi (C,L_{+,m})`$. Hence,
$$III=\underset{EE_{}A^n,E_{}W_{}(C),=0,1,n2g,0d_{,m}2g2}{}\frac{q^{h^0(C,L_{,m})+d_{,m}}}{(q1)^2}\left(q^{2n+}\right)^s.$$
Set now for $`=0,1`$,
$$A_{}(m):=\underset{EE_{}A^n,E_{}W_{}(C),n2g,0d_{,m}2g2}{}\frac{q^{h^0(C,L_{,m})+d_{,m}}}{(q1)^2}.$$
Then we have
$$A_{}(m)=\underset{L\mathrm{Pic}^d(C)(𝐅_q),0d2g2}{}\frac{q^{h^0(C,L)}+d(L)}{(q1)^2}:=A(C)$$
which is independent of $``$ and $`m`$. Therefore
$$III=A(C)\times \underset{=0,1,n2g}{}\left(q^{2n+}\right)^s,$$
which is easily seen to be convergent if $`\mathrm{Re}(s)>1`$ and rational as well. In this way, we complete the proof of the Theorem 3.
In fact, a functional equation also holds for these new zeta functions. But we will leave it to the reader. Clearly, from (the proof of) this theorem, we may conclude that zeta functions associated to non-stable bundles, and hence the so-called total zeta function associated to all bundles, both semi-stable and non-stable, are interesting but not as canonical as that for semi-stable bundles. Moreover, note that total zeta functions are closed related to automorphic $`L`$-functions. Hence, it would be of great interests if one does similar decompositions (as what we have done here) for them.
We end this paper by pointing out that this work is a continuation of our small note \[We1\], in which a different kind of restricted zeta function is formulated, and that all these works are motivated by our new non-abelian zeta functions for number fields \[We2\].
REFERENCES
\[A\] E. Artin, Collected Papers, edited by S. Lang and J. Tata, Addison-Wesley (1965)
\[AB\] M.F. Atiyah & R. Bott, The Yang-Mills equations over Riemann surfaces, Phil. Trans. R. Soc. Lond. A. 308, 523-615 (1982)
\[DR\] U.V. Desale & S. Ramanan, Poincaré polynomials of the variety of stable bundles, Maeh. Ann 216, 233-244 (1975)
\[HN\] G. Harder & M.S. Narasimhan, On the cohomology groups of moduli spaces of vector bundles over curves, Math Ann. 212, (1975) 215-248
\[Mu\] D. Mumford, Geometric Invariant Theory, Springer-Verlag, Berlin (1965)
\[Se\] J.P. Serre, Trees, Springer-Verlag, Berlin (1980)
\[W\] A. Weil, Sur les courbes algébriques et les variétés qui s’en déduisent, Herman, Paris (1948)
\[We1\] L. Weng, A result on zeta functions for curves over finite fields, preprint, Nagoya University, 2000
\[We2\] L. Weng, Riemann-Roch Theorem, Stability, and New Zeta Functions for Number Fields, preprint, Nagoya, 2000
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# A Rosetta Stone for Quantum Mechanics with an Introduction to Quantum Computation Version 1.5
## Part I Preamble <br>
### 1. <br>Introduction
These lecture notes were written for the American Mathematical Society (AMS) Short Course on Quantum Computation held 17-18 January 2000 in conjunction with the Annual Meeting of the AMS in Washington, DC in January 2000. The notes are intended for readers with some mathematical background but with little or no exposure to quantum mechanics. The purpose of these notes is to provide such readers with enough material in quantum mechanics and quantum computation to begin reading the vast literature on quantum computation, quantum cryptography, and quantum information theory.
The paper was written in an informal style. Whenever possible, each new topic was begun with the introduction of the underlying motivating intuitions, and then followed by an explanation of the accompanying mathematical finery. Hopefully, once having grasped the basic intuitions, the reader will find that the remaining material easily follows.
Since this paper is intended for a diverse audience, it was written at varying levels of difficulty and sophistication, from the very elementary to the more advanced. A large number of examples have been included. An index and table of contents are provided for those readers who prefer to “pick and choose.” Hopefully, this paper will provide something of interest for everyone.
Because of space limitations, these notes are, of necessity, far from a complete overview of quantum mechanics. For example, only finite dimensional Hilbert spaces are considered, thereby avoiding the many pathologies that always arise when dealing with infinite dimensional objects. Many important experiments that are traditionally part of the standard fare in quantum mechanics texts (such as for example, the Stern-Gerlach experiment, Young’s two slit experiment, the Aspect experiment) have not been mentioned in this paper. We leave it to the reader to decide if these notes have achieved their objective.
The final version of this paper together with all the other lecture notes of the AMS Short Course on Quantum Computation will be published as a book in the AMS PSAPM Series entitled “Quantum Computation.”
### 2. The classical world
#### 2.1. Introducing the Shannon bit.
Since one of the objectives of this paper is to discuss quantum information, we begin with a brief discussion of classical information.
The Shannon bit is so well known in our age of information that it needs little, if any, introduction. As we all know, the Shannon bit is like a very decisive individual. It is either 0 or 1, but by no means both at the same time. The Shannon bit has become so much a part of our every day lives that we take many of its properties for granted. For example, we take for granted that Shannon bits can be copied.
#### 2.2. <br>Polarized light: Part I. The classical perspective
Throughout this paper the quantum polarization states of light will be used to provide concrete illustrations of underlying quantum mechanical principles. So we also begin with a brief discussion of polarized light from the classical perspective.
Light waves in the vacuum are transverse electromagnetic (EM) waves with both electric and magnetic field vectors perpendicular to the direction of propagation and also to each other. (See figure 1.)
Figure 1. A linearly polarized electromagnetic wave.
If the electric field vector is always parallel to a fixed line, then the EM wave is said to be linearly polarized. If the electric field vector rotates about the direction of propagation forming a right-(left-)handed screw, it is said to be right (left) elliptically polarized. If the rotating electric field vector inscribes a circle, the EM wave is said to be right-or left-circularly polarized.
### 3. The quantum world
#### 3.1. Introducing the qubit – But what is a qubit?
Many of us may not be as familiar with the quantum bit of information, called a qubit. Unlike its sibling rival, the Shannon bit, the qubit can be both 0 and 1 at the same time. Moreover, unlike the Shannon bit, the qubit can not be duplicated<sup>1</sup><sup>1</sup>1This is a result of the no-cloning theorem of Wootters and Zurek. A proof of the no-cloning theorem is given in Section 10.8 of this paper.. As we shall see, qubits are like very slippery, irascible individuals, exceedingly difficult to deal with.
One example of a qubit is a spin $`\frac{1}{2}`$ particle which can be in a spin-up state $`|1`$ which we label as “$`1`$”, in a spin-down state $`|0`$ which we label as “$`0`$”, or in a superposition of these states, which we interpret as being both $`0`$ and $`1`$ at the same time. (The term “superposition” will be explained shortly.)
Another example of a qubit is the polarization state of a photon. A photon can be in a vertically polarized state $`|`$. We assign a label of “$`1`$” to this state. It can be in a horizontally polarized state $`|`$. We assign a label of “$`0`$” to this state. Or, it can be in a superposition of these states. In this case, we interpret its state as representing both $`0`$ and $`1`$ at the same time.
Anyone who has worn polarized sunglasses is familiar with the polarization states of light. Polarized sunglasses eliminate glare by letting through only vertically polarized light, while filtering out the horizontally polarized light. For that reason, they are often used to eliminate road glare, i.e., horizontally polarized light reflected from the road.
#### 3.2. Where do qubits live? – But what is a qubit?
But where do qubits live? They live in a Hilbert space $``$. By a Hilbert space, we mean:
A Hilbert space $``$ is a vector space over the complex numbers $``$ with a complex valued inner product
$$(,):\times $$
which is complete with respect to the norm
$$u=\sqrt{(u,u)}$$
induced by the inner product.
###### Remark 1.
By a complex valued inner product, we mean a map
$$(,):\times $$
from $`\times `$ into the complex numbers $``$ such that:
* $`(u,u)=0`$ if and only if $`u=0`$
* $`(u,v)=(v,u)^{}`$
* $`(u,v+w)=(u,v)+(u,w)`$
* $`(u,\lambda v)=\lambda (u,v)`$
where ‘’ denotes the complex conjugate.
###### Remark 2.
Please note that $`(\lambda u,v)=\lambda ^{}(u,v)`$.
#### 3.3. <br>A qubit is …
<sup>2</sup><sup>2</sup>2Barenco et al in define a qubit as a quantum system with a two dimensional Hilbert space, capable of existing in a superposition of Boolean states, and also capable of being entangled with the states of other qubits. Their more functional definition will take on more meaning as the reader progresses through this paper.
| A qubit is a quantum system $`𝒬`$ whose |
| --- |
| state lies in a two dimensional Hilbert space $``$. |
## Part II An Introduction to Quantum Mechanics <br>
### 4. The beginnings of quantum mechanics
#### 4.1. A Rosetta stone for Dirac notation: Part I. Bras, kets, and bra-(c)-kets
The elements of a Hilbert space $``$ will be called ket vectors, state kets, or simply kets. They will be denoted as:
$$|label$$
where ‘$`label`$’ denotes some label.
Let $`^{}`$ denote the Hilbert space of all Hilbert space morphisms of $``$ into the Hilbert space of all complex numbers $``$, i.e.,
$$^{}=Hom_{}(,)\text{.}$$
The elements of $`^{}`$ will be called bra vectors, state bras, or simply bras. They will be denoted as:
$$label|$$
where once again ‘$`label`$’ denotes some label.
Also please note that the complex number
$$label_1|\left(|label_2\right)$$
will simply be denoted by
$$label_1label_2$$
and will be called the bra-(c)-ket product of the bra $`label_1|`$ and the ket $`|label_2`$.
There is a monomorphism (which is an isomorphism if the underlying Hilbert space is finite dimensional)
$$\stackrel{}{}^{}$$
defined by
$$|label(|label,)$$
The bra $`(|label,)`$ is denoted by $`label|`$.
Hence,
$$label_1label_2=(|label_1,|label_2)$$
###### Remark 3.
Please note that $`\left(\lambda |label\right)^{}=\lambda ^{}label|`$.
The tensor product<sup>3</sup><sup>3</sup>3Readers well versed in homological algebra will recognize this informal definition as a slightly disguised version of the more rigorous universal definition of the tensor product. For more details, please refer to , or any other standard reference on homological algebra. $`𝒦`$ of two Hilbert spaces $``$ and $`𝒦`$ is simply the “simplest” Hilbert space such that
1. $`\left(h_1+h_2\right)k=h_1k+h_2k`$, for all $`h_1`$, $`h_2`$ and for all $`k𝒦`$, and
2. $`h\left(k_1+k_2\right)=hk_1+hk_2`$ for all $`h`$ and for all $`k_1`$, $`k_2𝒦`$.
3. $`\lambda \left(hk\right)\left(\lambda h\right)k=h\left(\lambda k\right)`$ for all $`\lambda `$, $`h`$, $`k𝒦`$.
###### Remark 4.
Hence, $`|label=\sqrt{labellabel}`$ and $`label_1label_2=(|label_1,|label_2)`$ .
It follows that, if $`\{e_1,e_2,\mathrm{},e_m\}`$ and $`\{f_1,f_2,\mathrm{},f_n\}`$ are respectively bases of the Hilbert spaces $``$ and $`𝒦`$, then $`\{e_if_j1im\text{}1jn\}`$ is a basis of $`𝒦`$. Hence, the dimension of the Hilbert space $`𝒦`$ is the product of the dimensions of the Hilbert spaces $``$ and $`𝒦`$, i.e.,
$$Dim\left(𝒦\right)=Dim\left(\right)Dim\left(𝒦\right)\text{ .}$$
Finally, if $`|label_1`$ and $`|label_2`$ are kets respectively in Hilbert spaces $`_1`$ and $`_2`$, then their tensor product will be written in any one of the following three ways:
$$\begin{array}{c}|label_1|label_2\\ \\ |label_1|label_2\\ \\ |label_1,label_2\end{array}$$
#### 4.2. Quantum mechanics: Part I. The state of a quantum system
The states of a quantum system $`𝒬`$ are represented by state kets in a Hilbert space $``$. Two kets $`|\alpha `$ and $`|\beta `$ represent the same state of a quantum system $`𝒬`$ if they differ by a non-zero multiplicative constant. In other words, $`|\alpha `$ and $`|\beta `$ represent the same quantum state $`𝒬`$ if there exists a non-zero $`\lambda `$ such that
$$|\alpha =\lambda |\beta $$
Hence, quantum states are simply elements of the manifold
$$/\mathrm{~}=P^{n1}$$
where $`n`$ denotes the dimension of $``$, and $`P^{n1}`$ denotes complex projective $`\left(n1\right)`$-space .
Since a quantum mechanical state is represented by a state ket up to a multiplicative constant, we will, unless stated otherwise, choose those kets $`|\alpha `$ which are of unit length, i.e., such that
$$\alpha \alpha =1|\alpha =1$$
##### 4.2.1. Polarized light: Part II. The quantum mechanical perspective
As an illustration of the above concepts, we consider the polarization states of a photon.
The polarization states of a photon are represented as state kets in a two dimensional Hilbert space $``$. One orthonormal basis of $``$ consists of the kets
$$|\text{ and }|$$
which represent respectively the quantum mechanical states of left- and right-circularly polarized photons. Another orthonormal basis consists of the kets
$$|\text{ and }|$$
representing respectively vertically and horizontally linearly polarized photons. And yet another orthonormal basis consists of the kets
$$|\text{ and }|$$
for linearly polarized photons at the angles $`\theta =\pi /4`$ and $`\theta =\pi /4`$ off the vertical, respectively.
These orthonormal bases are related as follows:
$`\{\begin{array}{ccc}|& =& \frac{1}{\sqrt{2}}\left(|+|\right)\\ & & \\ |& =& \frac{1}{\sqrt{2}}\left(||\right)\end{array}\{\begin{array}{ccc}|& =& \frac{1+i}{2}|+\frac{1i}{2}|\\ & & \\ |& =& \frac{1i}{2}|+\frac{1+i}{2}|\end{array}`$
$`\{\begin{array}{ccc}|& =& \frac{1}{\sqrt{2}}\left(|+|\right)\\ & & \\ |& =& \frac{1}{\sqrt{2}}\left(||\right)\end{array}\{\begin{array}{ccc}|& =& \frac{1}{\sqrt{2}}\left(|+|\right)\\ & & \\ |& =& \frac{i}{\sqrt{2}}\left(||\right)\end{array}`$
$`\{\begin{array}{ccc}|& =& \frac{1}{\sqrt{2}}\left(|i|\right)\\ & & \\ |& =& \frac{1}{\sqrt{2}}\left(|+i|\right)\end{array}\{\begin{array}{ccc}|& =& \frac{1i}{2}|+\frac{1+i}{2}|\\ & & \\ |& =& \frac{1+i}{2}|+\frac{1i}{2}|\end{array}`$
The bracket products of the various polarization kets are given in the table below:
| | $`|`$ | $`|`$ | $`|`$ | $`|`$ | $`|`$ | $`|`$ |
| --- | --- | --- | --- | --- | --- | --- |
| $`|`$ | $`1`$ | $`0`$ | $`\frac{1}{\sqrt{2}}`$ | $`\frac{1}{\sqrt{2}}`$ | $`\frac{1}{\sqrt{2}}`$ | $`\frac{1}{\sqrt{2}}`$ |
| $`|`$ | $`0`$ | $`1`$ | $`\frac{1}{\sqrt{2}}`$ | $`\frac{1}{\sqrt{2}}`$ | $`\frac{i}{\sqrt{2}}`$ | $`\frac{i}{\sqrt{2}}`$ |
| $`|`$ | $`\frac{1}{\sqrt{2}}`$ | $`\frac{1}{\sqrt{2}}`$ | $`1`$ | $`0`$ | $`\frac{1i}{2}`$ | $`\frac{1+i}{2}`$ |
| $`|`$ | $`\frac{1}{\sqrt{2}}`$ | $`\frac{1}{\sqrt{2}}`$ | $`0`$ | $`1`$ | $`\frac{1+i}{2}`$ | $`\frac{1i}{2}`$ |
| $`|`$ | $`\frac{1}{\sqrt{2}}`$ | $`\frac{i}{\sqrt{2}}`$ | $`\frac{1+i}{2}`$ | $`\frac{1i}{2}`$ | $`1`$ | $`0`$ |
| $`|`$ | $`\frac{1}{\sqrt{2}}`$ | $`\frac{i}{\sqrt{2}}`$ | $`\frac{1i}{2}`$ | $`\frac{1+i}{2}`$ | $`0`$ | $`1`$ |
In terms of the basis $`\{|,|\}`$ and the dual basis $`\{|,\left|\right\}`$, these kets and bras can be written as matrices as indicated below:
$`\{\begin{array}{ccccccc}|& =& \left(\begin{array}{cc}1& 0\end{array}\right),& & |& =& \left(\begin{array}{c}1\\ 0\end{array}\right)\\ & & & & & & \\ |& =& \left(\begin{array}{cc}0& 1\end{array}\right),& & |& =& \left(\begin{array}{c}0\\ 1\end{array}\right)\end{array}`$
$`\{\begin{array}{ccccccc}|& =& \frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\end{array}\right),& & |& =& \frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 1\end{array}\right)\\ & & & & & & \\ |& =& \frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\end{array}\right),& & |& =& \frac{1}{\sqrt{2}}\left(\begin{array}{c}\hfill 1\\ \hfill 1\end{array}\right)\end{array}`$
$`\{\begin{array}{ccccccc}|& =& \frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& i\end{array}\right),& & |& =& \frac{1}{\sqrt{2}}\left(\begin{array}{c}\hfill 1\\ \hfill i\end{array}\right)\\ & & & & & & \\ |& =& \frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& i\end{array}\right),& & |& =& \frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ i\end{array}\right)\end{array}`$
In this basis, for example, the tensor product $`|`$ is
$$|=(\begin{array}{c}\frac{1}{\sqrt{2}}\\ \frac{1}{\sqrt{2}}\end{array})(\begin{array}{c}\hfill \frac{1}{\sqrt{2}}\\ \hfill \frac{i}{\sqrt{2}}\end{array})=\frac{1}{2}(\begin{array}{c}\hfill 1\\ \hfill i\\ \hfill 1\\ \hfill i\end{array})$$
and the projection operator $`||`$ is:
$$||=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ i\end{array}\right)\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& i\end{array}\right)=\frac{1}{2}\left(\begin{array}{cc}\hfill 1& \hfill i\\ \hfill i& \hfill 1\end{array}\right)$$
#### 4.3. A Rosetta stone for Dirac notation: Part II. Operators
An (linear) operator or transformation $`𝒪`$ on a ket space $``$ is a Hilbert space morphism of $``$ into $``$, i.e., is an element of
$$Hom_{}(,)$$
The adjoint $`𝒪^{}`$ of an operator $`𝒪`$ is that operator such that
$$(𝒪^{}|label_1,|label_2)=(|label_1,𝒪|label_2)$$
for all kets $`|label_1`$ and $`|label_2`$.
In like manner, an (linear) operator or transformation on a bra space $`^{}`$ is an element of
$$Hom_{}(^{},^{})$$
Moreover, each operator $`𝒪`$ on $``$ can be identified with an operator, also denoted by $`𝒪`$, on $`^{}`$ defined by
$$label_1|label_1|𝒪$$
where $`label_1|𝒪`$ is the bra defined by
$$\left(label_1|𝒪\right)\left(|label_2\right)=label_1|\left(𝒪|label_2\right)$$
(This is sometimes called Dirac’s associativity law.) Hence, the expression
$$label_1\left|𝒪\right|label_2$$
is unambiguous.
###### Remark 5.
Please note that
$$\left(𝒪|label\right)^{}=label|𝒪^{}$$
#### 4.4. Quantum mechanics: Part II. Observables
In quantum mechanics, an observable is simply a Hermitian (also called self-adjoint) operator on a Hilbert space $``$, i.e., an operator $`𝒪`$ such that
$$𝒪^{}=𝒪\text{ .}$$
An eigenvalue $`a`$ of an operator $`A`$ is a complex number for which there is a ket $`|label`$ such that
$$A|label=a|label\text{ .}$$
The ket $`|label`$ is called an eigenket of $`A`$ corresponding to the eigenvalue $`a`$.
An important theorem about observables is given below:
###### Theorem 1.
The eigenvalues $`a_i`$ of an observable $`A`$ are all real numbers. Moreover, the eigenkets for distinct eigenvalues of an observable are orthogonal.
###### Definition 1.
An eigenvalue is degenerate if there are at least two linearly independent eigenkets for that eigenvalue. Otherwise, it is non-degenerate .
If all the eigenvalues $`a_i`$ of an observable $`A`$ are nondegenerate, then we can and do label the eigenkets of $`A`$ with the corresponding eigenvalues $`a_i`$. Thus, we can write:
$$A|a_i=a_i|a_i$$
for each eigenvalue $`a_i`$.
In this paper, unless stated otherwise, we assume that the eigenvalues of observables are non-degenerate.
One notable exception to the above convention is the measurement operator
$$|a_ia_i|$$
for the eigenvalue $`a_i`$, which is the outer product of ket $`|a_i`$ with its adjoint the bra $`a_i|`$, where we have assumed that $`|a_i`$ (and hence, $`a_i|`$) is of unit length. It has two eigenvalues $`0`$ and $`1`$. $`1`$ is a nondegenerate eigenvalue with eigenket $`|a_i`$. $`0`$ is a degenerate eigenvalue with corresponding eigenkets $`\left\{|a_j\right\}_{ji}`$ .
An observable $`A`$ is said to be complete if its eigenkets $`|a_i`$ form a basis of the Hilbert space $``$. Since by convention all the eigenkets are chosen to be of unit length, it follows that the eigenkets of a complete nondegenerate observable $`A`$ form an orthonormal basis of the underlying Hilbert space.
Moreover, given a complete nondegenerate observable $`A`$, every ket $`|\psi `$ in $``$ can be written as:
$$|\psi =\underset{i}{}|a_ia_i\psi $$
Thus, for a complete nondegenerate observable $`A`$, we have the following operator equation which expresses the completeness of $`A`$,
$$\underset{i}{}|a_ia_i|=1$$
In this notation, we also have
$$A=\underset{i}{}a_i|a_ia_i|\text{ ,}$$
where once again we have assumed that $`|a_i`$ and $`a_i|`$ are of unit length for all $`i`$.
###### Example 1.
The Pauli spin matrices
$$\sigma _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\text{}\sigma _2=\left(\begin{array}{cc}\hfill 0& \hfill i\\ \hfill i& \hfill 0\end{array}\right)\text{}\sigma _3=\left(\begin{array}{cc}\hfill 1& \hfill 0\\ \hfill 0& \hfill 1\end{array}\right)$$
are examples of observables that frequently appear in quantum mechanics and quantum computation. Their eigenvalues and eigenkets are given in the following table:
| Pauli Matrices | Eigenvalue/Eigenket |
| --- | --- |
| $`\sigma _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ | +1 $`\frac{|0+|1}{\sqrt{2}}=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 1\end{array}\right)`$ -1 $`\frac{|0|1}{\sqrt{2}}=\frac{1}{\sqrt{2}}\left(\begin{array}{c}\hfill 1\\ \hfill 1\end{array}\right)`$ |
| $`\sigma _2=\left(\begin{array}{cc}\hfill 0& \hfill i\\ \hfill i& \hfill 0\end{array}\right)`$ | +1 $`\frac{|0+i|1}{\sqrt{2}}=\frac{1}{\sqrt{2}}\left(\begin{array}{c}\hfill 1\\ \hfill i\end{array}\right)`$ -1 $`\frac{|0i|1}{\sqrt{2}}=\frac{1}{\sqrt{2}}\left(\begin{array}{c}\hfill 1\\ \hfill i\end{array}\right)`$ |
| $`\sigma _3=\left(\begin{array}{cc}\hfill 1& \hfill 0\\ \hfill 0& \hfill 1\end{array}\right)`$ | +1 $`|0=\left(\begin{array}{c}\hfill 1\\ \hfill 0\end{array}\right)`$ -1 $`|1=\left(\begin{array}{c}\hfill 0\\ \hfill 1\end{array}\right)`$ |
#### 4.5. Quantum mechanics: Part III. Quantum measurement – General principles
In this section, $`A`$ will denote a complete nondegenerate observable with eigenvalues $`a_i`$ and eigenkets $`|a_i`$. We will, on occasion, refer to $`\left\{|a_i\right\}`$ as the frame (or the basis) of the observable $`A`$.
According to quantum measurement theory, the measurement of an observable $`A`$ of a quantum system $`𝒬`$ in the state $`|\psi `$ produces the eigenvalue $`a_i`$ as the measured result with probability
$$Prob\left(\text{Value }a_i\text{ is\hspace{1em}observed}\right)=a_i\psi ^2\text{ ,}$$
and forces the state of the quantum system $`𝒬`$ into the state of the corresponding eigenket $`|a_i`$.
Since quantum measurement is such a hotly debated topic among physicists, we (in self-defense) quote P.A.M. Dirac:
> “A measurement always causes the (quantum mechanical) system to jump into an eigenstate of the dynamical variable that is being measured.”
Thus, the result of the above mentioned measurement of observable $`A`$ of a quantum system $`𝒬`$ which is in the state $`|\psi `$ before the measurement can be diagrammatically represented as follows:
$$\overline{)|\psi =_i|a_ia_i\psi \begin{array}{c}\text{First}\\ \text{Meas. of }A\\ \\ Prob=a_j\psi ^2\end{array}a_j|a_j|a_j\begin{array}{c}\text{Second}\\ \text{Meas. of }A\\ \\ Prob=1\end{array}|a_j}$$
Please note that the measured value is the eigenvalue $`a_j`$ with probability $`a_j\psi ^2`$ . If the same measurement is repeated on the quantum system $`𝒬`$ after the first measurement, then the result of the second measurement is no longer stochastic. It produces the previous measured value $`a_j`$ and the state of $`𝒬`$ remains the same, i.e., $`|a_j`$ .
The observable
$$|a_ia_i|$$
is frequently called a selective measurement operator (or a filtration) for $`a_i`$. As mentioned earlier, it has two eigenvalues $`0`$ and $`1`$. $`1`$ is a nondegenerate eigenvalue with eigenket $`|a_j`$, and $`0`$ is a degenerate eigenvalue with eigenkets $`\left\{|a_j\right\}_{ji}`$.
Thus,
$`|\psi \begin{array}{c}\text{ Meas. of }|a_ia_i|\\ \\ Prob=a_i\psi ^2\end{array}1|a_i=|a_i`$ ,
but for $`ji`$,
$$\overline{)|\psi \begin{array}{c}\text{Meas. of }|a_ia_i|\\ \\ Prob=a_j\psi ^2\end{array}0|a_j=0}$$
The above description of quantum measurement is not the most general possible. For the more advanced quantum measurement theory of probabilistic operator valued measures (POVMs) (a.k.a., positive operator valued measures), please refer to such books as for example and .
#### 4.6. Polarized light: Part III. Three examples of quantum measurement
We can now apply the above general principles of quantum measurement to polarized light. Three examples are given below:<sup>4</sup><sup>4</sup>4The last two examples can easily be verified experimentally with at most three pair of polarized sunglasses.
###### Example 2.
$$\overline{)\begin{array}{c}\text{Rt. Circularly}\\ \text{polarized photon}\\ \\ |=\frac{1}{\sqrt{2}}\left(|+i|\right)\end{array}\begin{array}{c}\text{Vertical}\\ \text{Polaroid}\\ \text{filter}\\ \text{}\\ \text{Measurement op.}\\ ||\end{array}\begin{array}{ccc}& & \text{Vertically}\\ & & \text{polarized}\\ Prob=\frac{1}{2}& & \text{photon}\\ & & |\\ & & \\ & & \\ & & 0\\ Prob=\frac{1}{2}& & \text{No photon}\\ & & \end{array}}$$
###### Example 3.
A vertically polarized filter followed by a horizontally polarized filter.
$`\begin{array}{c}\\ \text{photon}\\ \\ \alpha |+\beta |\\ \\ \text{Normalized so that}\\ \alpha ^2+\beta ^2=1\end{array}\begin{array}{c}\text{Vert.}\\ \text{polar.}\\ \text{filter}\\ \text{}\\ \\ ||\end{array}`$.$`\begin{array}{c}Prob=\alpha ^2\\ \\ \end{array}\begin{array}{c}\text{Vert.}\\ \text{polar.}\\ \text{photon}\\ \\ |\end{array}\begin{array}{c}\text{Horiz.}\\ \text{polar.}\\ \text{filter}\\ \text{}\\ \\ ||\end{array}\begin{array}{cc}& \text{No}\\ & \text{photon}\\ Prob=1& \\ & 0\\ & \\ & \\ .& \end{array}`$
###### Example 4.
But if we insert a diagonally polarized filter (by $`45^o`$ off the vertical) between the two polarized filters in the above example, we have:
$$\overline{)\begin{array}{c}\\ \\ \text{}\\ \\ ||\end{array}\begin{array}{c}\alpha ^2\\ \\ \end{array}|=\frac{1}{\sqrt{2}}\left(|+|\right)\begin{array}{c}\\ \\ \text{}\\ \\ ||\end{array}\begin{array}{c}\frac{1}{2}\\ \\ \end{array}|=\frac{1}{\sqrt{2}}\left(|+|\right)\begin{array}{c}\\ \\ \text{}\\ \\ ||\end{array}\begin{array}{c}\frac{1}{2}\\ \\ \end{array}|}$$
where the input to the first filter is $`\alpha |+\beta |`$.
#### 4.7. A Rosetta stone for Dirac notation: Part III. Expected values
The average value (expected value) of a measurement of an observable $`A`$ on a state $`|\alpha `$ is:
$$A=\alpha \left|A\right|\alpha $$
For, since
$$\underset{i}{}|a_ia_i|=1\text{ ,}$$
we have
$$A=\alpha \left|A\right|\alpha =\alpha |\left(\underset{i}{}|a_ia_i|\right)A\left(\underset{j}{}|a_ja_j|\right)|\alpha =\underset{i,j}{}\alpha a_ia_i\left|A\right|a_ja_j\alpha $$
But on the other hand,
$$a_i\left|A\right|a_j=a_ja_ia_j=a_i\delta _{ij}$$
Thus,
$$A=\underset{i}{}\alpha a_ia_ia_i\alpha =\underset{i}{}a_ia_i\alpha ^2$$
Hence, we have the standard expected value formula,
$$A=\underset{i}{}a_iProb\left(\text{Observing }a_i\text{ on input }|\alpha \right)\text{ .}$$
#### 4.8. Quantum Mechanics: Part IV. The Heisenberg uncertainty principle
There is, surprisingly enough, a limitation of what we can observe in the quantum world.
From classical probability theory, we know that one yardstick of uncertainty is the standard deviation, which measures the average fluctuation about the mean. Thus, the uncertainty involved in the measurement of a quantum observable $`A`$ is defined as the standard deviation of the observed eigenvalues. This standard deviation is given by the expression
$$Uncertainty(A)=\sqrt{\left(\mathrm{}A\right)^2}$$
where
$$\mathrm{}A=AA$$
Two observables $`A`$ and $`B`$ are said to be compatible if they commute, i.e., if
$$AB=BA\text{.}$$
Otherwise, they are said to be incompatible.
Let $`[A,B]`$, called the commutator of $`A`$ and $`B`$, denote the expression
$$[A,B]=ABBA$$
In this notation, two operators $`A`$ and $`B`$ are compatible if and only if $`[A,B]=0`$.
The following principle is one expression of how quantum mechanics places limits on what can be observed:
Heisenberg’s Uncertainty Principle<sup>5</sup><sup>5</sup>5We have assumed units have been chosen such that $`\mathrm{}=1`$.
$$\left(\mathrm{}A\right)^2\left(\mathrm{}B\right)^2\frac{1}{4}\left|[A,B]\right|^2$$
Thus, if $`A`$ and $`B`$ are incompatible, i.e., do not commute, then, by measuring $`A`$ more precisely, we are forced to measure $`B`$ less precisely, and vice versa. We can not simultaneously measure both $`A`$ and $`B`$ to unlimited precision. Measurement of $`A`$ somehow has an impact on the measurement of $`B`$, and vice versa.
#### 4.9. Quantum mechanics: Part V. Dynamics of closed quantum systems: Unitary transformations, the Hamiltonian, and Schrödinger’s equation
An operator $`U`$ on a Hilbert space $``$ is unitary if
$$U^{}=U^1\text{ .}$$
Unitary operators are of central importance in quantum mechanics for many reasons. We list below only two:
* Closed quantum mechanical systems transform only via unitary transformations
* Unitary transformations preserve quantum probabilities
Let $`|\psi (t)`$ denote the state as a function of time $`t`$ of a closed quantum mechanical system $`𝒬`$ . Then the dynamical behavior of the state of $`𝒬`$ is determined by the Schrödinger equation
$$i\mathrm{}\frac{}{t}|\psi (t)=H|\psi (t)\text{ ,}$$
where $`\mathrm{}`$ denotes Planck’s constant divided by $`2\pi `$, and where $`H`$ denotes an observable of $`𝒬`$ called the Hamiltonian. The Hamiltonian is the quantum mechanical analog of the Hamiltonian of classical mechanics. In classical physics, the Hamiltonian is the total energy of the system.
#### 4.10. <br>The mathematical perspective
From the mathematical perspective, Schrödinger’s equation is written as:
$$\frac{}{t}U(t)=\frac{i}{\mathrm{}}H(t)U(t)\text{ ,}$$
where
$$|\psi (t)=U|\psi (0)\text{ ,}$$
and where $`\frac{i}{\mathrm{}}H(t)`$ is a skew-Hermitian operator lying in the Lie algebra of the unitary group. The solution is given by a multiplicative integral, called the path-ordered integral,
$$U(t)=_\stackrel{}{t}\text{}_{0}e^{\frac{i}{\mathrm{}}H(t)dt},$$
which is taken over the path $`\frac{i}{\mathrm{}}H(t)`$ in the Lie algebra of the unitary group. The path-ordered integral is given by:
$`{}_{\stackrel{}{t}}{}^{}\text{}_{0}^{}e^{\frac{i}{\mathrm{}}H(t)dt}`$ $`=\underset{n\mathrm{}}{lim}{\displaystyle \underset{k=n}{\overset{0}{}}}e^{\frac{i}{\mathrm{}}H(k\frac{t}{n})\frac{t}{n}}`$
$`=\underset{n\mathrm{}}{lim}\left[e^{\frac{i}{\mathrm{}}H(n\frac{t}{n})}e^{\frac{i}{\mathrm{}}H((n1)\frac{t}{n})}\mathrm{}e^{\frac{i}{\mathrm{}}H(1\frac{t}{n})}e^{\frac{i}{\mathrm{}}H(0\frac{t}{n})}\right]`$
###### Remark 6.
The standard notation for the above path-ordered integral is
$$𝐏\mathrm{exp}\left(\frac{i}{\mathrm{}}\underset{0}{\overset{t}{}}H(t)𝑑t\right)$$
If the Hamiltonian $`H(t)=H`$ is independent of time, then all matrices commute and the above path-ordered integral simplifies to
$${}_{\stackrel{}{t}}{}^{}\text{}_{0}^{}e^{\frac{i}{\mathrm{}}Hdt}=e^{_0^t\frac{i}{\mathrm{}}Hdt}=e^{\frac{i}{\mathrm{}}Ht}$$
Thus, in this case, $`U(t)`$ is nothing more than a one parameter subgroup of the unitary group.
### 5. The Density Operator
#### 5.1. Introducing the density operator
John von Neumann suggested yet another way of representing the state of a quantum system.
Let $`|\psi `$ be a unit length ket (i.e., $`\psi \psi =1`$) in the Hilbert space $``$ representing the state of a quantum system<sup>6</sup><sup>6</sup>6Please recall that each of the kets in the set $`\left\{\lambda |\psi \lambda \text{}\lambda 0\right\}`$ represent the same state of a quantum system. Hence, we can always (and usually do) represent the state of a quantum system as a unit normal ket, i.e., as a ket such that $`\psi \psi =1`$ .. The density operator $`\rho `$ associated with the state ket $`|\psi `$ is defined as the outer product of the ket $`|\psi `$ (which can be thought of as a column vector) with the bra $`\psi |`$ (which can be thought of as a row vector), i.e.,
$$\rho =|\psi \psi |$$
The density operator formalism has a number of advantages over the ket state formalism. One advantage is that the density operator can also be used to represent hybrid quantum/classical states, i.e., states which are a classical statistical mixture of quantum states. Such hybrid states may also be thought of as quantum states for which we have incomplete information.
For example, consider a quantum system which is in the states (each of unit length)
$$|\psi _1,|\psi _2,\mathrm{},|\psi _n$$
with probabilities
$$p_1,p_2,\mathrm{},p_n$$
respectively, where
$$p_1+p_2+\mathrm{}+p_n=1$$
(Please note that the states $`|\psi _1,|\psi _2,\mathrm{},|\psi _n`$ need not be orthogonal.) Then the density operator representation of this state is defined as
$$\rho =p_1|\psi _1\psi _1\left|+p_2\right|\psi _2\psi _2\left|+\mathrm{}+p_n\right|\psi _n\psi _n|$$
If a density operator $`\rho `$ can be written in the form
$$\rho =|\psi \psi |\text{ ,}$$
it is said to represent a pure ensemble . Otherwise, it is said to represent a mixed ensemble .
#### 5.2. Properties of density operators
It can be shown that all density operators are positive semi-definite Hermitian operators of trace 1, and vice versa. As a result, we have the following crisp mathematical definition:
###### Definition 2.
An linear operator on a Hilbert space $``$ is a density operator if it is a positive semi-definite Hermitian operator of trace 1.
It can be shown that a density operator represents a pure ensemble if and only if $`\rho ^2=\rho `$, or equivalently, if and only if $`Trace(\rho ^2)=1`$. For all ensembles, both pure and mixed, $`Trace(\rho ^2)1`$.
From standard theorems in linear algebra, we know that, for every density operator $`\rho `$, there exists a unitary matrix $`U`$ which diagonalizes $`\rho `$, i.e., such that $`U\rho U^{}`$ is a diagonal matrix. The diagonal entries in this matrix are, of course, the eigenvalues of $`\rho `$. These are non-negative real numbers which all sum to 1.
Finally, if we let $`𝒟`$ denote the set of all density operators for a Hilbert space $``$, then $`i𝒟`$ is a convex subset of the Lie algebra of the unitary group associated with $``$.
#### 5.3. Quantum measurement in terms of density operators
Let $`\left\{a_i\right\}`$ denote the set of distinct eigenvalues $`a_i`$ of an observable $`A`$. Let $`P_{a_i}`$ denote the projection operator that projects the underlying Hilbert space onto the eigenspace determined by the eigenvalue $`a_i`$. For example, if $`a_i`$ is a non-degenerate eigenvalue, then
$$P_{a_i}=|a_ia_i|$$
Finally, let $`𝒬`$ be a quantum system with state given by the density operator $`\rho `$.
If the quantum system $`𝒬`$ is measured with respect to the observable $`A`$, then with probability
$$p_i=Trace\left(P_{a_i}\rho \right)$$
the resulting measured eigenvalue is $`a_i`$, and the resulting state of $`𝒬`$ is given by the density operator
$$\rho _i=\frac{P_{a_i}\rho P_{a_i}}{Trace\left(P_{a_i}\rho \right)}\text{ .}$$
Moreover, for an observable $`A`$, the averaged observed eigenvalue expressed in terms of the density operator is:
$$A=trace(\rho A)$$
Thus, we have extended the definition of $`A`$ so that it applies to mixed as well as pure ensembles, i.e., generalized the following formula to mixed ensembles:
$$A=\psi A\psi =trace\left(|\psi \psi |A\right)=trace(\rho A)\text{ .}$$
#### 5.4. Some examples of density operators
For example, consider the following mixed ensemble of the polarization state of a photon:
###### Example 5.
| Ket | $`\underset{}{\overset{}{|}}`$ | $`|`$ |
| --- | --- | --- |
| Prob. | $`\underset{}{\overset{}{\frac{3}{4}}}`$ | $`\frac{1}{4}`$ |
In terms of the basis $`|`$, $`|`$ of the two dimensional Hilbert space $``$, the density operator $`\rho `$ of the above mixed ensemble can be written as:
$$\begin{array}{ccc}\rho & =& \frac{3}{4}||+\frac{1}{4}||\hfill \\ & & \\ & =& \frac{3}{4}\left(\begin{array}{c}1\\ 0\end{array}\right)\left(\begin{array}{cc}1& 0\end{array}\right)+\frac{1}{4}\left(\begin{array}{c}1/\sqrt{2}\\ 1/\sqrt{2}\end{array}\right)\left(\begin{array}{cc}1/\sqrt{2}& 1/\sqrt{2}\end{array}\right)\hfill \\ & & \\ & =& \frac{3}{4}\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)+\frac{1}{8}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)=\left(\begin{array}{cc}\underset{}{\overset{}{\frac{7}{8}}}& \frac{1}{8}\\ \underset{}{\overset{}{\frac{1}{8}}}& \frac{1}{8}\end{array}\right)\hfill \end{array}$$
###### Example 6.
The following two preparations produce mixed ensembles with the same density operator:
| Ket | $`\underset{}{\overset{}{|}}`$ | $`|`$ |
| --- | --- | --- |
| Prob. | $`\underset{}{\overset{}{\frac{1}{2}}}`$ | $`\frac{1}{2}`$ |
and
| Ket | $`\underset{}{\overset{}{|}}`$ | $`|`$ |
| --- | --- | --- |
| Prob. | $`\underset{}{\overset{}{\frac{1}{2}}}`$ | $`\frac{1}{2}`$ |
For, for the left preparation, we have
$$\begin{array}{ccc}\rho & =& \frac{1}{2}||+\frac{1}{2}||\hfill \\ & & \\ & =& \frac{1}{2}\left(\begin{array}{c}1\\ 0\end{array}\right)\left(\begin{array}{cc}1& 0\end{array}\right)+\frac{1}{2}\left(\begin{array}{c}0\\ 1\end{array}\right)\left(\begin{array}{cc}0& 1\end{array}\right)\hfill \\ & & \\ & =& \frac{1}{2}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\hfill \end{array}$$
And for the right preparation, we have
$$\begin{array}{ccc}\rho & =& \frac{1}{2}||+\frac{1}{2}||\hfill \\ & & \\ & =& \frac{1}{2}\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 1\end{array}\right)\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\end{array}\right)+\frac{1}{2}\frac{1}{\sqrt{2}}\left(\begin{array}{c}\hfill 1\\ \hfill 1\end{array}\right)\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\end{array}\right)\hfill \\ & & \\ & =& \frac{1}{4}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)+\frac{1}{4}\left(\begin{array}{cc}\hfill 1& \hfill 1\\ \hfill 1& \hfill 1\end{array}\right)=\frac{1}{2}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\hfill \end{array}$$
There is no way of physically distinquishing the above two mixed ensembles which were prepared in two entirely different ways. For the density operator represents all that can be known about the state of the quantum system.
#### 5.5. The partial trace of a linear operator
In order to deal with a quantum system composed of many quantum subsystems, we need to define the partial trace.
Let
$$𝒪:Hom_{}(,)$$
be a linear operator on the Hilbert space $``$.
Since Hilbert spaces are free algebraic objects, it follows from standard results in abstract algebra<sup>7</sup><sup>7</sup>7See for example . that
$$Hom_{}(,)^{}\text{ ,}$$
where we recall that
$$^{}=Hom_{}(,)\text{ .}$$
Hence, such an operator $`𝒪`$ can be written in the form
$$𝒪=\underset{\alpha }{}a_\alpha |h_\alpha k_\alpha |\text{ ,}$$
where the kets $`|h_\alpha `$ lie in $``$ and the bras $`k_\alpha |`$ lie in $`^{}`$.
Thus, the standard trace of a linear operator
$$Trace:Hom_{}(,)$$
is nothing more than a contraction, i.e.,
$$Trace(𝒪)=\underset{\alpha }{}a_\alpha k_\alpha h_\alpha \text{ ,}$$
i.e., a replacement of each outer product $`|h_\alpha k_\alpha |`$ by the corresponding bracket $`k_\alpha h_\alpha `$.
We can generalize the $`Trace`$ as follows:
Let $``$ now be the tensor product of Hilbert spaces $`_1`$, $`_2`$, $`\mathrm{}`$ ,$`_n`$, i.e.,
$$=\underset{j=1}{\overset{n}{}}_j\text{ .}$$
Then it follows once again from standard results in abstract algebra that
$$Hom_{}(,)\underset{j=1}{\overset{n}{}}\left(_j_j^{}\right)\text{ .}$$
Hence, the operator $`𝒪`$ can be written in the form
$$𝒪=\underset{\alpha }{}a_\alpha \underset{j=1}{\overset{n}{}}|h_{\alpha ,j}k_{\alpha ,j}|\text{ ,}$$
where, for each $`j`$, the kets $`|h_{\alpha ,j}`$ lie in $`_j`$ and the bras $`k_{\alpha ,j}|`$ lie in $`_j^{}`$ for all $`\alpha `$.
Next we note that for every subset $``$ of the set of indices $`𝒥=\{1,2,\mathrm{},n\}`$, we can define the partial trace over $``$, written
$$Trace_{}:Hom_{}(\underset{j𝒥}{}_j,\underset{j𝒥}{}_j)Hom_{}(\underset{j𝒥}{}_j,\underset{j𝒥}{}_j)$$
as the contraction on the indices $``$, i.e.,
$$Trace_{}\left(𝒪\right)=\underset{\alpha }{}a_\alpha \left(\underset{j}{}k_{\alpha ,j}h_{\alpha ,j}\right)\underset{j𝒥}{}|h_{\alpha ,j}k_{\alpha ,j}|\text{ .}$$
For example, let $`_1`$ and $`_0`$ be two dimensional Hilbert spaces with selected orthonormal bases $`\{|0_1,|1_1\}`$ and $`\{|0_0,|1_0\}`$, respectively. Thus, $`\{|0_10_0,|0_11_0,|1_10_0,|1_11_0\}`$ is an orthonormal basis of $`=_1_0`$ .
Let $`\rho Hom_{}(,)`$ be the operator
$`\rho `$ $`=\left({\displaystyle \frac{|0_10_0|1_11_0}{\sqrt{2}}}\right)\left({\displaystyle \frac{0_10_0|1_11_0|}{\sqrt{2}}}\right)`$
$`={\displaystyle \frac{1}{2}}\left(|0_10_00_10_0||0_10_01_11_0||1_11_00_10_0|+|1_11_01_11_0|\right)`$
which in terms of the basis $`\{|0_10_0,|0_11_0,|1_10_0,|1_11_0\}`$ can be written as the matrix
$$\rho =\frac{1}{2}\left(\begin{array}{cccc}\hfill 1& \hfill 0& \hfill 0& \hfill 1\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 1& \hfill 0& \hfill 0& \hfill 1\end{array}\right)\text{ ,}$$
where the rows and columns are listed in the order $`|0_10_0`$, $`|0_11_0`$, $`|1_10_0`$, $`|1_11_0`$
The partial trace $`Trace_0`$ with respect to $`=\left\{0\right\}`$ of $`\rho `$ is
$`\rho _1`$ $`=Trace_0\left(\rho \right)`$
$`={\displaystyle \frac{1}{2}}Trace_0\left(|0_10_00_10_0||0_10_01_11_0||1_11_00_10_0|+|1_11_01_11_0|\right)`$
$`={\displaystyle \frac{1}{2}}\left(0_00_0|0_10_1\left|1_00_0\right|0_11_1\left|0_01_0\right|1_10_1\left|+1_01_0\right|1_11_1|\right)`$
$`={\displaystyle \frac{1}{2}}\left(|0_10_1||0_11_1|\right)`$
which in terms of the basis $`\{|0_1,|1_1\}`$ becomes
$$\rho _1=Trace_0(\rho )=\frac{1}{2}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\text{ ,}$$
where the rows and columns are listed in the order $`|0_1`$, $`|1_1`$ .
#### 5.6. Multipartite quantum systems
One advantage density operators have over kets is that they provide us with a means for dealing with multipartite quantum systems.
###### Definition 3.
Let $`𝒬_1`$, $`𝒬_2`$, $`\mathrm{}`$ , $`𝒬_n`$ be quantum systems with underlying Hilbert spaces $`_1`$, $`_2`$, $`\mathrm{}`$ , $`_n`$, respectively. The global quantum system $`𝒬`$ consisting of the quantum systems $`𝒬_1`$, $`𝒬_2`$, $`\mathrm{}`$ , $`𝒬_n`$ is called a multipartite quantum system. Each of the quantum systems $`𝒬_j`$ ($`j=1,2,\mathrm{},n`$) is called a constituent “part” of $`𝒬`$ . The underlying Hilbert space $``$ of $`𝒬`$ is the tensor product of the Hilbert spaces of the constituent “parts,” i.e.,
$$=\underset{j=1}{\overset{n}{}}_j\text{ .}$$
If the density operator $`\rho `$ is the state of a multipartite quantum system $`𝒬`$, then the state of each constituent “part” $`𝒬_j`$ is the density operator $`\rho _j`$ given by the partial trace
$$\rho _j=Trace_{𝒥\left\{j\right\}}\left(\rho \right)\text{ ,}$$
where $`𝒥=\{1,2,\mathrm{},n\}`$ is the set of indices.
Obviously, much more can be said about the states of multipartite systems and their constituent parts. However, we will forego that discussion until after we have had an opportunity introduce the concepts of quantum entanglement and von Neumann entropy.
#### 5.7. Quantum dynamics in density operator formalism
Under a unitary transformation $`U`$, a density operator $`\rho `$ transforms according to the rubric:
$$\rho U\rho U^{}$$
Moreover, in terms of the density operator, Schrödinger’s equation<sup>8</sup><sup>8</sup>8Schrödinger’s equation determines the dynamics of closed quantum systems. However, non-closed quantum systems are also of importance in quantum computation and quantum information theory. See for example the Schumacher’s work on superoperators, e.g., . becomes:
$$i\mathrm{}\frac{\rho }{t}=[H,\rho ]\text{ ,}$$
where $`[H,\rho ]`$ denotes the commutator of $`H`$ and $`\rho `$, i.e.,
$$[H,\rho ]=H\rho \rho H$$
#### 5.8. <br>The mathematical perspective
From the mathematical perspective, one works with $`i\rho `$ instead of $`\rho `$ because $`i\rho `$ lies in the Lie algebra of the unitary group. Thus, the density operator transforms under a unitary transformation $`U`$ according to the rubric:
$$i\rho Ad_U(i\rho )\text{ ,}$$
where $`Ad_U`$ denotes the big adjoint representation.
From the mathematical perspective, Schrödinger’s equation is in this case more informatively written as:
$$\frac{(i\rho )}{t}=\frac{1}{\mathrm{}}ad_{iH}(i\rho )\text{ ,}$$
where $`ad_{\frac{i}{\mathrm{}}H}`$ denotes the little adjoint representation . Thus, the solution to the above form of Schrödinger’s equation is given by the path ordered integral:
$$\rho =\left({}_{\stackrel{}{t}}{}^{}\text{}_{0}^{}e^{\frac{1}{\mathrm{}}\left(ad_{iH(t)}\right)dt}\right)\rho _0$$
where $`\rho _0`$ denotes the density operator at time $`t=0`$.
### 6. <br>The Heisenberg model of quantum mechanics
Consider a computing device with inputs and outputs for which we have no knowledge of the internal workings of the device. We are allowed to probe the device with inputs and observe the corresponding outputs. But we are given no information as to how the device performs its calculation. We call such a device a blackbox computing device.
For such blackboxes, we say that two theoretical models for blackboxes are equivalent provided both predict the same input/output behavior. We may prefer one model over the other for various reasons, such as simplicity, aesthetics, or whatever meets our fancy. But the basic fact is that each of the two equivalent models is just as “correct” as the other.
In like manner, two theoretical models of the quantum world are said to be equivalent if they both predict the same results in regard to quantum measurements.
Up to this point, we have been describing the Schrödinger model of quantum mechanics. However, shortly after Schrödinger proposed his model for the quantum world, called the Schrödinger picture, Heisenberg proposed yet another, called the Heisenberg picture. Both models were later proven to be equivalent.
In the Heisenberg picture, state kets remain stationary with time, but observables move with time. While state kets, and hence density operators, remain fixed with respect to time, the observables $`A`$ change dynamically as:
$$AU^{}AU$$
under a unitary transformation $`U=U(t)`$, where the unitary transformation is determined by the equation
$$i\mathrm{}\frac{U}{t}=HU$$
It follows that the equation of motion of observables is according to the following equation
$$i\mathrm{}\frac{A}{t}=[A,H]$$
One advantage the Heisenberg picture has over the Schrödinger picture is that the equations appearing in it are similar to those found in classical mechanics.
In summary, we have the following table which contrasts the two pictures:
| | Schrödinger Picture | Heisenberg Picture |
| --- | --- | --- |
| State ket | Moving $`\underset{}{\overset{}{|\psi _0|\psi =U|\psi _0}}`$ | Stationary $`\underset{}{\overset{}{|\psi _0}}`$ |
| Density Operator | Moving $`\underset{}{\overset{}{\rho _0\rho =U\rho _0U^{}=Ad_U\left(\rho _0\right)}}`$ | Stationary $`\underset{}{\overset{}{\rho _0}}`$ |
| Observable | Stationary $`\underset{}{\overset{}{A_0}}`$ | Moving $`\underset{}{\overset{}{A_0A=U^{}A_0U=Ad_U^{}\left(A_0\right)}}`$ |
| Observable Eigenvalues | Stationary $`\underset{}{\overset{}{a_j}}`$ | Stationary $`\underset{}{\overset{}{a_j}}`$ |
| Observable Frame | Stationary $`\underset{}{\overset{}{A_0=_ja_j|a_j_0a_j|_0}}`$ | Moving $`\underset{}{\overset{}{A_0=_ja_j|a_j_0a_j|_0}}`$ $``$ $`\underset{}{\overset{}{A_t=_ja_j|a_j_ta_j|_t}}`$ where $`\underset{}{\overset{}{|a_j_t=U^{}|a_j_0}}`$ |
| Dynamical Equations | $`\underset{}{\overset{}{i\mathrm{}\frac{U}{t}=H^{(S)}U}}`$ $`\underset{}{\overset{}{i\mathrm{}\frac{}{t}|\psi =H^{(S)}|\psi }}`$ | $`\underset{}{\overset{}{i\mathrm{}\frac{U}{t}=H^{(H)}U}}`$ $`\underset{}{\overset{}{i\mathrm{}\frac{A}{t}=[A,H^{(H)}]}}`$ |
| Measurement | Measurement of observable $`A_0`$ produces eigenvalue $`a_j`$ with probability $`\underset{}{\overset{}{\left|\left(a_j|_0\right)|\psi \right|^2=\left|\left(a_j|_0\right)|\psi \right|^2}}`$ | Measurement of observable $`A`$ produces eigenvalue $`a_j`$ with probability $`\underset{}{\overset{}{\left|\left(a_j|_t\right)|\psi _0\right|^2=\left|\left(a_j|_0\right)|\psi \right|^2}}`$ |
where
$$H^{(H)}=U^{}H^{(S)}U$$
It follows that the Schrödinger Hamiltonian $`H^{(S)}`$ and the Heisenberg Hamiltonian are related as follows:
$$\frac{H^{(S)}}{t}=U\frac{H^{(H)}}{t}U^{}\text{,}$$
where terms containing $`\frac{U}{t}`$ and $`\frac{U^{}}{t}`$ have cancelled out as a result of the Schrödinger equation.
We should also mention that the Schrödinger and Heisenberg pictures can be transformed into one another via the mappings:
| $`SH`$ | $`HS`$ |
| --- | --- |
| $`\underset{}{\overset{}{|\psi ^{(S)}|\psi ^{(H)}=U^{}|\psi ^{(S)}}}`$ | $`|\psi ^{(H)}|\psi ^{(S)}=U|\psi ^{(H)}`$ |
| $`\underset{}{\overset{}{\rho ^{(S)}\rho ^{(H)}=U^{}\rho ^{(S)}U}}`$ | $`\rho ^{(H)}\rho ^{(S)}=U\rho ^{(H)}U^{}`$ |
| $`\underset{}{\overset{}{A^{(S)}A^{(H)}=U^{}A^{(S)}U}}`$ | $`A^{(H)}A^{(S)}=UA^{(H)}U^{}`$ |
| $`\underset{}{\overset{}{A^{(S)}A^{(H)}=U^{}A^{(S)}U}}`$ | $`A^{(H)}A^{(S)}=UA^{(H)}U^{}`$ |
Obviously, much more could be said on this topic.
For quantum computation from the perspective of the Heisenberg model, please refer to the work of Deutsch and Hayden, and also to Gottesman’s “study of the ancient Hittites” :-) .
### 7. Quantum entanglement
#### 7.1. The juxtaposition of two quantum systems
Let $`𝒬_1`$ and $`𝒬_2`$ be two quantum systems that have been separately prepared respectively in states $`|\psi _1`$ and $`|\psi _2`$, and that then have been united without interacting. Because $`𝒬_1`$ and $`𝒬_2`$ have been separately prepared without interacting, their states $`|\psi _1`$ and $`|\psi _2`$ respectively lie in distinct Hilbert spaces $`_1`$ and $`_2`$. Moreover, because of the way in which $`𝒬_1`$ and $`𝒬_2`$ have been prepared, all physical predictions relating to one of these quantum systems do not depend in any way whatsoever on the other quantum system.
The global quantum system $`𝒬`$ consisting of the two quantum systems $`𝒬_1`$ and $`𝒬_2`$ as prepared above is called a juxtaposition of the quantum systems $`𝒬_1`$ and $`𝒬_2`$. The state of the global quantum system $`𝒬`$ is the tensor product of the states $`|\psi _1`$ and $`|\psi _2`$. In other words, the state of $`𝒬`$ is:
$$|\psi _1|\psi _2_1_2$$
#### 7.2. An example: An $`n`$-qubit register $`𝒬`$ consisting of the juxtaposition of $`n`$ qubits.
Let $``$ be a two dimensional Hilbert space, and let $`\{|0,|1\}`$ denote an arbitrarily selected orthonormal basis<sup>9</sup><sup>9</sup>9We obviously have chosen to label the basis elements in a suggestive way.. Let $`_{n1}`$, $`_{n2}`$, $`\mathrm{}`$ ,$`_0`$ be distinct Hilbert spaces, each isomorphic to $``$, with the obvious induced orthonormal bases
$$\{|0_{n1},|1_{n1}\},\{|0_{n2},|1_{n2}\},\mathrm{},\{|0_0,|1_0\}$$
respectively.
Consider $`n`$ qubits $`𝒬_{n1}`$, $`𝒬_{n2}`$, $`\mathrm{}`$ , $`𝒬_0`$ separately prepared in the states
$`{\displaystyle \frac{1}{\sqrt{2}}}\left(|0_{n1}+|1_{n1}\right)\text{}{\displaystyle \frac{1}{\sqrt{2}}}\left(|0_{n2}+|1_{n2}\right)\text{}\mathrm{}\text{ , }{\displaystyle \frac{1}{\sqrt{2}}}\left(|0_0+|1_0\right)\text{,}`$
respectively. Let $`𝒬`$ denote the global system consisting of the separately prepared (without interacting) qubits $`𝒬_{n1}`$, $`𝒬_{n2}`$, $`\mathrm{}`$ , $`𝒬_0`$. Then the state $`|\psi `$ of $`𝒬`$ is:
$`|\psi `$ $`={\displaystyle \frac{1}{\sqrt{2}}}\left(|0_{n1}+|1_{n1}\right){\displaystyle \frac{1}{\sqrt{2}}}\left(|0_{n2}+|1_{n2}\right)\mathrm{}{\displaystyle \frac{1}{\sqrt{2}}}\left(|0_0+|1_0\right)`$
$`=\left({\displaystyle \frac{1}{\sqrt{2}}}\right)^n\left(|0_{n1}0_{n2}\mathrm{}0_10_0+|0_{n1}0_{n2}\mathrm{}0_11_0+\mathrm{}+|1_{n1}1_{n2}\mathrm{}1_11_0\right)`$
which lies in the Hilbert space
$$=_{n1}_{n2}\mathrm{}_0.$$
We will usually omit subscripts whenever they can easily be inferred from context.
Thus, the global system $`𝒬`$ consisting of the $`n`$ qubits $`𝒬_{n1}`$, $`𝒬_{n2}`$,$`\mathrm{}`$, $`𝒬_0`$ is in the state
$$|\psi =\left(\frac{1}{\sqrt{2}}\right)^n\left(|00\mathrm{}00+|00\mathrm{}01+\mathrm{}+|11\mathrm{}11\right)\underset{0}{\overset{n1}{}}$$
The reader should note that the $`n`$-qubit register $`𝒬`$ is a superposition of kets with labels consisting of all the binary n-tuples. If each binary n-tuple $`b_{n1}b_{n2}\mathrm{}b_0`$ is identified with the integer
$$b_{n1}2^{n1}+b_{n2}2^{n2}+\mathrm{}+b_02^0\text{ ,}$$
i.e., if we interpret each binary n-tuple as the radix 2 representation of an integer, then we can rewrite the state as
$$|\psi =\left(\frac{1}{\sqrt{2}}\right)^n\left(|0+|1+|2+\mathrm{}+|2^n1\right)\text{.}$$
In other words, this n-qubit register contains all the integers from $`0`$ to $`2^n1`$ in superposition. But most importantly, it contains all the integers $`0`$ to $`2^n1`$ simultaneously!
This is an example of the massive parallelism that is possible within quantum computation. However, there is a downside. If we observe (measure) the register, then all the massive parallelism disappears. On measurement, the quantum world selects for us one and only one of the $`2^n`$ integers. The probability of observing any particular one of the integers is $`\left|\left(1/\sqrt{2}\right)^n\right|^2=(\frac{1}{2})^n`$. The selection of which integer is observed is unfortunately not made by us, but by the quantum world.
Thus, harnessing the massive parallelism of quantum mechanics is no easy task! As we will see, a more subtle approach is required.
#### 7.3. An example of the dynamic behavior of a 2-qubit register
We now consider the previous $`n`$-qubit register for $`n=2`$. In terms of the bases described in the previous section, we have:
$$\{\begin{array}{ccccccc}|0& =& |00& =& \left(\begin{array}{c}1\\ 0\end{array}\right)\left(\begin{array}{c}1\\ 0\end{array}\right)& =& \left(\begin{array}{c}1\\ 0\\ 0\\ 0\end{array}\right)\\ & & & & & & \\ |1& =& |01& =& \left(\begin{array}{c}1\\ 0\end{array}\right)\left(\begin{array}{c}0\\ 1\end{array}\right)& =& \left(\begin{array}{c}0\\ 1\\ 0\\ 0\end{array}\right)\\ & & & =& & & \\ |2& =& |10& =& \left(\begin{array}{c}0\\ 1\end{array}\right)\left(\begin{array}{c}1\\ 0\end{array}\right)& =& \left(\begin{array}{c}0\\ 0\\ 1\\ 0\end{array}\right)\\ & & & & & & \\ |3& =& |11& =& \left(\begin{array}{c}0\\ 1\end{array}\right)\left(\begin{array}{c}0\\ 1\end{array}\right)& =& \left(\begin{array}{c}0\\ 0\\ 0\\ 1\end{array}\right)\end{array}$$
Let us assume that the initial state $`|\psi _{t=0}`$ of our 2-qubit register is
$$|\psi _{t=0}=\left(\frac{|0|1}{\sqrt{2}}\right)|0=\frac{1}{\sqrt{2}}\left(|00|10\right)=\frac{1}{\sqrt{2}}\left(|0|2\right)=\frac{1}{\sqrt{2}}\left(\begin{array}{c}\hfill 1\\ \hfill 0\\ \hfill 1\\ \hfill 0\end{array}\right)$$
Let us also assume that from time $`t=0`$ to time $`t=1`$ the dynamical behavior of the above 2-qubit register is determined by a constant Hamiltonian $`H`$, which when written in terms of the basis $`\{|00,|01,|10,|11\}=\{|0,|1,|2,|3\}`$ is given by
$$H=\frac{\pi \mathrm{}}{2}\left(\begin{array}{cccc}\hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 1\\ \hfill 0& \hfill 0& \hfill 1& \hfill 1\end{array}\right)\text{ ,}$$
where the rows and the columns are listed in the order $`|00`$, $`|01`$, $`|10`$, $`|11`$, i.e., in the order $`|0`$, $`|1`$, $`|2`$, $`|3`$.
Then, as a consequence of Schrödinger’s equation, the Hamiltonian $`H`$ determines a unitary transformation
$`U_{CNOT}`$ $`=_\stackrel{}{t}\text{}_{0}e^{\frac{i}{\mathrm{}}Hdt}=e^{_0^1\frac{i}{\mathrm{}}Hdt}=e^{\frac{i}{\mathrm{}}H}`$
$`=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right)=|00|+|11|+|23|+|32|`$
which moves the 2-qubit register from the initial state $`|\psi _{t=0}`$ at time $`t=0`$ to $`|\psi _{t=1}=U_{CNOT}|\psi _{t=0}`$ at time $`t=1`$. Then
$`|\psi _{t=1}`$ $`=U_{CNOT}|\psi _{t=0}=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right){\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\hfill 1\\ \hfill 0\\ \hfill 1\\ \hfill 0\end{array}\right)`$
$`={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\hfill 1\\ \hfill 0\\ \hfill 0\\ \hfill 1\end{array}\right)={\displaystyle \frac{1}{\sqrt{2}}}\left(|00|11\right)={\displaystyle \frac{1}{\sqrt{2}}}\left(|0|3\right)`$
The resulting state (called an EPR pair of qubits for reasons we shall later explain) can no longer be written as a tensor product of two states. Consequently, we no longer have the juxtaposition of two qubits.
Somehow, the resulting two qubits have in some sense “lost their separate identities.” Measurement of any one of the qubits immediately impacts the other.
For example, if we measure the 0-th qubit (i.e., the right-most qubit), the EPR state in some sense “jumps” to one of two possible states. Each of the two possibilities occurs with probability $`\frac{1}{2}`$, as indicated in the table below:
| $`\underset{}{\overset{}{\frac{1}{\sqrt{2}}\left(|0_10_0|1_11_0\right)}}`$ | |
| --- | --- |
| $`\begin{array}{c}\text{Meas.}\\ \text{0-th}\\ \text{Qubit}\end{array}`$ | |
| $`\stackrel{}{Prob=\frac{1}{2}}`$ $`\underset{}{|0_10_0}`$ | $`\stackrel{}{Prob=\frac{1}{2}}`$ $`\underset{}{|1_11_0}`$ |
Thus we see that a measurement of one of the qubits causes a change in the other.
#### 7.4. Definition of quantum entanglement
The above mentioned phenomenon is so unusual and so non-classical that it warrants a name.
###### Definition 4.
Let $`𝒬_1`$, $`𝒬_2`$, $`\mathrm{}`$ , $`𝒬_n`$ be quantum systems with underlying Hilbert spaces $`_1`$, $`_2`$, $`\mathrm{}`$ , $`_n`$, respectively. Then the global quantum system $`𝒬`$ consisting of the quantum systems $`𝒬_1`$, $`𝒬_2`$, $`\mathrm{}`$ , $`𝒬_n`$ is said to be entangled if its state $`|\psi =_{j=1}^n_j`$ can not be written in the form
$$|\psi =\underset{j=1}{\overset{n}{}}|\psi _j\text{ ,}$$
where each ket $`|\psi _j`$ lies in the Hilbert space $`_j`$ for, $`j=1,2,\mathrm{},n`$. We also say that such a state $`|\psi `$ is entangled.
Thus, the state
$$|\psi _{t=1}=\frac{1}{\sqrt{2}}\left(|00|11\right)$$
of the 2-qubit register of the previous section is entangled.
###### Remark 7.
In terms of density operator formalism, a pure ensemble $`\rho `$ is entangled if it can not be written in the form
$$\rho =\underset{j=1}{\overset{n}{}}\rho _j\text{ ,}$$
where the $`\rho _j`$’s denote density operators.
Please note that we have defined entanglement only for pure ensembles. For mixed ensembles, entanglement is not well understood<sup>10</sup><sup>10</sup>10Quantum entanglement is not even well understood for pure ensembles.. As a result, the “right” definition of entanglement of mixed ensembles is still unresolved. We give one definition below:
###### Definition 5.
A density operator $`\rho `$ on a Hilbert space $``$ is said to be entangled with respect to the Hilbert space decomposition
$$=\underset{j=1}{\overset{n}{}}_j$$
if it can not be written in the form
$$\rho =\underset{k=1}{\overset{\mathrm{}}{}}\lambda _k\left(\underset{j=1}{\overset{n}{}}\rho _{(j,k)}\right)\text{ ,}$$
for some positive integer $`\mathrm{}`$, where the $`\lambda _k`$’s are positive real numbers such that
$$\underset{k=1}{\overset{\mathrm{}}{}}\lambda _k=1\text{ .}$$
and where each $`\rho _{(j,k)}`$ is a density operator on the Hilbert space $`_{j\text{.}}`$
Readers interested in pursuing this topic further should refer to the works of Bennett , the Horodecki’s, Nielsen, Smolin, Wootters , and others, , , .
#### 7.5. Einstein, Podolsky, Rosen’s (EPR’s) grand challenge to quantum mechanics.
Albert Einstein was skeptical of quantum mechanics, so skeptical that he together with Podolsky and Rosen wrote a joint paper appearing in 1935 challenging the very foundations of quantum mechanics. Their paper hit the scientific community like a bombshell. For it delivered a direct frontal attack at the very heart and center of quantum mechanics.
At the core of their objection was quantum entanglement. Einstein and his colleagues had insightfully recognized the central importance of this quantum phenomenon.
Their argument centered around the fact that quantum mechanics violated either the principle of non-locality<sup>11</sup><sup>11</sup>11We will later explain the principle of non-locality. or the principle of reality<sup>12</sup><sup>12</sup>12For an explanation of the principle of reality as well as the principle of non-localty, please refer, for example, to , . . They argued that, as a result, quantum mechanics must be incomplete, and that quantum entanglement could be explained by missing hidden variables.
For many years, no one was able to conceive of an experiment that could determine which of the two theories, i.e., quantum mechanics or EPR’s hidden variable theory, was correct. In fact, many believed that the two theories were not distinguishable on physical grounds.
It was not until Bell developed his famous inequalities ,, , that a physical criterion was found to distinquish the two theories. Bell developed inequalities which, if violated, would clearly prove that quantum mechanics is correct, and hidden variable theories are not. Many experiments were performed. Each emphatically supported quantum mechanics, and clearly demonstrated the incorrectness of hidden variable theory. Quantum mechanics was the victor!
#### 7.6. Why did Einstein, Podolsky, Rosen (EPR) object?
But why did Einstein and his colleagues object so vehemently to quantum entanglement?
As a preamble to our answer to this question, we note that Einstein and his colleagues were convinced of the validity of the following two physical orinciples:
* The principle of local interactions , i.e., that all the known forces of nature are local interactions,
* The principle of non-locality, i.e., that spacelike separated regions of spacetime are physically independent of one another.
Their conviction in regard to principle 1) was based on the fact that all four known forces of nature, i.e., gravitational, electromagnetic, weak, and strong forces, are local interactions. By this we mean:
* They are mediated by another entity, e.g., graviton, photon, etc.
* They propagate no faster than the speed $`c`$ of light
* Their strength drops off with distance
Their conviction in regard to principle 2) was based on the following reasoning:
Two points in spacetime $`P_1=(x_1,y_1,z_1,t_1)`$ and $`P_2=(x_2,y_2,z_2,t_2)`$ are separated by a spacelike distance provided the distance between $`(x_1,y_1,z_1)`$ and $`(x_2,y_2,z_2)`$ is greater than $`c\left|t_2t_1\right|`$, i.e.,
$$Distance((x_1,y_1,z_1),(x_2,y_2,z_2))>c\left|t_2t_1\right|\text{ ,}$$
where $`c`$ denotes the speed of light. In other words, no signal can travel between points that are said to be separated by a spacelike distance unless the signal travels faster than the speed of light. But because of the basic principles of relativity, such superluminal communication is not possible.
Hence we have:
Spacelike separated regions of spacetime are physically independent. In other words, spacelike separated regions can not influence one another.
##### 7.6.1. EPR’s objection
We now are ready to explain why Einstein and his colleagues objected so vehemently to quantum entanglement. We explain Bohm’s simplified version of their argument.
Consider a two qubit quantum system that has been prepared by Alice <sup>13</sup><sup>13</sup>13Alice is a well known personality in quantum computation, quantum cryptography, and quantum information theory. in her laboratory in the state
$$|\psi =\frac{1}{\sqrt{2}}\left(|0_10_0|1_11_0\right)\text{ .}$$
After the preparation, she decides to keep qubit #1 in her laboratory, but enlists Captain James T. Kirk of the Starship Enterprise to transport qubit #0 to her friend Bob <sup>14</sup><sup>14</sup>14Bob is another well known personality in quantum computation, quantum cryptography, and quantum information theory. who is at some far removed distant part of the universe, such as at a Federation outpost orbiting about the double star Alpha Centauri in the constellation Centaurus.
After Captain Kirk has delivered qubit #0, Alice’s two qubits are now separated by a spacelike distance. Qubit #1 is located in her Earth based laboratory. Qubits #0 is located with Bob at a Federation outpost orbiting Alpha Centauri. But the two qubits are still entangled, even in spite of the fact that they are separated by a spacelike distance.
If Alice now measures qubit #1 (which is located in her Earth based laboratory), then the principles of quantum mechanics force her to conclude that instantly, without any time lapse, both qubits are “effected.” As a result of the measurement, both qubits will be either in the state $`|0_10_0`$ or the state $`|1_11_0`$, each possibility occurring with probability 1/2.
This is a non-local “interaction.” For,
* The “interaction” occurred without the presence of any force. It was not mediated by anything.
* The measurement produced an instantaneous change, which was certainly faster than the speed of light.
* The strength of the “effect” of the measurement did not drop off with distance.
No wonder Einstein was highly skeptical of quantum entanglement. Yet puzzlingly enough, since no information is exchanged by the process, the principles of general relativity are not violated. As a result, such an “effect” can not be used for superluminal communication.
For a more in-depth discussion of the EPR paradox and the foundations of quantum mechanics, the reader should refer to .
#### 7.7. <br>Quantum entanglement: The Lie group perspective
Many aspects of quantum entanglement can naturally be captured in terms of Lie groups and their Lie algebras.
Let
$$=_{n1}_{n2}\mathrm{}_0=_{\stackrel{}{n1}}_0_j$$
be a decomposition of a Hilbert space $``$ into the tensor product of the Hilbert spaces $`_{n1}`$, $`_{n2}`$, $`\mathrm{}`$ ,$`_0`$. Let $`𝕌=𝕌()`$, $`𝕌_{n1}=𝕌(_{n1})`$, $`𝕌_{n2}=𝕌(_{n2})`$, $`\mathrm{}`$,$`𝕌_0=𝕌(_0)`$, denote respectively the Lie groups of all unitary transformations on $``$, $`_{n1}`$, $`_{n2}`$, $`\mathrm{}`$ ,$`_0`$. Moreover, let $`𝐮=𝐮()`$, $`𝐮_{n1}=𝐮_{n1}(_{n1})`$, $`𝐮_{n2}=𝐮_{n2}(_{n2})`$, $`\mathrm{},𝐮_0=𝐮_0(_0)`$ denote the corresponding Lie algebras.
###### Definition 6.
The local subgroup $`𝕃=𝕃()`$ of $`𝕌=𝕌()`$ is defined as the subgroup
$$𝕃=𝕌_{n1}𝕌_{n2}\mathrm{}𝕌_0=_{\stackrel{}{n1}}_0𝕌_j\text{ .}$$
The elements of $`𝕃`$ are called local unitary transformations . Unitary transformations which are in $`𝕌`$ but not in $`𝕃`$ are called global unitary transformations. The corresponding lie algebra
$$\mathrm{}=𝐮_{n1}𝐮_{n2}\mathrm{}𝐮_0$$
is called the local Lie algebra, where ‘$``$’ denotes the Kronecker sum<sup>15</sup><sup>15</sup>15The Kronecker sum $`\mathrm{A}\mathrm{B}`$ is defined as
$$\mathrm{A}\mathrm{B}=\mathrm{A}1+1\mathrm{B}\text{ ,}$$
where $`1`$ denotes the identity transformation..
Local unitary transformations can not entangle quantum systems with respect to the above tensor product decomposition. However, global unitary transformations are those unitary transformations which can and often do produce interactions which entangle quantum systems. This leads to the following definition:
###### Definition 7.
Two states $`|\psi _1`$ and $`|\psi _2`$ in $``$ are said to be locally equivalent ( or, of the same entanglement type) , written
$$|\psi _1\underset{local}{}|\psi _2\text{ ,}$$
if there exists a local unitary transformation $`U𝕃`$ such that
$$U|\psi _1=|\psi _2\text{ .}$$
The equivalence classes of local equivalence $`\underset{local}{}`$ are called the entanglement classes of $``$. Two density operators $`\rho _1`$ and $`\rho _2`$, (and hence, the corresponding two skew Hermitian operators $`i\rho _1`$ and $`i\rho _2`$ lying in $`𝐮`$) are said to be locally equivalent ( or, of the same entanglement type), written
$$\rho _1\underset{local}{}\rho _2\text{ ,}$$
if there exists a local unitary transformation $`U𝕃`$ such that
$$Ad_U(\rho _1)=\rho _2\text{ ,}$$
where $`Ad_U`$ denotes the big adjoint representation, i.e., $`Ad_U(i\rho )=U(i\rho )U^{}`$. The equivalence classes under this relation are called entanglement classes of the Lie algebra $`𝐮()`$.
Thus, the entanglement classes of the Hilbert space $``$ are just the orbits of the group action of $`𝕃()`$ on $``$. In like manner, the entanglement classes of the Lie algebra $`𝐮()`$ are the orbits of the big adjoint action of $`𝕃()`$ on $`𝐮()`$. Two states are entangled in the same way if and only if they lie in the same entanglement class, i.e., the same orbit.
For example, let us assume that Alice and Bob collectively possess two qubits $`𝒬_{AB}`$ which are in the entangled state
$$|\psi _1=\frac{|0_B0_A+|1_B1_A}{\sqrt{2}}=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 0\\ 0\\ 1\end{array}\right)\text{ ,}$$
and moreover that Alice possesses qubit labeled $`A`$, but not the qubit labeled $`B`$, and that Bob holds qubit $`B`$, but not qubit $`A`$. Let us also assume that Alice and Bob are also separated by a spacelike distance. As a result, they can only apply local unitary transformations to the qubits that they possess.
Alice could, for example, apply the local unitary transformation
$$U_A=\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)=\left(\begin{array}{cccc}\hfill 0& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1\\ \hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0\end{array}\right)$$
to her qubit to move Alice’s and Bob’s qubits $`A`$ and $`B`$ respectively into the state
$$|\psi _2=\frac{|0_B1_A|1_B0_A}{\sqrt{2}}=\frac{1}{\sqrt{2}}\left(\begin{array}{c}\hfill 0\\ \hfill 1\\ \hfill 1\\ \hfill 0\end{array}\right)\text{ ,}$$
Bob also could accomplish the same by applying the local unitary transformation
$$U_B=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right)=\left(\begin{array}{cccc}\hfill 0& \hfill 1& \hfill 0& \hfill 0\\ \hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0\end{array}\right)$$
to his qubit.
By local unitary transformations, Alice and Bob can move the state of their two qubits to any other state within the same entanglement class. But with local unitary transformations, there is no way whatsoever that Alice and Bob can transform the two qubits into a state lying in a different entanglement class (i.e., a different orbit), such as
$$|\psi _3=|0_B0_A.$$
The only way Alice and Bob could transform the two qubits from state $`|\psi _1`$ to the state $`|\psi _3`$ is for Alice and Bob to come together, and make the two qubits interact with one another via a global unitary transformation such as
$$U_{AB}=\frac{1}{\sqrt{2}}\left(\begin{array}{cccc}\hfill 1& \hfill 0& \hfill 0& \hfill 1\\ \hfill 0& \hfill 1& \hfill 1& \hfill 0\\ \hfill 0& \hfill 1& \hfill 1& \hfill 0\\ \hfill 1& \hfill 0& \hfill 0& \hfill 1\end{array}\right)$$
The main objective of this approach to quantum entanglement is to determine when two states lie in the same orbit or in different orbits? In other words, what is needed is a complete set of invariants, i.e., invariants that completely specify all the orbits ( i.e., all the entanglement classes). We save this topic for another lecture.
At first it would seem that state kets are a much better vehicle than density operators for the study of quantum entanglement. After all, state kets are much simpler mathematical objects. So why should one deal with the additional mathematical luggage of density operators?
Actually, density operators have a number of advantages over state kets. The most obvious advantage is that density operators certainly have an upper hand over state kets when dealing with mixed ensembles. But their most important advantage is that the orbits of the adjoint action are actually manifolds, which have a very rich and pliable mathematical structure. Needless to say, this topic is beyond the scope of this paper.
###### Remark 8.
It should also be mentioned that the mathematical approach discussed in this section by no means captures every aspect of the physical phenomenon of quantum entanglement. The use of ancilla and of classical communication have not been considered. For an in-depth study of the relation between quantum entanglement and classical communication (including catalysis), please refer to the work of Jonathan, Nielson, and others.
In regard to describing the locality of unitary operations, we will later have need for a little less precision than that given above in the above definitions. So we give the following (unfortunately rather technical) definitions:
###### Definition 8.
Let $``$, $`_{n1}`$, $`_{n2}`$, $`\mathrm{}`$ ,$`_0`$ be as stated above. Let $`𝒫=\left\{B_\alpha \right\}`$ be a partition of the set of indices $`\{0,1,2,\mathrm{},n1\}`$, i.e., $`𝒫`$ is a collection of disjoint subsets $`B_\alpha `$ of $`\{0,1,2,\mathrm{},n1\}`$, called blocks, such that $`_\alpha B_\alpha =\{0,1,2,\mathrm{},n1\}`$. Then the $`𝒫`$-tensor product decomposition of $``$ is defined as
$$=\underset{B_\alpha 𝒫}{}_{B_\alpha }\text{ ,}$$
where
$$_{B_\alpha }=\underset{jB_\alpha }{}_j\text{ ,}$$
for each block $`B_\alpha `$ in $`𝒫`$. Also the subgroup of $`𝒫`$-local unitary transformations $`𝕃_𝒫()`$ is defined as the subgroup of local unitary transformations of $``$ corresponding to the $`𝒫`$-tensor decomposition of $``$.
We define the fineness of a partition $`𝒫`$, written $`fineness(𝒫)`$, as the maximum number of indices in a block of $`𝒫`$. We say that a unitary transformation $`U`$ of $``$ is sufficiently local if there exists a partition $`𝒫`$ with sufficiently small $`fineness(𝒫)`$ (e.g., $`fineness(𝒫)3`$) such that $`U𝕃_𝒫()`$.
###### Remark 9.
The above lack of precision is needed because there is no way to know what kind (if any) of quantum computing devices will be implemented in the future. Perhaps we will at some future date be able to construct quantum computing devices that locally manipulate more than 2 or 3 qubits at a time?
### 8. Entropy and quantum mechanics
#### 8.1. Classical entropy, i.e., Shannon Entropy
Let $`𝒮`$ be a probability distribution on a finite set $`\{s_1,s_2,\mathrm{},s_n\}`$ of elements called symbols given by
$$\text{Prob}\left(s_j\right)=p_j\text{ ,}$$
where $`_{j=1}^np_j=1`$. Let $`s`$ denote the random variable (i.e., finite memoryless stochastic source) that produces the value $`s_j`$ with probability $`p_j`$.
###### Definition 9.
The classical entropy (also called the Shannon entropy) $`\mathrm{H}(\mathrm{S})`$ of a probability distribution $`𝒮`$ (or of the source $`s`$) is defined as:
$$H(𝒮)=H(s)=\underset{j=1}{\overset{n}{}}p_j\mathrm{lg}(p_j)\text{ ,}$$
where ‘$`\mathrm{lg}`$’ denotes the $`\mathrm{log}`$ to the base 2 .
Classical entropy $`H(𝒮)`$ is a measure of the uncertainty inherent in the probability distribution $`𝒮`$. Or in other words, it is the measure of the uncertainty of an observer before the source $`s`$ “outputs” a symbol $`s_j`$.
One property of such classical stochastic sources we often take for granted is that the output symbols $`s_j`$ are completely distinguishable from one another. We will see that this is not necessarily the case in the strange world of the quantum.
#### 8.2. Quantum entropy, i.e., Von Neumann entropy
Let $`𝒬`$ be a quantum system with state given by the desity operator $`\rho `$.
Then there are many preparations
|ψ1
|ψ2
|ψnp1p2pnPreparationPreparationfragments
|ψ1
|ψ2
|ψnfragmentsp1fragmentsp2fragmentsp𝑛\overset{\text{{Preparation}}}{\begin{tabular}[c]{||c||c||c||c||}\hline\cr\hline\cr$\left|\psi_{1}\right\rangle$&$\left|\psi_{2}\right\rangle$&$\ \ldots\ $&$\left|\psi_{n}\right\rangle$\\
\hline\cr\hline\cr$p_{1}$&$p_{2}$&$\ldots$&$p_{n}$\\
\hline\cr\hline\cr\end{tabular}}
which will produce the same state $`\rho `$. These preparations are classical stochastic sources with classical entropy given by
$$H=p_j\mathrm{lg}(p_j)\text{ .}$$
Unfortunately, the classical entropy $`H`$ of a preparation does not necessarily reflect the uncertainty in the resulting state $`\rho `$. For two different preparations $`𝒫_1`$ and $`𝒫_2`$, having different entropies $`H\left(𝒫_1\right)`$ and $`H\left(𝒫_2\right)`$, can (and often do) produce the same state $`\rho `$. The problem is that the states of the preparation my not be completely physically distinguishable from one another. This happens when the states of the preparation are not orthogonal. (Please refer to the Heisenberg uncertainty principle.)
John von Neumann found that the true measure of quantum entropy can be defined as follows:
###### Definition 10.
Let $`𝒬`$ be a quantum system with state given by the density operator $`\rho `$. Then the quantum entropy (also called the von Neumann entropy) of $`𝒬`$, written $`S(𝒬)`$, is defined as
$$S(𝒬)=Trace\left(\rho \mathrm{lg}\rho \right)\text{ ,}$$
where ‘$`\mathrm{lg}\rho `$’ denotes the log to the base 2 of the operator $`\rho `$.
###### Remark 10.
The operator $`\mathrm{lg}\rho `$ exists and is an analytic map $`\rho \mathrm{lg}\rho `$ given by the power series
$$\mathrm{lg}\rho =\frac{1}{\mathrm{ln}2}\underset{n=1}{\overset{\mathrm{}}{}}(1)^{n+1}\frac{\left(\rho I\right)^n}{n}$$
provided that $`\rho `$ is sufficiently close to the identity operator $`I`$, i.e., provided
$$\rho I<1\text{ ,}$$
where
$$A=\underset{v}{sup}\frac{Av}{v}\text{ .}$$
It can be shown that this is the case for all positive definite Hermitian operators of trace $`1`$.
For Hermitian operators $`\rho `$ of trace $`1`$ which are not positive definite, but only positive semi-definite (i.e., which have a zero eigenvalue), the logarithm $`\mathrm{lg}(\rho )`$ does not exist. However, there exists a sequence $`\rho _1,\rho _2,\rho _3,\mathrm{}`$ of positive definite Hermitian operators of trace $`1`$ which converges to $`\rho `$, i.e., such that
$$\rho =\underset{k\mathrm{}}{lim}\rho _k$$
It can then be shown that the limit
$$\underset{k\mathrm{}}{lim}\rho _k\mathrm{lg}\rho _k$$
exists.
Hence, $`S(\rho )`$ is defined and exists for all density operators $`\rho `$.
Quantum entropy is a measure of the uncertainty at the quantum level. As we shall see, it is very different from the classical entropy that arises when a measurement is made.
One important feature of quantum entropy $`S(\rho )`$ is that it is invariant under the adjoint action of unitary transformations, i.e.,
$$S\left(Ad_U(\rho )\right)=S\left(U\rho U^{}\right)=S(\rho )\text{ .}$$
It follows that, for closed quantum systems, it is a dynamical invariant. As the state $`\rho `$ moves according to Schrödinger’s equation, the quantum entropy $`S(\rho )`$ of $`\rho `$ remains constant. It does not change unless measurement is made, or, as we shall see, unless we ignore part of the quantum system.
Because of unitary invariance, the quantum entropy can be most easily computed by first diagonalizing $`\rho `$ with a unitary transformation $`U`$, i.e.,
$$U\rho U^{}=\mathrm{}(\stackrel{}{\lambda })\text{ ,}$$
where $`\mathrm{}(\stackrel{}{\lambda })`$ denotes the diagonal matrix with diagonal $`\stackrel{}{\lambda }=(\lambda _1,\lambda _2,\mathrm{},\lambda _n)`$.
Once $`\rho `$ has been diagonalized , we have
$`S(\rho )`$ $`=Trace\left(\mathrm{}(\stackrel{}{\lambda })\mathrm{lg}\mathrm{}(\stackrel{}{\lambda })\right)`$
$`=Trace\left(\mathrm{}(\lambda _1\mathrm{lg}\lambda _1,\lambda _2\mathrm{lg}\lambda _2,\mathrm{},\lambda _n\mathrm{lg}\lambda _n)\right)`$
$`={\displaystyle \underset{j=1}{\overset{n}{}}}\lambda _j\mathrm{lg}\lambda _j\text{ ,}`$
where the $`\lambda _j`$’s are the eigenvalues of $`\rho `$, and where $`0\mathrm{lg}00`$.
Please note that, because $`\rho `$ is positive semi-definite Hermitian of trace $`1`$, all the eigenvalues of $`\rho `$ are non-negative real numbers such that
$$\underset{j=1}{\overset{n}{}}\lambda _j=1\text{ .}$$
As an immediate corollary we have that the quantum entropy of a pure ensemble must be zero, i.e.,
$`\rho `$ pure ensemble $`S(\rho )=0`$
There is no quantum uncertainty in a pure ensemble. However, as expected, there is quantum uncertainty in mixed ensembles.
#### 8.3. How is quantum entropy related to classical entropy?
But how is classical entropy $`H`$ related to quantum entropy $`S`$?
Let $`A`$ be an observable of the quantum system $`𝒬`$. Then a measurement of $`A`$ of $`𝒬`$ produces an eigenvalue $`a_i`$ with probability
$$p_i=Trace\left(P_{a_i}\rho \right)\text{ ,}$$
where $`P_{a_i}`$ denotes the projection operator for the eigenspace of the eigenvalue $`a_i`$. For example, if $`a_i`$ is a non-degenerate eigenvalue, then $`P_{a_i}=|a_ia_i|`$ .
In other words, measurement of $`A`$ of the quantum system $`𝒬`$ in state $`\rho `$ can be identified with a classical stochastic source with the eigenvalues $`a_i`$ as output symbols occurring with probability $`p_i`$. We denote this classical stochastic source simply by $`(\rho ,A)`$ .
The two entropies $`S(\rho )`$ and $`H(\rho ,A)`$ are by no means the same. One is a measure of quantum uncertainty before measurement, the other a measure of the classical uncertainty that results from measurement. The quantum entropy $`S(\rho )`$ is usually a lower bound for the classical entropy, i.e.,
$$S(\rho )H(\rho ,A)\text{ .}$$
If $`A`$ is a complete observable (hence, non-degenerate), and if $`A`$ is compatible with $`\rho `$, i.e., $`[\rho ,A]=0`$, then $`S(\rho )=H(\rho ,A)`$.
#### 8.4. <br>When a part is greater than the whole – Ignorance = uncertainty
Let $`𝒬`$ be a multipartite quantum system with constituent parts $`𝒬_{n1}`$, $`\mathrm{}`$ ,$`𝒬_1`$, $`𝒬_0`$, and let the density operator $`\rho `$ denote the state of $`𝒬`$. Then from section 5.6 of this paper we know that the state $`\rho _j`$ of each constituent “part” $`𝒬_j`$ is given by the partial trace over all degrees of freedom except $`𝒬_j`$, i.e., by
$$\rho _j=\underset{\begin{array}{c}0kn1\\ kj\end{array}}{Trace\left(\rho \right)}\text{ .}$$
By applying the above partial trace, we are focusing only on the quantum system $`𝒬_j`$, and literally ignoring the remaining constituent “parts” of $`𝒬`$. By taking the partial trace, we have done nothing physical to the quantum system. We have simply ignored parts of the quantum system.
What is surprising is that, by intentionally ignoring “part” of the quantum system, we can in some cases create more quantum uncertainty. This happens when the constituent “parts” of $`𝒬`$ are quantum entangled.
For example, let $`𝒬`$ denote the bipartite quantum system consisting of two qubits $`𝒬_1`$ and $`𝒬_0`$ in the entangled state
$$|\mathrm{\Psi }_𝒬=\frac{|0_10_0|1_11_0}{\sqrt{2}}\text{ .}$$
The corresponding density operator $`\rho _𝒬`$ is
$`\rho _𝒬`$ $`={\displaystyle \frac{1}{2}}\left(|0_10_00_10_0||0_10_01_11_0||1_11_00_10_0|+|1_11_01_11_0|\right)`$
$`={\displaystyle \frac{1}{2}}\left(\begin{array}{cccc}\hfill 1& \hfill 0& \hfill 0& \hfill 1\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 1& \hfill 0& \hfill 0& \hfill 1\end{array}\right)`$
Since $`\rho _𝒬`$ is a pure ensemble, there is no quantum uncertainty, i.e.,
$$S\left(\rho _𝒬\right)=0\text{ .}$$
Let us now focus on qubit #0 (i.e., $`𝒬_0`$). The resulting density operator $`\rho _0`$ for qubit #0 is obtained by tracing over $`𝒬_1`$, i.e.,
$$\rho _0=Trace_1\left(\rho _𝒬\right)=\frac{1}{2}\left(|00|+|11|\right)=\frac{1}{2}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\text{ .}$$
Hence, the quantum uncertainty of qubit #0 is
$$S(\rho _0)=1\text{ .}$$
Something most unusual, and non-classical, has happened. Simply by ignoring part of the quantum system, we have increased the quantum uncertainty. The quantum uncertainty of the constituent “part” $`𝒬_0`$ is greater than that of he whole quantum system $`𝒬`$. This is not possible in the classical world, i.e., not possible for Shannon entropy. (For more details, see .)
### 9. <br>There is much more to quantum mechanics
There is much more to quantum mechanics. For more in-depth overviews, there are many outstanding books. Among such books are , , , , , , , , , , , , , , , and many more.
## Part III Part of a Rosetta Stone for Quantum Computation <br>
### 10. The Beginnings of Quantum Computation - Elementary Quantum Computing Devices
We begin this section with some examples of quantum computing devices. By a quantum computing device<sup>16</sup><sup>16</sup>16Unfortunately, Physicists have “stolen” the akronym QCD. :-) we mean a unitary transformation $`U`$ that is the composition of finitely many sufficiently local unitary transformations, i.e.,
$$U=U_{n1}U_{n2}\mathrm{}U_1U_0\text{}$$
where $`U_{n1}`$, $`U_{n2}`$, $`\mathrm{}`$ ,$`U_1`$ ,$`U_0`$ are sufficiently local<sup>17</sup><sup>17</sup>17See Definition 8 in Section 7.7 of this paper for a definition of the term ‘sufficiently local’. unitary transformations. Each $`U_j`$ is called a computational step of the device.
Our first examples will be obtained by embedding classical computing devices within the realm of quantum mechanics. We will then look at some other quantum computing devices that are not the embeddings of classical devices.
#### 10.1. Embedding classical (memoryless) computation in quantum mechanics
One objective in this section is to represent<sup>18</sup><sup>18</sup>18Double meaning is intended. classical computing computing devices as unitary transformations. Since unitary transformations are invertible, i.e., reversible, it follows that the only classical computing devices that can be represented as such transformations must of necessity be reversible devices. Hence, the keen interest in reversible computation.
For a more in depth study of reversible computation, please refer to the work of Bennett and others.
#### 10.2. Classical reversible computation without memory
Each classical $`n`$-input/$`n`$-output (binary memoryless) reversible computing device (CRCD<sub>n</sub>) can be identified with a bijection
$$\pi :\{0,1\}^n\{0,1\}^n$$
on the set $`\{0,1\}^n`$ of all binary $`n`$-tuples. Thus, we can in turn identify each CRCD<sub>n</sub> with an element of the permutation group $`S_{2^n}`$ on the $`2^n`$ symbols
$$\{\stackrel{}{a}|\stackrel{}{a}\{0,1\}^n\}\text{ .}$$
Let
$$_n=x_0,x_1,\mathrm{},x_{n1}$$
denote the free Boolean ring on the symbols $`x_0,x_1,\mathrm{},x_{n1}`$ . Then the binary $`n`$-tuples $`\stackrel{}{a}\{0,1\}^n`$ are in one-to-one correspondence with the minterms of $`_n`$, i.e.,
$$\stackrel{}{a}x^\stackrel{}{a}=\underset{j=0}{\overset{n1}{}}x_j^{a_j}\text{ ,}$$
where
$$\{\begin{array}{ccc}x_j^0& =& \overline{x}_j\\ & & \\ x_j^1& =& x_j\end{array}$$
Since there is a one-to-one correspondence between the automorphisms of $`_n`$ and the permutations on the set of minterms, it follows that CRCD<sub>n</sub>’s can also be identified with the automorphism group $`Aut\left(_n\right)`$ of the free Boolean ring $`_n`$.
Moreover, since the set of binary $`n`$-tuples $`\{0,1\}^n`$ is in one-to-one correspondence with the set of integers $`\{0,1,2,\mathrm{},2^n1\}`$ via the radix 2 representation of integers, i.e.,
$$(b_{n1},b_{n2},\mathrm{},b_1,b_0)\underset{j=0}{\overset{n1}{}}b_j2^j\text{ ,}$$
we can, and frequently do, identify binary $`n`$-tuples with integers.
For example, consider the Controlled-NOT gate, called $`\mathrm{𝐂𝐍𝐎𝐓}`$ , which is defined by the following wiring diagram:
$$\mathrm{𝐂𝐍𝐎𝐓}=\overline{)\begin{array}{ccc}c& & b+c\\ & & \\ b& & b\\ & & \\ a& & a\end{array}}\text{ ,}$$
where ‘$``$’ and ‘$``$’ denote respectively a control bit and a target bit, and where ‘$`a+b`$’ denotes the exclusive ‘or’ of bits $`a`$ and $`b`$. This corresponds to the permutation $`\pi =(26)(37)`$, i.e.,
$$\{\begin{array}{ccccc}|0=& |000& & |000& =|0\\ |1=& |001& & |001& =|1\\ |2=& |010& & |110& =|6\\ |3=& |011& & |111& =|7\\ & & & & \\ |4=& |100& & |100& =|4\\ |5=& |101& & |101& =|5\\ |6=& |110& & |010& =|2\\ |7=& |111& & |011& =|3\end{array},$$
where we have used the following indexing conventions:
$$\{\begin{array}{c}\text{First=Right=Bottom}\hfill \\ \text{Last=Left=Top}\hfill \end{array}$$
As another example, consider the Toffoli gate , which is defined by the following wiring diagram:
$$\mathrm{𝐓𝐨𝐟𝐟𝐨𝐥𝐢}=\overline{)\begin{array}{ccc}c& & c+ab\\ & & \\ b& & b\\ & & \\ a& & a\end{array}}\text{ ,}$$
where ‘$`ab`$’ denotes the logical ‘and’ of $`a`$ and $`b`$. As before, ‘$`+`$’ denotes exclusive ‘or’. This gate corresponds to the permutation $`\pi =(67)`$.
In summary, we have:
$$\overline{)\underset{}{\overset{}{\left\{CRCD_n\right\}}}=S_2^n=Aut\left(_n\right)}$$
#### 10.3. Embedding classical irreversible computation within classical reversible computation
A classical 1-input/n-output (binary memoryless) irreversible computing device can be thought of as a Boolean function $`f=f(x_{n2},\mathrm{},x_1,x_0)`$ in $`_{n1}=x_0,x_1,\mathrm{},x_{n2}`$. Such irreversible computing devices can be transformed into reversible computing devices via the monomorphism
$$\iota :_{n1}Aut(_n),$$
where $`\iota (f)`$ is the automorphism in $`Aut(_n)`$ defined by
$$(x_{n1},x_{n2},\mathrm{},x_1,x_0)(x_{n1}f,x_{n2},\mathrm{},x_1,x_0),$$
and where ‘$``$’ denotes exclusive ‘or’. Thus, the image of each Boolean function $`f`$ is a product of disjoint transpositions in $`S_{2^n}`$.
As an additive group (ignoring ring structure), $`_{n1}`$ is the abelian group $`_{j=0}^{2^{(n1)}1}_2`$, where $`_2`$ denotes the cyclic group of order two.
Classical Binary Memoryless Computation is summarized in the table below:
| $`\stackrel{}{\text{Summary}}`$ |
| --- |
| $`\underset{}{\text{Classical Binary Memoryless Computation}}`$ |
| $`\underset{}{\overset{}{_{n1}=_{j=0}^{2^{(n1)}1}_2\stackrel{𝜄}{}S_{2^n}=Aut(_n)}}`$ |
#### 10.4. The unitary representation of reversible computing devices
It is now a straight forward task to represent CRCD<sub>n</sub>’s as unitary transformations. We simply use the standard unitary representation
$`\nu :S_2^n𝕌(2^n;)`$
of the symmetric group $`S_{2^n}`$ into the group of $`2^n\times 2^n`$ unitary matrices $`𝕌(2^n;)`$. This is the representation defined by
$$\pi \left(\delta _{k,\pi k}\right)_{2^n\times 2^n}\text{ ,}$$
where $`\delta _k\mathrm{}`$ denotes the Kronecker delta, i.e.,
$$\delta _k\mathrm{}=\{\begin{array}{cc}1& \text{if }k=\mathrm{}\hfill \\ & \\ 0& \text{otherwise}\hfill \end{array}$$
We think of such unitary transformations as quantum computing devices.
For example, consider the controlled-NOT gate $`\mathrm{𝐂𝐍𝐎𝐓}^{}=(45)(67)S_8`$ given by the wiring diagram
$$\mathrm{𝐂𝐍𝐎𝐓}^{}=\overline{)\begin{array}{ccc}c& & c\\ & & \\ b& & b\\ & & \\ a& & a+c\end{array}}$$
This corresponds to the unitary transformation
$$U_{\mathrm{𝐂𝐍𝐎𝐓}^{}}=\nu (\mathrm{𝐂𝐍𝐎𝐓}^{})=\left(\begin{array}{cccccccc}1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\end{array}\right)$$
Moreover, consider the Toffoli gate $`\mathrm{𝐓𝐨𝐟𝐟𝐨𝐥𝐢}^{}=(57)S_8`$ given by the wiring diagram
$$\mathrm{𝐓𝐨𝐟𝐟𝐨𝐥𝐢}^{}=\overline{)\begin{array}{ccc}c& & c\\ & & \\ b& & b+ac\\ & & \\ a& & a\end{array}}\text{ }$$
This corresponds to the unitary transformation
$$U_{\mathrm{𝐓𝐨𝐟𝐟𝐨𝐥𝐢}^{}}=\nu (\mathrm{𝐓𝐨𝐟𝐟𝐨𝐥𝐢}^{})=\left(\begin{array}{cccccccc}1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\end{array}\right)$$
Whenever it is clear from context, we will use the name of a CRCD<sub>n</sub> to also refer to the unitary transformation corresponding to the CRCD<sub>n</sub>. For example, we will denote $`\nu (CNOT)`$ and $`\nu (Toffoli)`$ simply by $`CNOT`$ and $`Toffoli`$. Moreover we will also use the wiring diagram of a CRCD<sub>n</sub> to refer to the unitary transformation corresponding to the CRCD<sub>n</sub>. For quantum computation beginners, this can lead to some confusion. Be careful!
#### 10.5. Some other simple quantum computing devices
After CRCD<sub>n</sub>’s are embedded as quantum computing devices, they are no longer classical computing devices. After the embedding, they suddenly have acquired much more computing power. Their inputs and outputs can be a superposition of many states. They can entangle their outputs. It is misleading to think of their input qubits as separate, for they could be entangled.
As an illustration of this fact, please note that the quantum computing device $`\mathrm{𝐂𝐍𝐎𝐓}^{\prime \prime }`$ given by the wiring diagram
$$\mathrm{𝐂𝐍𝐎𝐓}^{\prime \prime }=\overline{)\begin{array}{ccc}b& & a+b\\ & & \\ a& & a\end{array}}=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right)$$
is far from classical. It is more than a permutation. It is a linear operator that respects quantum superposition.
For example, $`\mathrm{𝐂𝐍𝐎𝐓}^{\prime \prime }`$ can take two non-entangled qubits as input, and then produce two entangled qubits as output. This is something no classical computing device can do. For example,
$$\frac{|0|1}{\sqrt{2}}|0=\frac{1}{\sqrt{2}}\left(|00|10\right)\frac{1}{\sqrt{2}}\left(|00|11\right)$$
For completeness, we list two other quantum computing devices that are embeddings of CRCD<sub>n</sub>’s, $`\mathrm{𝐍𝐎𝐓}`$ and $`\mathrm{𝐒𝐖𝐀𝐏}`$:
$$\mathrm{𝐍𝐎𝐓}=\overline{)\begin{array}{ccc}a& \text{NOT}& a+1\end{array}}=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)=\sigma _1$$
and
$$\mathrm{𝐒𝐖𝐀𝐏}=\overline{)\begin{array}{ccccc}b& & & & a\\ & & & & \\ a& & & & b\end{array}}=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 1\end{array}\right)$$
#### 10.6. Quantum computing devices that are not embeddings
We now consider quantum computing devices that are not embeddings of CRCD<sub>n</sub>’s.
The Hadamard gate $`𝐇`$ is defined as:
$$𝐇=\overline{)\begin{array}{c}𝐇\end{array}}=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}\hfill 1& \hfill 1\\ \hfill 1& \hfill 1\end{array}\right)\text{ .}$$
Another quantum gate is the square root of NOT, i.e., $`\sqrt{\mathrm{𝐍𝐎𝐓}}`$, which is given by
$$\sqrt{\mathrm{𝐍𝐎𝐓}}=\overline{)\begin{array}{c}\sqrt{\mathrm{𝐍𝐎𝐓}}\end{array}}=\frac{1i}{2}\left(\begin{array}{cc}\hfill i& \hfill 1\\ \hfill 1& \hfill i\end{array}\right)=\frac{1+i}{2}\left(\begin{array}{cc}\hfill 1& \hfill i\\ \hfill i& \hfill 1\end{array}\right)\text{ .}$$
There is also the square root of swap $`\sqrt{\mathrm{𝐒𝐖𝐀𝐏}}`$ which is defined as:
$$\sqrt{\mathrm{𝐒𝐖𝐀𝐏}}=\overline{)\begin{array}{c}\sqrt{\mathrm{𝐒𝐖𝐀𝐏}}\end{array}}=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& \frac{1+i}{2}& \frac{1i}{2}& 0\\ 0& \frac{1i}{2}& \frac{1+i}{2}& 0\\ 0& 0& 0& 1\end{array}\right)\text{ .}$$
Three frequently used unary quantum gates are the rotations:
$$\overline{)\begin{array}{c}\overline{)e^{i\theta \sigma _1}}\end{array}}=\left(\begin{array}{cc}\hfill \mathrm{cos}\theta & \hfill i\mathrm{sin}\theta \\ \hfill i\mathrm{sin}\theta & \hfill \mathrm{cos}\theta \end{array}\right)=e^{i\theta \sigma _1}$$
$$\overline{)\begin{array}{c}\overline{)e^{i\theta \sigma _2}}\end{array}}=\left(\begin{array}{cc}\hfill \mathrm{cos}\theta & \hfill \mathrm{sin}\theta \\ \hfill \mathrm{sin}\theta & \hfill \mathrm{cos}\theta \end{array}\right)=e^{i\theta \sigma 2}$$
$$\overline{)\begin{array}{c}\overline{)e^{i\theta \sigma _3}}\end{array}}=\left(\begin{array}{cc}e^{i\theta }& 0\\ 0& e^{i\theta }\end{array}\right)=e^{i\theta \sigma _3}$$
#### 10.7. <br>The implicit frame of a wiring diagram
Wiring diagrams have the advantage of being a simple means of describing some rather complicated unitary transformations. However, they do have their drawbacks, and they can, if we are not careful, be even misleading.
One problem with wiring diagrams is that they are not frame (i.e., basis) independent descriptions of unitary transformations. Each wiring diagram describes a unitary transformation using an implicitly understood basis.
For example, consider $`\mathrm{𝐂𝐍𝐎𝐓}^{\prime \prime }`$ given by the wiring diagram:
$$\mathrm{𝐂𝐍𝐎𝐓}^{\prime \prime }=\overline{)\begin{array}{ccc}b& & a+b\\ & & \\ a& & a\end{array}}\text{ .}$$
The above wiring diagram defines $`\mathrm{𝐂𝐍𝐎𝐓}^{\prime \prime }`$ in terms of the implicitly understood basis
$$\left\{|0=\left(\begin{array}{c}1\\ 0\end{array}\right),|1=\left(\begin{array}{c}0\\ 1\end{array}\right)\right\}\text{ .}$$
This wiring diagram suggests that qubit #1 controls qubit #0, and that qubit #1 is not effected by qubit #0. But this is far from the truth. For, $`\mathrm{𝐂𝐍𝐎𝐓}^{\prime \prime }`$ transforms
$$\frac{|0+|1}{\sqrt{2}}\frac{|0|1}{\sqrt{2}}$$
into
$$\frac{|0|1}{\sqrt{2}}\frac{|0|1}{\sqrt{2}}\text{ ,}$$
where we have used our indexing conventions
$$\{\begin{array}{c}\text{First=Right=Bottom}\hfill \\ \text{Last=Left=Top}\hfill \end{array}\text{ .}$$
In fact, in the basis
$$\left\{|0^{}=\frac{|0+|1}{\sqrt{2}},|1^{}=\frac{|0|1}{\sqrt{2}}\right\}$$
the wiring diagram of the same unitary transformation $`\mathrm{𝐂𝐍𝐎𝐓}^{\prime \prime }`$ is:
$$\overline{)\begin{array}{ccc}b& & a+b\\ & & \\ a& & a\end{array}}$$
The roles of the target and control qubits appeared to have switched!
### 11. <br>The No-Cloning Theorem
In this section, we prove the no-cloning theorem of Wootters and Zurek . The theorem states that there can be no device that produces exact replicas or copies of a quantum state. (See also for an elegant proof using the creation operators of quantum electrodynamics.)
The proof is an amazingly simple application of the linearity of quantum mechanics. The key idea is that copying is an inherently quadratic transformation, while the unitary transformations of quantum mechanics are inherently linear. Ergo, copying can not be a unitary transformation.
But what do we mean by a quantum replicator?
###### Definition 11.
Let $``$ be a Hilbert space. Then a quantum replicator consists of an auxiliary Hilbert space $`_A`$, a fixed state $`|\psi _0_A`$ (called the initial state of replicator), and a unitary transformation
$$U:_A_A$$
such that, for some fixed state $`|blank`$,
$$U|\psi _0|a|blank=|\psi _a|a|a\text{ ,}$$
for all states $`|a`$, where $`|\psi _a_A`$ (called the replicator state after replication of $`|a`$) depends on $`|a`$.
Since a quantum state is determined by a ket up to a multiplicative non-zero complex number, we can without loss of generality assume that $`|\psi _0`$, $`|a`$, $`|blank`$ are all of unit length. From unitarity, it follows that $`|\psi _a`$ is also of unit length.
Let $`|a`$, $`|b`$ be two kets of unit length in $``$ such that
$$0<\left|ab\right|<1\text{ .}$$
Then
$$\{\begin{array}{ccc}U|\psi _0|a|blank& =& |\psi _a|a|a\\ & & \\ U|\psi _0|b|blank& =& |\psi _b|b|b\end{array}$$
Hence,
$`blank\left|a\left|\psi _0\left|U^{}U\right|\psi _0\right|b\right|blank`$ $`=blank\left|a\left|\psi _0\psi _0\right|b\right|blank`$
$`=ab`$
On the other hand,
$`blank\left|a\left|\psi _0\left|U^{}U\right|\psi _0\right|b\right|blank`$ $`=a\left|a\left|\psi _a\psi _b\right|b\right|b`$
$`=ab^2\psi _a\psi _b`$
Thus,
$$ab^2\psi _a\psi _b=ab\text{ .}$$
And so,
$$ab\psi _a\psi _b=1\text{ .}$$
But this equation can not be satisfied since
$$\left|ab\right|<1$$
and
$$\left|\psi _a\psi _b\right||\psi _a|\psi _b=1$$
Hence, a quantum replicator cannot exist.
### 12. <br>Quantum teleportation
We now give a brief description of quantum teleportation, a means possibly to be used by future quantum computers to bus qubits from one location to another.
As stated earlier, qubits can not be copied as a result of the no-cloning theorem. (Please refer to the previous section.) However, they can be teleported, as has been demonstrated in laboratory settings. Such a mechanism could be used to bus qubits from one computer location to another. It could be used to create devices called quantum repeaters.
But what do we mean by teleportation?
Teleportation is the transferring of an object from one location to another by a process that:
* Firstly dissociates (i.e., destroys) the object to obtain information. – The object to be teleported is first scanned to extract sufficient information to reassemble the original object.
* Secondly transmits the acquired information from one location to another.
* Lastly reconstructs the object at the new location from the received information. – An exact replicas re-assembled at the destination out of locally available materials.
Two key effects of teleportation should be noted:
* The original object is destroyed during the process of teleportation. Hence, the no-cloning theorem is not violated.
* An exact replica of the original object is created at the intended destination.
Scotty of the Starship Enterprise was gracious enough to loan me the following teleportation manual. So I am passing it on to you.
Quantum Teleportation Manual
At location A, construct an EPR pair of qubits (qubits #2 and #3) in $`_2_3`$.
| $`|00`$ | $`\underset{}{\overset{}{\begin{array}{c}\text{Unitary}\\ \text{Matrix}\end{array}}}`$ | $`\frac{|01|10}{\sqrt{2}}`$ |
| --- | --- | --- |
| $`_2_3`$ | $``$ | $`_2_3`$ |
Physically transport entangled qubit #3 from location A to location B.
The qubit to be teleported, i.e., qubit #1, is delivered to location A in an unknown state
$$a|0+b|1$$
As a result of Steps 1 - 3, we have:
* Locations A and B share an EPR pair, i.e.
+ The qubit which is to be teleported, i.e., qubit #1, is at Location A
+ Qubit #2 is at Location A
+ Qubit #3 is at Location B
+ Qubits #2 & #3 are entangled
* The current state $`|\mathrm{\Phi }`$ of all three qubits is:
$$|\mathrm{\Phi }=\left(a|0+b|1\right)\left(\frac{|01|10}{\sqrt{2}}\right)_1_2_3$$
To better understand what is about to happen, we re-express the state $`|\mathrm{\Phi }`$ of the three qubits in terms of the following basis (called the Bell basis) of $`_1_2`$ :
$$\{\begin{array}{ccc}|\mathrm{\Psi }_A& =& \frac{|10|01}{\sqrt{2}}\\ & & \\ |\mathrm{\Psi }_B& =& \frac{|10+|01}{\sqrt{2}}\\ & & \\ |\mathrm{\Psi }_C& =& \frac{|00|11}{\sqrt{2}}\\ & & \\ |\mathrm{\Psi }_D& =& \frac{|00+|11}{\sqrt{2}}\end{array}$$
The result is:
$$\begin{array}{ccc}|\mathrm{\Phi }& =& \begin{array}{c}\frac{1}{2}[|\mathrm{\Psi }_A(a|0b|1)\\ +|\mathrm{\Psi }_B\left(a|0+b|1\right)\\ +|\mathrm{\Psi }_C\left(a|1+b|0\right)\\ +|\mathrm{\Psi }_D\left(a\right|1b|0)]\end{array}\end{array}\text{ ,}$$
where, as you might have noticed, we have written the expression in a suggestive way.
###### Remark 11.
Please note that since the completion of Step 3, we have done nothing physical. We have simply performed some algebraic manipulation of the expression representing the state $`|\mathrm{\Phi }`$ of the three qubits.
Let $`U:_1_2_1_2`$ be the unitary transformation defined by
$$\{\begin{array}{ccc}\underset{}{\overset{}{|\mathrm{\Psi }_A}}& & |00\\ \underset{}{\overset{}{|\mathrm{\Psi }_B}}& & |01\\ \underset{}{\overset{}{|\mathrm{\Psi }_C}}& & |10\\ \underset{}{\overset{}{|\mathrm{\Psi }_D}}& & |11\end{array}$$
<sup>19</sup><sup>19</sup>19Actually, there is no need to apply the unitary transformation $`U`$. We could have instead made a complete Bell state measurement, i.e., a measurement with respect to the compatible observables $`|\mathrm{\Psi }_A\mathrm{\Psi }_A|`$, $`|\mathrm{\Psi }_B\mathrm{\Psi }_B|`$, $`|\mathrm{\Psi }_C\mathrm{\Psi }_C|`$, $`|\mathrm{\Psi }_D\mathrm{\Psi }_D|`$. We have added the additional step 4 to make quantum teleportation easier to understand for quantum computation beginners. Please note that a complete Bell state measurement has, of this writing, yet to be achieved in a laboratoy setting.
Apply the local unitary transformation $`UI:_1_2_3_1_2_3`$ to the three qubits (actually more precisely, to qubits #1 and #2). Thus, under $`UI`$ the state $`|\mathrm{\Phi }`$ of all three qubits becomes
$$\begin{array}{ccc}|\mathrm{\Phi }^{}& =& \begin{array}{c}\frac{1}{2}[|00(a|0b|1)\\ +|01\left(a|0+b|1\right)\\ +|10\left(a|1+b|0\right)\\ +|11\left(a\right|1b|0)]\end{array}\hfill \end{array}$$
Measure qubits #1 and #2 to obtain two bits of classical information. The result of this measurement will be one of the bit pairs $`\{00,01,10,11\}`$.
Send from location A to location B (via a classical communication channel) the two classical bits obtained in Step 6.
As an intermediate summary, we have:
* Qubit #1 has been disassembled, and
* The information obtained during disassembly (two classical bits) has been sent to location B.
The two bits $`(i,j)`$ received from location A are used to select from the following table a unitary transformation $`U^{(i,j)}`$ of $`_3`$, (i.e., a local unitary transformation $`I_4U^{(i,j)}`$ on $`_1_2_3`$)
| Rec. Bits | $`U^{(i,j)}`$ | Future effect on qubit #3 |
| --- | --- | --- |
| $`00`$ | $`U^{(00)}=\left(\begin{array}{cc}\hfill 1& \hfill 0\\ \hfill 0& \hfill 1\end{array}\right)`$ | $`a|0b|1a|0+b|1`$ |
| $`01`$ | $`U^{(01)}=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ | $`a|0+b|1a|0+b|1`$ |
| $`10`$ | $`U^{(10)}=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ | $`a|1+b|0a|0+b|1`$ |
| $`11`$ | $`U^{(11)}=\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right)`$ | $`a|1b|0a|0+b|1`$ |
The unitary transformation $`U^{(i,j)}`$ selected in Step 7 is applied to qubit #3.
As a result, qubit #3 is at location B and has the original state of qubit #1 when qubit #1 was first delivered to location A, i.e., the state
$$a|0+b|1$$
It is indeed amazing that no one knows the state of the quantum teleported qubit except possibly the individual that prepared the qubit. Knowledge of the actual state of the qubit is not required for teleportaton. If its state is unknown before the teleportation, it remains unknown after the teleportation. All that we know is that the states before and after the teleportation are the same.
### 13. Shor’s algorithm
The following description of Shor’s algorithm is based on , , , , and .
#### 13.1. Preamble to Shor’s algorithm
There are cryptographic systems (such as RSA<sup>20</sup><sup>20</sup>20RSA is a public key cryptographic system invented by Rivest, Shamir, Adleman. Hence the name. For more information, please refer to .) that are extensively used today (e.g., in the banking industry) which are based on the following questionable assumption, i.e., conjecture:
Conjecture(Assumption). Integer factoring is computationally much harder than integer multiplication. In other words, while there are obviously many polynomial time algorithms for integer multiplication, there are no polynomial time algorithms for integer factoring. I.e., integer factoring computationally requires super-polynomial time.
This assumption is based on the fact that, in spite of the intensive efforts over many centuries of the best minds to find a polynomial time factoring algorithm, no one has succeeded so far. As of this writing, the most asymptotically efficient *classical* algorithm is the number theoretic sieve , , which factors an integer $`N`$ in time $`O\left(\mathrm{exp}\left[\left(\mathrm{lg}N\right)^{1/3}\left(\mathrm{lg}\mathrm{lg}N\right)^{2/3}\right]\right)`$. Thus, this is a super-polynomial time algorithm in the number $`O\left(\mathrm{lg}N\right)`$ of digits in $`N`$.
However, … Peter Shor suddenly changed the rules of the game.
Hidden in the above conjecture is the unstated, but implicitly understood, assumption that all algorithms run on computers based on the principles of classical mechanics, i.e., on classical computers. But what if a computer could be built that is based not only on classical mechanics, but on quantum mechanics as well? I.e., what if we could build a quantum computer?
Shor, starting from the works of Benioff, Bennett, Deutsch , Feynman, Simon, and others, created an algorithm to be run on a quantum computer, i.e., a quantum algorithm, that factors integers in polynomial time! Shor’s algorithm takes asymptotically $`O\left(\left(\mathrm{lg}N\right)^2\left(\mathrm{lg}\mathrm{lg}N\right)\left(\mathrm{lg}\mathrm{lg}\mathrm{lg}N\right)\right)`$ steps on a quantum computer, which is polynomial time in the number of digits $`O\left(\mathrm{lg}N\right)`$ of $`N`$.
#### 13.2. Number theoretic preliminaries
Since the time of Euclid, it has been known that every positive integer $`N`$ can be uniquely (up to order) factored into the product of primes. Moreover, it is a computationally easy (polynomial time) task to determine whether or not $`N`$ is a prime or composite number. For the primality testing algorithm of Miller-Rabin makes such a determination at the cost of $`O\left(s\mathrm{lg}N\right)`$ arithmetic operations \[$`O\left(s\mathrm{lg}^3N\right)`$ bit operations\] with probability of error $`Prob_{Error}2^s`$.
However, once an odd positive integer $`N`$ is known to be composite, it does not appear to be an easy (polynomial time) task on a classical computer to determine its prime factors. As mentioned earlier, so far the most asymptotically efficient *classical* algorithm known is the number theoretic sieve , , which factors an integer $`N`$ in time $`O\left(\mathrm{exp}\left[\left(\mathrm{lg}N\right)^{1/3}\left(\mathrm{lg}\mathrm{lg}N\right)^{2/3}\right]\right)`$.
Prime Factorization Problem. Given a composite odd positive integer $`N`$, find its prime factors.
It is well known that factoring $`N`$ can be reduced to the task of choosing at random an integer $`m`$ relatively prime to $`N`$, and then determining its modulo $`N`$ multiplicative order $`P`$, i.e., to finding the smallest positive integer $`P`$ such that
$$m^P=1\mathrm{mod}N\text{ .}$$
It was precisely this approach to factoring that enabled Shor to construct his factoring algorithm.
#### 13.3. Overview of Shor’s algorithm
But what is Shor’s quantum factoring algorithm?
Let $`=\{0,1,2,3,\mathrm{}\}`$ denote the set of natural numbers.
Shor’s algorithm provides a solution to the above problem. His algorithm consists of the five steps (steps 1 through 5), with only $`𝕊𝕋𝔼`$ 2 requiring the use of a quantum computer. The remaining four other steps of the algorithm are to be performed on a classical computer.
We begin by briefly describing all five steps. After that, we will then focus in on the quantum part of the algorithm, i.e., $`𝕊𝕋𝔼`$ 2.
* Choose a random positive even integer $`m`$. Use the polynomial time Euclidean algorithm<sup>21</sup><sup>21</sup>21The Euclidean algorithm is $`O\left(\mathrm{lg}^2N\right)`$. For a description of the Euclidean algorithm, see for example or . to compute the greatest common divisor $`\mathrm{gcd}(m,N)`$ of $`m`$ and $`N`$. If the greatest common divisor $`\mathrm{gcd}(m,N)1`$, then we have found a non-trivial factor of $`N`$, and we are done. If, on the other hand, $`\mathrm{gcd}(m,N)=1`$, then proceed to $`𝕊𝕋𝔼`$ 2.
* Use a quantum computer to determine the unknown period $`P`$ of the function
$$\begin{array}{ccc}& \stackrel{f_N}{}& \\ a& & m^a\mathrm{mod}N\end{array}$$
* If $`P`$ is an odd integer, then goto Step 1. \[The probability of $`P`$ being odd is $`(\frac{1}{2})^k`$, where $`k`$ is the number of distinct prime factors of $`N`$.\] If $`P`$ is even, then proceed to Step 4.
* Since $`P`$ is even,
$$\left(m^{P/2}1\right)\left(m^{P/2}+1\right)=m^P1=0\mathrm{mod}N\text{ .}$$
If $`m^{P/2}+1=0\mathrm{mod}N`$, then goto Step 1. If $`m^{P/2}+10\mathrm{mod}N`$, then proceed to Step 5. It can be shown that the probability that $`m^{P/2}+1=0\mathrm{mod}N`$ is less than $`(\frac{1}{2})^{k1}`$, where $`k`$ denotes the number of distinct prime factors of $`N`$.
* Use the Euclidean algorithm to compute $`d=\mathrm{gcd}(m^{P/2}1,N)`$. Since $`m^{P/2}+10\mathrm{mod}N`$, it can easily be shown that $`d`$ is a non-trivial factor of $`N`$. Exit with the answer $`d`$.
Thus, the task of factoring an odd positive integer $`N`$ reduces to the following problem:
Problem. Given a periodic function
$$f:\text{ ,}$$
find the period $`P`$ of $`f`$.
#### 13.4. Preparations for the quantum part of Shor’s algorithm
Choose a power of 2
$$Q=2^L$$
such that
$$N^2Q=2^L<2N^2\text{ ,}$$
and consider $`f`$ restricted to the set
$$S_Q=\{0,1,\mathrm{},Q1\}$$
which we also denote by $`f`$, i.e.,
$$f:S_QS_Q\text{ .}$$
In preparation for a discussion of $`𝕊𝕋𝔼`$ 2 of Shor’s algorithm, we construct two $`L`$-qubit quantum registers, Register1 and Register2 to hold respectively the arguments and the values of the function $`f`$, i.e.,
$$|\text{Reg1}|\text{Reg2}=|a|f(a)=|a|b=|a_0a_1\mathrm{}a_{L1}|b_0b_1\mathrm{}b_{L1}$$
In doing so, we have adopted the following convention for representing integers in these registers:
Notation Convention. In a quantum computer, we represent an integer $`a`$ with radix $`2`$ representation
$$a=\underset{j=0}{\overset{L1}{}}a_j2^j\text{ , }$$
as a quantum register consisting of the $`2^n`$ qubits
$$|a=|a_0a_1\mathrm{}a_{L1}=\underset{j=0}{\overset{L1}{}}|a_j$$
For example, the integer $`23`$ is represented in our quantum computer as $`n`$ qubits in the state:
$$|23=|10111000\mathrm{}0$$
Before continuing, we remind the reader of the classical definition of the $`Q`$-point Fourier transform.
###### Definition 12.
Let $`\omega `$ be a primitive $`Q`$-th root of unity, e.g., $`\omega =e^{2\pi i/Q}`$. Then the $`Q`$-point Fourier transform is the map
$$Map(S_Q,)\stackrel{}{}Map(S_Q,)$$
$$[f:S_Q][\widehat{f}:S_Q]$$
where
$$\widehat{f}\left(y\right)=\frac{1}{\sqrt{Q}}\underset{xS_Q}{}f(x)\omega ^{xy}$$
We implement the Fourier transform $``$ as a unitary transformation, which in the standard basis
$$|0,|1,\mathrm{},|Q1$$
is given by the $`Q\times Q`$ unitary matrix
$$=\frac{1}{\sqrt{Q}}\left(\omega ^{xy}\right)\text{ .}$$
This unitary transformation can be factored into the product of $`O\left(\mathrm{lg}^2Q\right)=O\left(\mathrm{lg}^2N\right)`$ sufficiently local unitary transformations. (See , .)
#### 13.5. The quantum part of Shor’s algorithm
The quantum part of Shor’s algorithm, i.e., $`𝕊𝕋𝔼`$ 2, is the following:
* Initialize registers 1 and 2, i.e.,
$$|\psi _0=|\text{Reg1}|\text{Reg2}=|0|0=|00\mathrm{}0|0\mathrm{}0$$
* <sup>22</sup><sup>22</sup>22In this step we could have instead applied the Hadamard transform to Register1 with the same result, but at the computational cost of $`O\left(\mathrm{lg}N\right)`$ sufficiently local unitary transformations.
Apply the $`Q`$-point Fourier transform $``$ to Register1.
$$|\psi _0=|0|0\stackrel{I}{}|\psi _1=\frac{1}{\sqrt{Q}}\underset{x=0}{\overset{Q1}{}}\omega ^{0x}|x|0=\frac{1}{\sqrt{Q}}\underset{x=0}{\overset{Q1}{}}|x|0$$
###### Remark 12.
Hence, Register1 now holds all the integers
$$0,1,2,\mathrm{},Q1$$
in superposition.
* Let $`U_f`$ be the unitary transformation that takes $`|x|0`$ to $`|x|f(x)`$. Apply the linear transformation $`U_f`$ to the two registers. The result is:
$$|\psi _1=\frac{1}{\sqrt{Q}}\underset{x=0}{\overset{Q1}{}}|x|0\stackrel{U_f}{}|\psi _2=\frac{1}{\sqrt{Q}}\underset{x=0}{\overset{Q1}{}}|x|f(x)$$
###### Remark 13.
The state of the two registers is now more than a superposition of states. In this step, we have quantum entangled the two registers.
* Apply the $`Q`$-point Fourier transform $``$ to Reg1. The resulting state is:
$$\begin{array}{ccc}|\psi _2=\frac{1}{\sqrt{Q}}\underset{x=0}{\overset{Q1}{}}|x|f(x)& \stackrel{I}{}& |\psi _3=\frac{1}{Q}\underset{x=0}{\overset{Q1}{}}\underset{y=0}{\overset{Q1}{}}\omega ^{xy}|y|f(x)\hfill \\ & & \\ & & =\frac{1}{Q}\underset{y=0}{\overset{Q1}{}}|\mathrm{{\rm Y}}(y)|y\frac{|\mathrm{{\rm Y}}(y)}{|\mathrm{{\rm Y}}(y)}\text{ ,}\hfill \end{array}$$
where
$$|\mathrm{{\rm Y}}(y)=\underset{x=0}{\overset{Q1}{}}\omega ^{xy}|f(x)\text{}$$
* Measure Reg1, i.e., perform a measurement with respect to the orthogonal projections
$$|00|I,|11|I,|22|I,\mathrm{},|Q1Q1|I\text{ ,}$$
where $`I`$ denotes the identity operator on the Hilbert space of the second register Reg2.
As a result of this measurement, we have, with probability
$$Prob\left(y_0\right)=\frac{|\mathrm{{\rm Y}}(y_0)^2}{Q^2}\text{ ,}$$
moved to the state
$$|y_0\frac{|\mathrm{{\rm Y}}(y_0)}{|\mathrm{{\rm Y}}(y_0)}$$
and measured the value
$$y_0\{0,1,2,\mathrm{},Q1\}\text{ . }$$
If after this computation, we ignore the two registers Reg1 and Reg2, we see that what we have created is nothing more than a classical probability distribution $`𝒮`$ on the sample space
$$\{0,1,2,\mathrm{},Q1\}\text{ .}$$
In other words, the sole purpose of executing STEPS 2.1 to 2.4 is to create a classical finite memoryless stochastic source $`𝒮`$ which outputs a symbol $`y_0\{0,1,2,\mathrm{},Q1\}`$ with the probability
$$Prob(y_0)=\frac{|\mathrm{{\rm Y}}(y_0)^2}{Q^2}\text{ .}$$
(For more details, please refer to section 8.1 of this paper.)
As we shall see, the objective of the remander of Shor’s algorithm is to glean information about the period $`P`$ of $`f`$ from the just created stochastic source $`𝒮`$. The stochastic source was created exactly for that reason.
#### 13.6. Peter Shor’s stochastic source $`𝒮`$
Before continuing to the final part of Shor’s algorithm, we need to analyze the probability distribution $`Prob\left(y\right)`$ a little more carefully.
###### Proposition 1.
Let $`q`$ and $`r`$ be the unique non-negative integers such that $`Q=Pq+r`$ , where $`0r<P`$ ; and let $`Q_0=Pq`$. Then
$$Prob\left(y\right)=\{\begin{array}{ccc}\frac{r\mathrm{sin}^2\left(\frac{\pi Py}{Q}\left(\frac{Q_0}{P}+1\right)\right)+\left(Pr\right)\mathrm{sin}^2\left(\frac{\pi Py}{Q}\frac{Q_0}{P}\right)}{Q^2\mathrm{sin}^2\left(\frac{\pi Py}{Q}\right)}\hfill & \hfill \text{if}& Py0\mathrm{mod}Q\hfill \\ & & \\ \frac{r\left(Q_0+P\right)^2+\left(Pr\right)Q_0^2}{Q^2P^2}\hfill & \hfill \text{if}& Py=0\mathrm{mod}Q\hfill \end{array}$$
###### Proof.
We begin by deriving a more usable expression for $`|\mathrm{{\rm Y}}(y)`$.
$$\begin{array}{ccc}\hfill |\mathrm{{\rm Y}}(y)& \hfill =& \underset{x=0}{\overset{Q1}{}}\omega ^{xy}|f(x)=\underset{x=0}{\overset{Q_01}{}}\omega ^{xy}|f(x)+\underset{x=Q_0}{\overset{Q1}{}}\omega ^{xy}|f(x)\hfill \\ & & \\ & \hfill =& \underset{x_0=0}{\overset{P1}{}}\underset{x_1=0}{\overset{\frac{Q_0}{P}1}{}}\omega ^{\left(Px_1+x_0\right)y}|f(Px_1+x_0)+\underset{x_0=0}{\overset{r1}{}}\omega ^{\left[P\left(\frac{Q_0}{P}\right)+x_0\right]y}|f(Px_1+x_0)\hfill \\ & & \\ & \hfill =& \underset{x_0=0}{\overset{P1}{}}\omega ^{x_0y}\left(\underset{x_1=0}{\overset{\frac{Q_0}{P}1}{}}\omega ^{Pyx_1}\right)|f(x_0)+\underset{x_0=0}{\overset{r1}{}}\omega ^{x_0y}\omega ^{Py\left(\frac{Q_0}{P}\right)}|f(x_0)\hfill \\ & & \\ & \hfill =& \underset{x_0=0}{\overset{r1}{}}\omega ^{x_0y}\left(\underset{x_1=0}{\overset{\frac{Q_0}{P}}{}}\omega ^{Pyx_1}\right)|f(x_0)+\underset{x_0=r}{\overset{P1}{}}\omega ^{x_0y}\left(\underset{x_1=0}{\overset{\frac{Q_0}{P}1}{}}\omega ^{Pyx_1}\right)|f(x_0)\hfill \end{array}$$
where we have used the fact that $`f`$ is periodic of period $`P`$.
Since $`f`$ is one-to-one when restricted to its period $`0,1,2,\mathrm{},P1`$, all the kets
$$|f(0),|f(1),|f(2),\mathrm{},|f(P1),$$
are mutually orthogonal. Hence,
$$\mathrm{{\rm Y}}(y)\mathrm{{\rm Y}}(y)=r\left|\underset{x_1=0}{\overset{\frac{Q_0}{P}}{}}\omega ^{Pyx_1}\right|^2+(Pr)\left|\underset{x_1=0}{\overset{\frac{Q_0}{P}1}{}}\omega ^{Pyx_1}\right|^2\text{ .}$$
If $`Py=0\mathrm{mod}Q`$, then since $`\omega `$ is a $`Q`$-th root of unity, we have
$$\mathrm{{\rm Y}}(y)\mathrm{{\rm Y}}(y)=r\left(\frac{Q_0}{P}+1\right)^2+\left(Pr\right)\left(\frac{Q_0}{P}\right)^2\text{ .}$$
On the other hand, if $`Py0\mathrm{mod}Q`$, then we can sum the geometric series to obtain
$`\mathrm{{\rm Y}}(y)\mathrm{{\rm Y}}(y)`$ $`=|{\displaystyle \frac{\omega ^{Py\left(\frac{Q_0}{P}+1\right)}1}{\omega ^{Py}1}}|^2+(Pr))\left|{\displaystyle \frac{\omega ^{Py\left(\frac{Q_0}{P}\right)}1}{\omega ^{Py}1}}\right|^2`$
$`=|{\displaystyle \frac{e^{\frac{2\pi i}{Q}Py\left(\frac{Q_0}{P}+1\right)}1}{e^{\frac{2\pi i}{Q}Py}1}}|^2+(Pr))\left|{\displaystyle \frac{e^{\frac{2\pi i}{Q}Py\left(\frac{Q_0}{P}\right)}1}{e^{\frac{2\pi i}{Q}Py}1}}\right|^2`$
where we have used the fact that $`\omega `$ is the primitive $`Q`$-th root of unity given by
$$\omega =e^{2\pi i/Q}\text{ .}$$
The remaining part of the proposition is a consequence of the trigonometric identity
$$\left|e^{i\theta }1\right|^2=4\mathrm{sin}^2\left(\frac{\theta }{2}\right)\text{ .}$$
As a corollary, we have
###### Corollary 1.
If $`P`$ is an exact divisor of $`Q`$, then
$$Prob\left(y\right)=\{\begin{array}{ccc}0\hfill & \hfill \text{if}& Py0\mathrm{mod}Q\hfill \\ & & \\ \frac{1}{P}\hfill & \hfill \text{if}& Py=0\mathrm{mod}Q\hfill \end{array}$$
#### 13.7. A momentary digression: Continued fractions
We digress for a moment to review the theory of continued fractions. (For a more in-depth explanation of the theory of continued fractions, please refer to and .)
Every positive rational number $`\xi `$ can be written as an expression in the form
$$\xi =a_0+\frac{1}{a_1+\frac{\underset{}{\overset{}{1}}}{a_2+\frac{\underset{}{\overset{}{1}}}{a_3+\frac{\underset{}{\overset{}{1}}}{\mathrm{}+\frac{\underset{}{\overset{}{1}}}{\stackrel{}{a_N}}}}}}\text{ ,}$$
where $`a_0`$ is a non-negative integer, and where $`a_1,\mathrm{},a_N`$ are positive integers. Such an expression is called a (finite, simple) continued fraction , and is uniquely determined by $`\xi `$ provided we impose the condition $`a_N>1`$. For typographical simplicity, we denote the above continued fraction by
$$[a_0,a_1,\mathrm{},a_N]\text{ .}$$
The continued fraction expansion of $`\xi `$ can be computed with the following recurrence relation, which always terminates if $`\xi `$ is rational:
$$\underset{}{\overset{}{\begin{array}{ccc}\{\begin{array}{c}\hfill a_0=\xi \\ \\ \hfill \xi _0=\xi a_0\end{array}\text{ ,}\hfill & \text{and if }\xi _n0\text{, then}\hfill & \{\begin{array}{c}a_{n+1}=1/\xi _n\hfill \\ \\ \xi _{n+1}=\frac{1}{\xi _n}a_{n+1}\hfill \end{array}\hfill \end{array}}}$$
The $`n`$-th convergent ($`0nN`$) of the above continued fraction is defined as the rational number $`\xi _n`$ given by
$$\xi _n=[a_0,a_1,\mathrm{},a_n]\text{ .}$$
Each convergent $`\xi _n`$ can be written in the for, $`\xi _n=\frac{p_n}{q_n}`$, where $`p_n`$ and $`q_n`$ are relatively prime integers ( $`\mathrm{gcd}(p_n,q_n)=1`$). The integers $`p_n`$ and $`q_n`$ are determined by the recurrence relation
$$\overline{)\begin{array}{ccc}p_0=a_0,\hfill & p_1=a_1a_0+1,\hfill & p_n=a_np_{n1}+p_{n2},\hfill \\ & & \\ q_0=1,\hfill & q_1=a_1,\hfill & q_n=a_nq_{n1}+q_{n2}\text{ .}\hfill \end{array}}$$
#### 13.8. Preparation for the final part of Shor’s algorithm
###### Definition 13.
<sup>23</sup><sup>23</sup>23$`\left\{a\right\}_Q=aQround\left(\frac{a}{Q}\right)=aQ\frac{a}{Q}+\frac{1}{2}`$.
For each integer $`a`$, let $`\left\{a\right\}_Q`$ denote the residue of $`a`$ modulo $`Q`$ of smallest magnitude. In other words, $`\left\{a\right\}_Q`$ is the unique integer such that
$$\{\begin{array}{c}a=\left\{a\right\}_Q\mathrm{mod}Q\hfill \\ \\ Q/2<\left\{a\right\}_QQ/2\hfill \end{array}\text{ .}$$
###### Proposition 2.
Let $`y`$ be an integer lying in $`S_Q`$. Then
$$Prob\left(y\right)\{\begin{array}{ccc}\frac{4}{\pi ^2}\frac{1}{P}\left(1\frac{1}{N}\right)^2\hfill & \hfill \text{if}& 0<\left|\left\{Py\right\}_Q\right|\frac{P}{2}\left(1\frac{1}{N}\right)\hfill \\ & & \\ \frac{1}{P}\left(1\frac{1}{N}\right)^2\hfill & \hfill \text{if}& \left\{Py\right\}_Q=0\hfill \end{array}$$
###### Proof.
We begin by noting that
$$\begin{array}{cc}\left|\frac{\pi \left\{Py\right\}_Q}{Q}\left(\frac{Q_0}{P}+1\right)\right|\hfill & \frac{\pi }{Q}\frac{P}{2}\left(1\frac{1}{N}\right)\left(\frac{Q_0+P}{P}\right)\frac{\pi }{2}\left(1\frac{1}{N}\right)\left(\frac{Q+P}{Q}\right)\hfill \\ & \\ & \frac{\pi }{2}\left(1\frac{1}{N}\right)\left(1+\frac{P}{Q}\right)\frac{\pi }{2}\left(1\frac{1}{N}\right)\left(1+\frac{N}{N^2}\right)<\frac{\pi }{2}\text{ ,}\hfill \end{array}$$
where we have made use of the inequalities
$$N^2Q<2N^2\text{ and }0<PN\text{ .}$$
It immediately follows that
$$\left|\frac{\pi \left\{Py\right\}_Q}{Q}\frac{Q_0}{P}\right|<\frac{\pi }{2}\text{ .}$$
As a result, we can legitimately use the inequality
$$\frac{4}{^{\pi ^2}}\theta ^2\mathrm{sin}^2\theta \theta ^2\text{, for }\left|\theta \right|<\frac{\pi }{2}$$
to simplify the expression for $`Prob\left(y\right)`$.
Thus,
$$\begin{array}{ccc}Prob\left(y\right)\hfill & =\hfill & \frac{r\mathrm{sin}^2\left(\frac{\pi \left\{Py\right\}_Q}{Q}\left(\frac{Q_0}{P}+1\right)\right)+\left(Pr\right)\mathrm{sin}^2\left(\frac{\pi \left\{Py\right\}_Q}{Q}\frac{Q_0}{P}\right)}{Q^2\mathrm{sin}^2\left(\frac{\pi Py}{Q}\right)}\hfill \\ & & \\ & \hfill & \frac{r\frac{4}{\pi ^2}\left(\frac{\pi \left\{Py\right\}_Q}{Q}\left(\frac{Q_0}{P}+1\right)\right)^2+\left(Pr\right)\frac{4}{\pi ^2}\left(\frac{\pi \left\{Py\right\}_Q}{Q}\frac{Q_0}{P}\right)^2}{Q^2\left(\frac{\pi \left\{Py\right\}_Q}{Q}\right)^2}\hfill \\ & & \\ & \hfill & \frac{4}{\pi ^2}\frac{P\left(\frac{Q_0}{P}\right)^2}{Q^2}=\frac{4}{\pi ^2}\frac{1}{P}\left(\frac{Qr}{Q}\right)^2\hfill \\ & & \\ & =\hfill & \frac{4}{\pi ^2}\frac{1}{P}\left(1\frac{r}{Q}\right)^2\frac{4}{\pi ^2}\frac{1}{P}\left(1\frac{1}{N}\right)^2\hfill \end{array}$$
The remaining case, $`\left\{Py\right\}_Q=0`$ is left to the reader. ∎
###### Lemma 1.
Let
$$Y=\{yS_Q\left|\left\{Py\right\}_Q\right|\frac{P}{2}\}\text{ and }S_P=\{dS_Q0d<P\}\text{ .}$$
Then the map
$$\begin{array}{ccc}Y\hfill & \hfill & S_P\hfill \\ y\hfill & \hfill & d=d(y)=round\left(\frac{P}{Q}y\right)\hfill \end{array}$$
is a bijection with inverse
$$y=y(d)=round\left(\frac{Q}{P}d\right)\text{ .}$$
Hence, $`Y`$ and $`S_P`$ are in one-to-one correspondence. Moreover,
$$\left\{Py\right\}_Q=PyQd(y)\text{ .}$$
###### Remark 14.
Moreover, the following two sets of rationals are in one-to-one correspondence
$$\{\frac{y}{Q}yY\}\{\frac{d}{P}0d<P\}$$
As a result of the measurement performed in $`𝕊𝕋𝔼`$ 2.4, we have in our possession an integer $`yY`$. We now show how $`y`$ can be use to determine the unknown period $`P`$.
We now need the following theorem<sup>24</sup><sup>24</sup>24See \[42, Theorem 184, Section 10.15\]. from the theory of continued fractions:
###### Theorem 2.
Let $`\xi `$ be a real number, and let $`a`$ and $`b`$ be integers with $`b>0`$. If
$$\left|\xi \frac{a}{b}\right|\frac{1}{2b^2}\text{ ,}$$
then the rational number $`a/b`$ is a convergent of the continued fraction expansion of $`\xi `$.
As a corollary, we have:
###### Corollary 2.
If $`\left|\left\{Py\right\}_Q\right|\frac{P}{2}`$, then the rational number $`\frac{d(y)}{P}`$ is a convergent of the continued fraction expansion of $`\frac{y}{Q}`$.
###### Proof.
Since
$$PyQd(y)=\left\{Py\right\}_Q\text{ ,}$$
we know that
$$\left|PyQd(y)\right|\frac{P}{2}\text{}$$
which can be rewritten as
$$\left|\frac{y}{Q}\frac{d(y)}{P}\right|\frac{1}{2Q}\text{ .}$$
But, since $`QN^2`$, it follows that
$$\left|\frac{y}{Q}\frac{d(y)}{P}\right|\frac{1}{2N^2}\text{ .}$$
Finally, since $`PN`$ (and hence $`\frac{1}{2N}\frac{1}{2P^2})`$, the above theorem can be applied. Thus, $`\frac{d(y)}{P}`$ is a convergent of the continued fraction expansion of $`\xi =\frac{y}{Q}`$. ∎
Since $`\frac{d(y)}{P}`$ is a convergent of the continued fraction expansion of $`\frac{y}{Q}`$, it follows that, for some $`n`$,
$$\frac{d(y)}{P}=\frac{p_n}{q_n}\text{ ,}$$
where $`p_n`$ and $`q_n`$ are relatively prime positive integers given by a recurrence relation found in the previous subsection. So it would seem that we have found a way of deducing the period $`P`$ from the output $`y`$ of $`𝕊𝕋𝔼`$ 2.4, and so we are done.
Not quite!
We can determine $`P`$ from the measured $`y`$ produced by $`𝕊𝕋𝔼`$ 2.4, only if
$$\{\begin{array}{c}p_n=d(y)\hfill \\ \\ q_n=P\hfill \end{array}\text{ ,}$$
which is true only when $`d(y)`$ and $`P`$ are relatively prime.
So what is the probability that the $`yY`$ produced by $`𝕊𝕋𝔼`$ 2.4 satisfies the additional condition that
$$\mathrm{gcd}(P,d(y))=1\text{ ?}$$
###### Proposition 3.
The probability that the random $`y`$ produced by $`𝕊𝕋𝔼`$ 2.4 is such that $`d(y)`$ and $`P`$ are relatively prime is bounded below by the following expression
$$Prob\{yY\mathrm{gcd}(d(y),P)=1\}\frac{4}{\pi ^2}\frac{\varphi (P)}{P}\left(1\frac{1}{N}\right)^2\text{ ,}$$
where $`\varphi (P)`$ denotes Euler’s totient function, i.e., $`\varphi (P)`$ is the number of positive integers less than $`P`$ which are relatively prime to $`P`$.
The following theorem can be found in \[42, Theorem 328, Section 18.4\]:
###### Theorem 3.
$$liminf\frac{\varphi (N)}{N/\mathrm{ln}\mathrm{ln}N}=e^\gamma \text{,}$$
where $`\gamma `$ denotes Euler’s constant $`\gamma =0.57721566490153286061\mathrm{}`$ , and where $`e^\gamma =0.5614594836\mathrm{}`$ .
As a corollary, we have:
###### Corollary 3.
$$Prob\{yY\mathrm{gcd}(d(y),P)=1\}\frac{4}{\pi ^2\mathrm{ln}2}\frac{e^\gamma ϵ\left(P\right)}{\mathrm{lg}\mathrm{lg}N}\left(1\frac{1}{N}\right)^2\text{ ,}$$
where $`ϵ\left(P\right)`$ is a monotone decreasing sequence converging to zero. In terms of asymptotic notation,
$$Prob\{yY\mathrm{gcd}(d(y),P)=1\}=\mathrm{\Omega }\left(\frac{1}{\mathrm{lg}\mathrm{lg}N}\right)\text{ .}$$
Thus , if $`𝕊𝕋𝔼`$ 2.4 is repeated $`O(\mathrm{lg}\mathrm{lg}N)`$ times, then the probability of success is $`\mathrm{\Omega }\left(1\right)`$.
###### Proof.
From the above theorem, we know that
$$\frac{\varphi (P)}{P/\mathrm{ln}\mathrm{ln}P}e^\gamma ϵ\left(P\right)\text{ .}$$
where $`ϵ\left(P\right)`$ is a monotone decreasing sequence of positive reals converging to zero. Thus,
$$\frac{\varphi (P)}{P}\frac{e^\gamma ϵ\left(P\right)}{\mathrm{ln}\mathrm{ln}P}\frac{e^\gamma ϵ\left(P\right)}{\mathrm{ln}\mathrm{ln}N}=\frac{e^\gamma ϵ\left(P\right)}{\mathrm{ln}\mathrm{ln}2+\mathrm{ln}\mathrm{lg}N}\frac{e^\gamma ϵ\left(P\right)}{\mathrm{ln}2}\frac{1}{\mathrm{lg}\mathrm{lg}N}$$
###### Remark 15.
$`\mathrm{\Omega }(\frac{1}{\mathrm{lg}\mathrm{lg}N})`$ denotes an asymptotic lower bound. Readers not familiar with the big-oh $`O()`$ and big-omega $`\mathrm{\Omega }()`$ notation should refer to \[19, Chapter 2\] or \[11, Chapter 2\].
###### Remark 16.
For the curious reader, lower bounds $`LB(P)`$ of $`e^\gamma ϵ\left(P\right)`$ for $`3P841`$ are given in the following table:
| $`P`$ | $`LB(P)`$ |
| --- | --- |
| 3 | 0.062 |
| 4 | 0.163 |
| 5 | 0.194 |
| 7 | 0.303 |
| 13 | 0.326 |
| 31 | 0.375 |
| 61 | 0.383 |
| 211 | 0.411 |
| 421 | 0.425 |
| 631 | 0.435 |
| 841 | 0.468 |
Thus, if one wants a reasonable bound on the $`Prob\{yY\mathrm{gcd}(d(y),P)=1\}`$ before continuing with Shor’s algorithm, it would pay to first use a classical algorithm to verify that the period $`P`$ of the randomly chosen integer $`m`$ is not too small.
#### 13.9. The final part of Shor’s algorithm
We are now prepared to give the last step in Shor’s algorithm. This step can be performed on a classical computer.
* Compute the period $`P`$ from the integer $`y`$ produced by $`𝕊𝕋𝔼`$ 2.4.
* + Loop for each $`n`$ from $`n=1`$ Until $`\xi _n=0`$.
* + - Use the recurrence relations given in subsection 13.7, to compute the $`p_n`$ and $`q_n`$ of the $`n`$-th convergent $`\frac{p_n}{q_n}`$ of $`\frac{y}{Q}`$.
* + - Test to see if $`q_n=P`$ by computing<sup>25</sup><sup>25</sup>25The indicated algorithm for computing $`m^{q_n}\mathrm{mod}N`$ requires $`O(\mathrm{lg}q_n)`$ arithmetic operations.
$$m^{q_n}=\underset{i}{}\left(m^{2^i}\right)^{q_{n,i}}\mathrm{mod}N\text{ ,}$$
where $`q_n=_iq_{n,i}2^i`$ is the binary expansion of $`q_n`$.
- If $`m^{q_n}=1\mathrm{mod}N`$, then exit with the answer $`P=q_n`$, and proceed to Step 3. If not, then continue the loop.
* + End of Loop
* + If you happen to reach this point, you are a very unlucky quantum computer scientist. You must start over by returning to $`𝕊𝕋𝔼`$ 2.0. But don’t give up hope! The probability that the integer $`y`$ produced by $`𝕊𝕋𝔼`$ 2.4 will lead to a successful completion of Step 2.5 is bounded below by
$$\frac{4}{\pi ^2\mathrm{ln}2}\frac{e^\gamma ϵ\left(P\right)}{\mathrm{lg}\mathrm{lg}N}\left(1\frac{1}{N}\right)^2>\frac{0.232}{\mathrm{lg}\mathrm{lg}N}\left(1\frac{1}{N}\right)^2\text{ ,}$$
provided the period $`P`$ is greater than $`3`$. \[ $`\gamma `$ denotes Euler’s constant.\]
#### 13.10. <br>An example of Shor’s algorithm
Let us now show how $`N=91(=713)`$ can be factored using Shor’s algorithm.
We choose $`Q=2^{14}=16384`$ so that $`N^2Q<2N^2`$.
* Choose a random positive integer $`m`$, say $`m=3`$. Since $`\mathrm{gcd}(91,3)=1`$, we proceed to $`𝕊𝕋𝔼`$ 2 to find the period of the function $`f`$ given by
$$f(a)=3^a\mathrm{mod}91$$
###### Remark 17.
Unknown to us, $`f`$ has period $`P=6`$. For,
| $`\begin{array}{ccccccccccc}a& & 0& 1& 2& 3& 4& 5& 6& 7& \mathrm{}\\ & & & & & & & & & & \\ f(a)& & 1& 3& 9& 27& 81& 61& 1& 3& \mathrm{}\end{array}`$ |
| --- |
| $`\text{ Unknown period }P=6`$ |
* Initialize registers 1 and 2. Thus, the state of the two registers becomes:
$$|\psi _0=|0|0$$
* Apply the $`Q`$-point Fourier transform $``$ to register #1, where
$$|k=\frac{1}{\sqrt{16384}}\underset{x=0}{\overset{16383}{}}\omega ^{kj}|x\text{ ,}$$
and where $`\omega `$ is a primitive $`Q`$-th root of unity, e.g., $`\omega =e^{\frac{2\pi i}{16384}}`$. Thus the state of the two registers becomes:
$$|\psi _1=\frac{1}{\sqrt{16384}}\underset{x=0}{\overset{16383}{}}|x|0$$
* Apply the unitary transformation $`U_f`$ to registers #1 and #2, where
$$U_f|x|\mathrm{}=|x|f(x)\mathrm{}\mathrm{mod}91\text{ .}$$
(Please note that $`U_f^2=I`$.) Thus, the state of the two registers becomes:
$$\begin{array}{cccc}\hfill |\psi _2& \hfill =& \hfill \frac{1}{\sqrt{16384}}& _{x=0}^{16383}|x|3^x\mathrm{mod}91\hfill \\ & & & \\ & \hfill =& \hfill \frac{1}{\sqrt{16384}}(& |0|1+|1|3+|2|9+|3|27+|4|81+|5|61\hfill \\ & & & \\ & & & +|6|1+|7|3+|8|9+|9|27+|10|81+|11|61\hfill \\ & & & \\ & & & +|12|1+|13|3+|14|9+|15|27+|16|81+|17|61\hfill \\ & & & \\ & & & +\mathrm{}\hfill \\ & & & \\ & & & +|16380|1+|16381|3+|16382|9+|16383|27\hfill \\ & & \hfill )& \end{array}$$
###### Remark 18.
The state of the two registers is now more than a superposition of states. We have in the above step quantum entangled the two registers.
* Apply the $`Q`$-point $``$ again to register #1. Thus, the state of the system becomes:
$$\begin{array}{ccc}\hfill |\psi _3& \hfill =& \frac{1}{\sqrt{16384}}_{x=0}^{16383}\frac{1}{\sqrt{16384}}_{y=0}^{16383}\omega ^{xy}|y|3^x\mathrm{mod}91\hfill \\ & & \\ & \hfill =& \frac{1}{16384}_{x=0}^{16383}|y_{x=0}^{16383}\omega ^{xy}|3^x\mathrm{mod}91\hfill \\ & & \\ & \hfill =& \frac{1}{16384}_{x=0}^{16383}|y|\mathrm{{\rm Y}}\left(y\right)\text{ ,}\hfill \end{array}$$
where
$$|\mathrm{{\rm Y}}\left(y\right)=\underset{x=0}{\overset{16383}{}}\omega ^{xy}|3^x\mathrm{mod}91$$
Thus,
$$\begin{array}{cc}\hfill |\mathrm{{\rm Y}}\left(y\right)=& |1+\omega ^y|3+\omega ^{2y}|9+\omega ^{3y}|27+\omega ^{4y}|81+\omega ^{5y}|61\hfill \\ & \\ & +\omega ^{6y}|1+\omega ^{7y}|3+\omega ^{8y}|9+\omega ^{9y}|27+\omega ^{10y}|81+\omega ^{11y}|61\hfill \\ & \\ & +\omega ^{12y}|1+\omega ^{13y}|3+\omega ^{14y}|9+\omega ^{15y}|27+\omega ^{16y}|81+\omega ^{17y}|61\hfill \\ & \\ & +\mathrm{}\hfill \\ & \\ & +\omega ^{16380y}|1+\omega ^{16381y}|3+\omega ^{16382y}|9+\omega ^{16383y}|27\hfill \end{array}$$
* Measure Reg1. The result of our measurement just happens to turn out to be
$$y=13453$$
Unknown to us, the probability of obtaining this particular $`y`$ is:
$$0.3189335551\times 10^6\text{ . }$$
Moreover, unknown to us, we’re lucky! The corresponding $`d`$ is relatively prime to $`P`$, i.e.,
$$d=d(y)=round(\frac{P}{Q}y)=5$$
However, we do know that the probability of $`d(y)`$ being relatively prime to $`P`$ is greater than
$$\frac{0.232}{\mathrm{lg}\mathrm{lg}N}\left(1\frac{1}{N}\right)^28.4\%\text{ (provided }P>3\text{),}$$
and we also know that
$$\frac{d(y)}{P}$$
is a convergent of the continued fraction expansion of
$$\xi =\frac{y}{Q}=\frac{13453}{16384}$$
So with a reasonable amount of confidence, we proceed to Step 2.5.
* Using the recurrence relations found in subsection 13.7 of this paper, we successively compute (beginning with $`n=0`$) the $`a_n`$’s and $`q_n`$’s for the continued fraction expansion of
$$\xi =\frac{y}{Q}=\frac{13453}{16384}\text{ .}$$
For each non-trivial $`n`$ in succession, we check to see if
$$3^{q_n}=1\mathrm{mod}91\text{}$$
If this is the case, then we know $`q_n=P`$, and we immediately exit from Step 2.5 and proceed to Step 3.
* In this example, $`n=0`$ and $`n=1`$ are trivial cases.
* For $`n=2`$, $`a_2=4`$ and $`q_2=5`$ . We test $`q_2`$ by computing
$$3^{q_2}=3^5=\left(3^{2^0}\right)^1\left(3^{2^1}\right)^0\left(3^{2^0}\right)^1=611\mathrm{mod}91\text{ .}$$
Hence, $`q_2P`$.
* We proceed to $`n=3`$, and compute
$$a_3=1\text{ and }q_3=6\text{}$$
We then test $`q_3`$ by computing
$$3^{q_3}=3^6=\left(3^{2^0}\right)^0\left(3^{2^1}\right)^1\left(3^{2^0}\right)^1=1\mathrm{mod}91\text{ .}$$
Hence, $`q_3=P`$. Since we now know the period $`P`$, there is no need to continue to compute the remaining $`a_n`$’s and $`q_n`$’s. We proceed immediately to Step 3.
To satisfy the reader’s curiosity we have listed in the table below all the values of $`a_n`$, $`p_n`$, and $`q_n`$ for $`n=0,1,\mathrm{},14`$. But it should be mentioned again that we need only to compute $`a_n`$ and $`q_n`$ for $`n=0,1,2,3`$, as indicated above.
| $`n`$ | 0 | 1 | 2 | 3 | 4 | 5 | 6 | 7 | 8 | 9 | 10 | 11 | 12 | 13 | 14 |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`a_n`$ | 0 | 1 | 4 | 1 | 1 | 2 | 3 | 1 | 1 | 3 | 1 | 1 | 1 | 1 | 3 |
| $`p_n`$ | 0 | 1 | 4 | 5 | 9 | 23 | 78 | 101 | 179 | 638 | 817 | 1455 | 2272 | 3727 | 13453 |
| $`q_n`$ | 1 | 1 | 5 | 6 | 11 | 28 | 95 | 123 | 218 | 777 | 995 | 1772 | 2767 | 4539 | 16384 |
* Since $`P=6`$ is even, we proceed to Step 4.
* Since
$$3^{P/2}=3^3=271\mathrm{mod}91\text{}$$
we goto Step 5.
* With the Euclidean algorithm, we compute
$$\mathrm{gcd}(3^{P/2}1,91)=\mathrm{gcd}(3^31,91)=\mathrm{gcd}(26,91)=13\text{ .}$$
We have succeeded in finding a non-trivial factor of $`N=91`$, namely $`13`$. We exit Shor’s algorithm, and proceed to celebrate!
### 14. Grover’s Algorithm
The the following description of Grover’s algorithm is based on , , and .
#### 14.1. Problem definition
We consider the problem of searching an unstructured database of $`N=2^n`$ records for exactly one record which has been specifically marked. This can be rephrased in mathematical terms as an oracle problem as follows:
Label the records of the database with the integers
$$0,1,2,\mathrm{},N1\text{ ,}$$
and denote the label of the unknown marked record by $`x_0`$. We are given an oracle which computes the $`n`$ bit binary function
$$f:\{0,1\}^n\{0,1\}$$
defined by
$$f(x)=\{\begin{array}{cc}1& \text{if }x=x_0\hfill \\ & \\ 0& \text{otherwise}\hfill \end{array}$$
We remind the readers that, as a standard oracle idealization, we have no access to the internal workings of the function $`f`$. It operates simply as a blackbox function, which we can query as many times as we like. But with each such a query comes an associated computational cost.
Search Problem for an Unstructured Database. Find the record labeled as $`x_0`$ with the minimum amount of computational work, i.e., with the minimum number of queries of the oracle $`f`$.
From probability theory, we know that if we examine $`k`$ records, i.e., if we compute the oracle $`f`$ for $`k`$ randomly chosen records, then the probability of finding the record labeled as $`x_0`$ is $`k/N`$. Hence, on a classical computer it takes $`O(N)=O(2^n)`$ queries to find the record labeled $`x_0`$.
#### 14.2. The quantum mechanical perspective
However, as Luv Grover so astutely observed, on a quantum computer the search of an unstructured database can be accomplished in $`O(\sqrt{N})`$ steps, or more precisely, with the application of $`O(\sqrt{N}\mathrm{lg}N)`$ sufficiently local unitary transformations. Although this is not exponentially faster, it is a significant speedup.
Let $`_2`$ be a 2 dimensional Hilbert space with orthonormal basis
$$\{|0,|1\}\text{ ;}$$
and let
$$\{|0,|1,\mathrm{},|N1\}$$
denote the induced orthonormal basis of the Hilbert space
$$=\underset{0}{\overset{N1}{}}_2\text{ .}$$
From the quantum mechanical perspective, the oracle function $`f`$ is given as a blackbox unitary transformation $`U_f`$, i.e., by
$$\begin{array}{ccc}_2& \stackrel{U_f}{}& _2\\ & & \\ |x|y& & |x|f(x)y\end{array}$$
where ‘$``$’ denotes exclusive ‘OR’, i.e., addition modulo 2.<sup>26</sup><sup>26</sup>26Please note that $`U_f=\left(\nu \iota \right)(f)`$, as defined in sections 10.3 and 10.4 of this paper.
Instead of $`U_f`$, we will use the computationally equivalent unitary transformation
$$I_{|x_0}\left(|x\right)=(1)^{f(x)}|x=\{\begin{array}{cc}|x_0& \text{if }x=x_0\hfill \\ & \\ |x& \text{otherwise}\hfill \end{array}$$
That $`I_{|x_0}`$ is computationally equivalent to $`U_f`$ follows from the easily verifiable fact that
$$U_f\left(|x\frac{|0|1}{\sqrt{2}}\right)=\left(I_{|x_0}\left(|x\right)\right)\frac{|0|1}{\sqrt{2}}\text{ ,}$$
and also from the fact that $`U_f`$ can be constructed from a controlled-$`I_{|x_0}`$ and two one qubit Hadamard transforms. (For details, please refer to , .)
The unitary transformation $`I_{|x_0}`$ is actually an inversion in $``$ about the hyperplane perpendicular to $`|x_0`$. This becomes evident when $`I_{|x_0}`$ is rewritten in the form
$$I_{|x_0}=I2|x_0x_0|\text{ ,}$$
where ‘$`I`$’ denotes the identity transformation. More generally, for any unit length ket $`|\psi `$, the unitary transformation
$$I_{|\psi }=I2|\psi \psi |\text{ }$$
is an inversion in $``$ about the hyperplane orthogonal to $`|\psi `$.
#### 14.3. Properties of the inversion $`I_{|\psi }`$
We digress for a moment to discuss the properties of the unitary transformation $`I_{|\psi }`$. To do so, we need the following definition.
###### Definition 14.
Let $`|\psi `$ and $`|\chi `$ be two kets in $``$ for which the bracket product $`\psi \chi `$ is a real number. We define
$$𝒮_{}=Span_{}(|\psi ,|\chi )=\left\{\alpha |\psi +\beta |\chi \alpha ,\beta \right\}$$
as the sub-Hilbert space of $``$ spanned by $`|\psi `$ and $`|\chi `$. We associate with the Hilbert space $`𝒮_{}`$ a real inner product space lying in $`𝒮_{}`$ defined by
$$𝒮_{}=Span_{}(|\psi ,|\chi )=\left\{a|\psi +b|\chi a,b\right\}\text{ ,}$$
where the inner product on $`𝒮_{}`$ is that induced by the bracket product on $``$. If $`|\psi `$ and $`|\chi `$ are also linearly independent, then $`𝒮_{}`$ is a 2 dimensional real inner product space (i.e., the 2 dimensional Euclidean plane) lying inside of the complex 2 dimensional space $`𝒮_{}`$.
###### Proposition 4.
Let $`|\psi `$ and $`|\chi `$ be two linearly independent unit length kets in $``$ with real bracket product; and let $`𝒮_{}=Span_{}(|\psi ,|\chi )`$ and $`𝒮_{}=Span_{}(|\psi ,|\chi )`$. Then
* Both $`𝒮_{}`$ and $`𝒮_{}`$ are invariant under the transformations $`I_{|\psi }`$, $`I_{|\chi }`$, and hence $`I_{|\psi }I_{|\chi }`$, i.e.,
$$\overline{)\begin{array}{ccc}I_{|\psi }\left(𝒮_{}\right)=𝒮_{}\hfill & \hfill \text{and }& I_{|\psi }\left(𝒮_{}\right)=𝒮_{}\hfill \\ & & \\ I_{|\chi }\left(𝒮_{}\right)=𝒮_{}\hfill & \hfill \text{and }& I_{|\chi }\left(𝒮_{}\right)=𝒮_{}\hfill \\ & & \\ I_{|\psi }I_{|\chi }\left(𝒮_{}\right)=𝒮_{}\hfill & \hfill \text{and }& I_{|\psi }I_{|\chi }\left(𝒮_{}\right)=𝒮_{}\hfill \end{array}}$$
* If $`L_{|\psi ^{}}`$ is the line in the plane $`𝒮_{}`$ which passes through the origin and which is perpendicular to $`|\psi `$, then $`I_{|\psi }`$ restricted to $`𝒮_{}`$ is a reflection in (i.e., a Möbius inversion about) the line $`L_{|\psi ^{}}`$. A similar statement can be made in regard to $`|\chi `$.
* If $`|\psi ^{}`$ is a unit length vector in $`𝒮_{}`$ perpendicular to $`|\psi `$, then
$$I_{|\psi }=I_{|\psi ^{}}\text{ .}$$
(Hence, $`\psi ^{}\chi `$ is real.)
Finally we note that, since $`I_{|\psi }=I2|\psi \psi |`$, it follows that
###### Proposition 5.
If $`|\psi `$ is a unit length ket in $``$, and if $`U`$ is a unitary transformation on $``$, then
$$UI_{|\psi }U^1=I_{U|\psi }\text{ .}$$
#### 14.4. The method in Luv’s “madness”
Let $`H:`$ be the Hadamard transform, i.e.,
$$H=\underset{0}{\overset{n1}{}}H^{(2)}\text{ , }$$
where
$$H^{(2)}=\left(\begin{array}{cc}\hfill 1& \hfill 1\\ \hfill 1& \hfill 1\end{array}\right)$$
with respect to the basis $`|0`$, $`|1`$.
We begin by using the Hadamard transform $`H`$ to construct a state $`|\psi _0`$ which is an equal superposition of all the standard basis states $`|0`$, $`|1`$,$`\mathrm{}`$,$`|N1`$ (including the unknown state $`|x_0`$), i.e.,
$$|\psi _0=H|0=\frac{1}{\sqrt{N}}\underset{k=0}{\overset{N1}{}}|k\text{ .}$$
Both $`|\psi _0`$ and the unknown state $`|x_0`$ lie in the Euclidean plane $`𝒮_{}=Span_{}(|\psi _0,|x_0)`$. Our strategy is to rotate within the plane $`𝒮_{}`$ the state $`|\psi _0`$ about the origin until it is as close as possible to $`|x_0`$. Then a measurement with respect to the standard basis of the state resulting from rotating $`|\psi _0`$, will produce $`|x_0`$ with high probability.
To achieve this objective, we use the oracle $`I_{|x_0}`$ to construct the unitary transformation
$$Q=HI_{|0}H^1I_{|x_0}\text{ ,}$$
which by proposition 2 above, can be reexpressed as
$$Q=I_{|\psi _0}I_{|x_0}\text{ .}$$
Let $`|x_0^{}`$ and $`|\psi _0^{}`$ denote unit length vectors in $`𝒮_{}`$ perpendicular to $`|x_0`$ and $`|\psi _0`$, respectively. There are two possible choices for each of $`|x_0^{}`$ and $`|\psi _0^{}`$ respectively. To remove this minor, but nonetheless annoying, ambiguity, we select $`|x_0^{}`$ and $`|\psi _0^{}`$ so that the orientation of the plane $`𝒮_{}`$ induced by the ordered spanning vectors $`|\psi _0`$, $`|x_0`$ is the same orientation as that induced by each of the ordered bases $`|x_0^{}`$, $`|x_0`$ and $`|\psi _0`$, $`|\psi _0^{}`$. (Please refer to Figure 2.)
###### Remark 19.
The removal of the above ambiguities is really not essential. However, it does simplify the exposition given below.
Figure 2. The linear transformation $`Q|_𝒮_{}`$ is reflection in the line $`L_{|x_0^{}}`$ followed by reflection in the line $`L_{|\psi _0}`$ which is the same as rotation by the angle $`2\beta `$. Thus, $`Q|_𝒮_{}`$ rotates $`|\psi _0`$ by the angle $`2\beta `$ toward $`|x_0`$.
We proceed by noting that, by the above proposition 1, the plane $`𝒮_{}`$ lying in $``$ is invariant under the linear transformation $`Q`$, and that, when $`Q`$ is restricted to the plane $`𝒮_{}`$, it can be written as the composition of two inversions, i.e.,
$$Q|_𝒮_{}=I_{|\psi _0^{}}I_{|x_0}\text{ .}$$
In particular, $`Q|_𝒮_{}`$ is the composition of two inversions in $`𝒮_{}`$, the first in the line $`L_{|x_0^{}}`$ in $`𝒮_{}`$ passing through the origin having $`|x_0`$ as normal, the second in the line $`L_{|\psi _0}`$ through the origin having $`|\psi _0^{}`$ as normal.<sup>27</sup><sup>27</sup>27The line $`L_{|x_0^{}}`$ is the intersection of the plane $`𝒮_{}`$ with the hyperplane in $``$ orthogonal to $`|x_0`$. A similar statement can be made in regard to $`L_{|\psi _0}`$.
We can now apply the following theorem from plane geometry:
###### Theorem 4.
If $`L_1`$ and $`L_2`$ are lines in the Euclidean plane $`^2`$ intersecting at a point $`O`$; and if $`\beta `$ is the angle in the plane from $`L_1`$ to $`L_2`$, then the operation of reflection in $`L_1`$ followed by reflection in $`L_2`$ is just rotation by angle $`2\beta `$ about the point $`O`$.
Let $`\beta `$ denote the angle in $`S_{}`$ from $`L_{|x_0^{}}`$ to $`L_{|\psi _0}`$, which by plane geometry is the same as the angle from $`|x_0^{}`$ to $`|\psi _0`$, which in turn is the same as the angle from $`|x_0`$ to $`|\psi _0^{}`$. Then by the above theorem $`Q|_𝒮_{}=I_{|\psi _0^{}}I_{|x_0}`$ is a rotation about the origin by the angle $`2\beta `$.
The key idea in Grover’s algorithm is to move $`|\psi _0`$ toward the unknown state $`|x_0`$ by successively applying the rotation $`Q`$ to $`|\psi _0`$ to rotate it around to $`|x_0`$. This process is called amplitude amplification. Once this process is completed, the measurement of the resulting state (with respect to the standard basis) will, with high probability, yield the unknown state $`|x_0`$. This is the essence of Grover’s algorithm.
But how many times $`K`$ should we apply the rotation $`Q`$ to $`|\psi _0`$? If we applied $`Q`$ too many or too few times, we would over- or undershoot our target state $`|x_0`$.
We determine the integer $`K`$ as follows:
Since
$$|\psi _0=\mathrm{sin}\beta |x_0+\mathrm{cos}\beta |x_0^{}\text{ ,}$$
the state resulting after $`k`$ applications of $`Q`$ is
$$|\psi _k=Q^k|\psi _0=\mathrm{sin}\left[\left(2k+1\right)\beta \right]|x_0+\mathrm{cos}\left[\left(2k+1\right)\beta \right]|x_0^{}\text{ .}$$
Thus, we seek to find the smallest positive integer $`K=k`$ such that
$$\mathrm{sin}\left[\left(2k+1\right)\beta \right]$$
is as close as possible to $`1`$. In other words, we seek to find the smallest positive integer $`K=k`$ such that
$$\left(2k+1\right)\beta $$
is as close as possible to $`\pi /2`$. It follows that<sup>28</sup><sup>28</sup>28The reader may prefer to use the $`floor`$ function instead of the $`round`$ function.
$$K=k=round\left(\frac{\pi }{4\beta }\frac{1}{2}\right)\text{ ,}$$
where “$`round`$” is the function that rounds to the nearest integer.
We can determine the angle $`\beta `$ by noting that the angle $`\alpha `$ from $`|\psi _0`$ and $`|x_0`$ is complementary to $`\beta `$, i.e.,
$$\alpha +\beta =\pi /2\text{ ,}$$
and hence,
$$\frac{1}{\sqrt{N}}=x_0\psi _0=\mathrm{cos}\alpha =\mathrm{cos}(\frac{\pi }{2}\beta )=\mathrm{sin}\beta \text{ .}$$
Thus, the angle $`\beta `$ is given by
$$\beta =\mathrm{sin}^1\left(\frac{1}{\sqrt{N}}\right)\frac{1}{\sqrt{N}}\text{ (for large }N\text{) ,}$$
and hence,
$$K=k=round\left(\frac{\pi }{4\mathrm{sin}^1\left(\frac{1}{\sqrt{N}}\right)}\frac{1}{2}\right)round\left(\frac{\pi }{4}\sqrt{N}\frac{1}{2}\right)\text{ (for large }N\text{).}$$
#### 14.5. Summary of Grover’s algorithm
In summary, we provide the following outline of Grover’s algorithm:
Grover’s Algorithm $`\begin{array}{c}\hfill 𝕊𝕋𝔼\text{ 0.}\\ \end{array}`$ $`\begin{array}{c}\text{(Initialization)}\hfill \\ |\psi H|0=\frac{1}{\sqrt{N}}{\displaystyle \underset{j=0}{\overset{N1}{}}}|j\hfill \\ k0\hfill \end{array}`$ $`\begin{array}{c}\hfill 𝕊𝕋𝔼\text{ 1.}\\ \end{array}`$ $`\begin{array}{c}\hfill \text{Loop until }k=\underset{}{round\left(\frac{\pi }{4\mathrm{sin}^1\left(1/\sqrt{N}\right)}\frac{1}{2}\right)}round\left(\frac{\pi }{4}\sqrt{N}\frac{1}{2}\right)\\ |\psi \underset{}{Q}|\psi =HI_{|0}HI_{|x_0}|\psi \hfill \\ kk+1\hfill \end{array}`$ $`\begin{array}{c}\hfill 𝕊𝕋𝔼\text{ 2.}\end{array}`$ $`\begin{array}{c}\text{Measure }|\psi \text{ with respect to the standard basis}\hfill \\ |0,|1,\mathrm{},|N1\text{ to obtain the marked unknown }\hfill \\ \text{state }|x_0\text{ with probability }1\frac{1}{N}\text{.}\hfill \end{array}`$
We complete our summary with the following theorem:
###### Theorem 5.
With a probability of error<sup>29</sup><sup>29</sup>29If the reader prefers to use the $`floor`$ function rather than the $`round`$ function, then probability of error becomes $`Prob_E\frac{4}{N}\frac{4}{N^2}`$.
$$Prob_E\frac{1}{N}\text{}$$
Grover’s algorithm finds the unknown state $`|x_0`$ at a computational cost of
$$O\left(\sqrt{N}\mathrm{lg}N\right)$$
###### Proof.
* The probability of error $`Prob_E`$ of finding the hidden state $`|x_0`$ is given by
$$Prob_E=\mathrm{cos}^2\left[\left(2K+1\right)\beta \right]\text{ ,}$$
where
$$\{\begin{array}{ccc}\hfill \beta & \hfill =& \mathrm{sin}^1\left(\frac{1}{\sqrt{N}}\right)\hfill \\ & & \\ \hfill K& \hfill =& round\left(\frac{\pi }{4\beta }\frac{1}{2}\right)\hfill \end{array}\text{,}$$
where “$`round`$” is the function that rounds to the nearest integer. Hence,
$$\begin{array}{ccc}\hfill \frac{\pi }{4\beta }1K\frac{\pi }{4\beta }& \hfill & \frac{\pi }{2}\beta \left(2K+1\right)\beta \frac{\pi }{2}+\beta \hfill \\ & & \\ & \hfill & \mathrm{sin}\beta =\mathrm{cos}\left(\frac{\pi }{2}\beta \right)\mathrm{cos}\left[\left(2K+1\right)\beta \right]\mathrm{cos}\left(\frac{\pi }{2}+\beta \right)=\mathrm{sin}\beta \hfill \end{array}$$
Thus,
$$Prob_E=\mathrm{cos}^2\left[\left(2K+1\right)\beta \right]\mathrm{sin}^2\beta =\mathrm{sin}^2\left(\mathrm{sin}^1\left(\frac{1}{\sqrt{N}}\right)\right)=\frac{1}{N}$$
* The computational cost of the Hadamard transform $`H=_0^{n1}H^{(2)}`$ is $`O(n)=O(\mathrm{lg}N)`$ single qubit operations. The transformations $`I_{|0}`$ and $`I_{|x_0}`$ each carry a computational cost of $`O(1)`$.
$`𝕊𝕋𝔼`$ 1 is the computationally dominant step. In $`𝕊𝕋𝔼`$ 1 there are $`O\left(\sqrt{N}\right)`$ iterations. In each iteration, the Hadamard transform is applied twice. The transformations $`I_{|0}`$ and $`I_{|x_0}`$ are each applied once. Hence, each iteration comes with a computational cost of $`O\left(\mathrm{lg}N\right)`$, and so the total cost of $`𝕊𝕋𝔼`$ 1 is $`O(\sqrt{N}\mathrm{lg}N)`$.
#### 14.6. <br>An example of Grover’s algorithm
As an example, we search a database consisting of $`N=2^n=8`$ records for an unknown record with the unknown label $`x_0=5`$. The calculations for this example were made with OpenQuacks , which is an open source quantum simulator Maple package developed at UMBC and publically available.
We are given a blackbox computing device
In
I|?
OutIn
I|?
Out\text{In}\rightarrow\framebox{\framebox{\begin{tabular}[c]{l}$I_{\left|?\right\rangle}$\end{tabular}
}}\rightarrow\text{Out}
that implements as an oracle the unknown unitary transformation
$$I_{|x_0}=I_{|5}=\left(\begin{array}{ccccccccc}\hfill 1& \hfill 0& \hfill 0& \hfill 0& & \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0& & \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0& & \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& & \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ & & & & & & & & \\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& & \hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& & \hfill 0& \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& & \hfill 0& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& & \hfill 0& \hfill 0& \hfill 0& \hfill 1\end{array}\right)$$
We cannot open up the blackbox
I|?
absent
I|?
absent\rightarrow\framebox{$\framebox{\begin{tabular}[c]{l}$I_{\left|?\right\rangle}$\end{tabular}
}$}\rightarrow to see what is inside. So we do not know what $`I_{|x_0}`$ and $`x_0`$ are. The only way that we can glean some information about $`x_0`$ is to apply some chosen state $`|\psi `$ as input, and then make use of the resulting output.
Using of the blackbox
I|?
absent
I|?
absent\rightarrow\framebox{\framebox{\begin{tabular}[c]{l}$I_{\left|?\right\rangle}$\end{tabular}
}}\rightarrow as a component device, we construct a computing device
-HI|0HI|?
absent
-HI|0HI|?
absent\rightarrow\framebox{\framebox{\begin{tabular}[c]{l}$-HI_{\left|0\right\rangle}HI_{\left|?\right\rangle}$\end{tabular}
}}\rightarrow which implements the unitary operator
$$Q=HI_{|0}HI_{|x_0}=\frac{1}{4}\left(\begin{array}{ccccccccc}\hfill 3& \hfill 1& \hfill 1& \hfill 1& & \hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 3& \hfill 1& \hfill 1& & \hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 3& \hfill 1& & \hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 3& & \hfill 1& \hfill 1& \hfill 1& \hfill 1\\ & & & & & & & & \\ \hfill 1& \hfill 1& \hfill 1& \hfill 1& & \hfill 3& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1& & \hfill 1& \hfill 3& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1& & \hfill 1& \hfill 1& \hfill 3& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1& & \hfill 1& \hfill 1& \hfill 1& \hfill 3\end{array}\right)$$
We do not know what unitary transformation $`Q`$ is implemented by the device
-HI|0HI|?
absent
-HI|0HI|?
absent\rightarrow\framebox{\framebox{\begin{tabular}[c]{l}$-HI_{\left|0\right\rangle}HI_{\left|?\right\rangle}$\end{tabular}
}}\rightarrow because the blackbox
I|?
absent
I|?
absent\rightarrow\framebox{\framebox{\begin{tabular}[c]{l}$I_{\left|?\right\rangle}$\end{tabular}
}}\rightarrow is one of its essential components.
* We begin by preparing the known state
$$\overline{)|\psi _0=H|0=\frac{1}{\sqrt{8}}(1,1,1,1,1,1,1,1)^{transpose}}$$
* We proceed to loop
$$K=round\left(\frac{\pi }{4\mathrm{sin}^1\left(1/\sqrt{8}\right)}\frac{1}{2}\right)=2$$
times in $`𝕊𝕋𝔼`$ 1.
+ On the first iteration, we obtain the unknown state
$$\overline{)|\psi _1=Q|\psi _0=\frac{1}{4\sqrt{2}}(1,1,1,1,5,1,1,1)^{transpose}}$$
+ On the second iteration, we obtain the unknown state
$$\overline{)|\psi _2=Q|\psi _1=\frac{1}{8\sqrt{2}}(1,1,1,1,11,1,1,1)^{transpose}}$$
and branch to $`𝕊𝕋𝔼`$ 2.
* We measure the unknown state $`|\psi _2`$ to obtain either
$$|5$$
with probability
$$Prob_{Success}=\mathrm{sin}^2\left(\left(2K+1\right)\beta \right)=\frac{121}{128}=0.9453$$
or some other state with probability
$$Prob_{Failure}=\mathrm{cos}^2\left(\left(2K+1\right)\beta \right)=\frac{7}{128}=0.0547$$
and then exit.
### 15. <br>There is much more to quantum computation
Needles to say, there is much more to quantum computation. I hope that you found this introductory paper useful.
### Index
* Adjoint §4.3
+ Big §5.8
+ Little §5.8
* Alice §7.6.1, §7.7
* Ancilla Remark 8
* Automorphism Group §10.2
* Bennett §10.1, §7.4
* Blackbox §14.1, §14.6
* Bob §7.6.1, §7.7
* Boolean Ring $`_n`$ §10.2
* Bra §4.1
* Bracket §4.1
* CNOT, see Controlled-NOT
* Commutator §4.8, §5.7
* Complex Projective Space $`P^{n1}`$ §4.2
* Computation
+ Irreversible §10.3
+ Reversible §10.1
* Computational Step §10
* Computing Device
+ Classical §10
+ Quantum §10
* Constituent “Part” Definition 3
* Continued Fraction §13.7
+ Convergent of §13.7
* Controlled-NOT §10.2, §10.4, §10.5
* Convergent, see Continued Fraction
* CRCD, see Computing Device, Classlical Reversible
* Density Operator §5§5.8
* Deutsch §13.1, §6
* Diagonalization §5.2, §8.2
* Dirac Notation §4.1§4.3, §4.7§4.7
* Dynamic Invariants §8.2
* Eigenket §4.4
* Eigenvalue §4.4
+ Degenerate Definition 1
+ Non-degenerate Definition 1
* Einstein §7.5
* Embedding §10, §10.3, §10.4
* Ensemble
+ Mixed §5.1, §5.4, Example 6, Example 6
+ Pure §5.1, §8.2
* Entangled, Quantum §7.1, §7.7
* Entanglement
+ Classes Definition 7
+ Quantum §7.4
+ Type Definition 7
* Entropy, Classical §8.1, §8.2
* Entropy, Shannon §8.1
* Entropy, von Neumann §8.2, §8.2
* EPR §7.3, §7.5, §7.6.1
* Euler’s constant Theorem 3
* Expected Value §4.7
* Filtration §4.5
* Fourier Transform Definition 12
* $`\gamma `$, see Euler’s Constant
* Global Quantum System Definition 3
* Global Unitary Transformation Definition 6
* Gottesman §6
* Grover’s Algorithm §14, §14.6
* Hamiltonian §4.9
* Heisenberg
+ Uncertainty Principle §4.8
* Heisenberg, Picture §6, §6
* Hermitian Operator §4.4
* Hilbert Space §3.2
* Inversion §14.2
* Juxtaposition §7.1
* Ket §4.1
* Kronecker Sum Definition 6
* $`\mathrm{lg}`$ Definition 9
* Lie Group §7.7
* Local
+ Sufficiently Definition 8
* Local Equivalence Definition 7
* Local Interaction item 1)
* Local Lie Algebra Definition 6
* Local Subgroup Definition 6
* Local Unitary Transformation Definition 6
* Measurement §4.5, §4.6, §5.3
* Mobius Inversion §14.2
* Multipartite Quantum System §5.6
* Nielsen §7.4
* No-Cloning Theorem §11
* Non-Locality §7.6
* NOT gate §10.5
* Observable §4.4
+ Complete §4.4
+ Incompatible Operators §4.8
+ Selective Measurement §4.5
* Observable, Measurement §4.4
* Observables
+ Compatible Operators §4.8
* OpenQuacks Public Domain Software §14.6
* Operator
+ Compatible §4.8
+ Hermitian §4.4
+ Incompatible §4.8
+ Measurement §4.4
+ Self-Adjoint §4.4
+ Unitary §4.9
* Orbits §7.7
* Partial Trace §5.5
* Partition Definition 8
* Path-Ordered Integral §4.10
* Pauli Spin Matrices Example 1
* Permutation Group §10.2
* Planck’s Constant §4.9
* Podolsky §7.5
* Polarized Light §2.2
* Positive Operator Valued Measure §4.5
* POVM §4.5
* Primality Testing §13.2
* Principle of Non-Locality §7.5, §7.6
* Probabilistic Operator Valued Measure §4.5
* Probability Distribution Definition 9
* Projection Operator §8.3
* Quantum
+ Repeater §12
+ Replicator Definition 11
* Quantum Entangled §7.1, §7.7, Definition 4
* Quantum Register §7.2
* Qubit §3.1, §3.3
* Radix 2 Representation §13.4
* Reality
+ Principle of §7.5
* Repeater
+ Quantum §12
* Replicator
+ Quantum Definition 11
* Reversible Computation §10.1
* Rosen §7.5
* Rotation §10.6
* $`round`$ function §14.4
* Schrodinger Picture §6
* Selective Measurement Operator §4.5
* Self-Adjoint Operator §4.4
* Shor §13
* Shor’s Algorithm §13.1, §13.10
* Spacelike Distance §7.6
* Square Root of
+ NOT §10.6
+ SWAP §10.6
* Standard Deviation §4.8
* Standard Unitary Representation §10.4
* Stochastic Source §8.1, Definition 9
* Sufficiently Local Unitary Transformation, see Unitary
* Superluminal Communication §7.6
* Superposition §3.1, §7.2, §7.2
* SWAP Gate §10.5
* Symbols
+ Output §8.1, §8.3
* Symmetric Group §10.4
* Teleportation §12
* Teleportation, Quantum §12, §12
* Tensor Product §4.1
* Toffoli Gate §10.2, §10.4
* Trace §5.5
* Unitary
+ Operator §4.9
+ Sufficiently Local Definition 8
+ Transformation §4.9
* $`\mathrm{\Omega }()`$, asymptotic lower bound Corollary 3
* Wiring Diagram §10.2§10.7
* Wiring Diagrams
+ Implicit Frame §10.7
* Wootters §11, §7.4
* Zurek §11
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# 1 Table
> To be submitted to:
Nucl. Phys A.
The $`\pi \pi `$ interaction in nuclear matter from a study of the $`\pi ^+A\pi ^+\pi ^\pm A^{}`$ reactions
F. Bonutti<sup>a,b</sup>, P. Camerini<sup>a,b</sup>, E. Fragiacomo<sup>a,b</sup>, N. Grion<sup>a,</sup><sup>1</sup><sup>1</sup>1Corresponding author, electronic mail: Nevio.Grion@ts.INFN.it, R. Rui<sup>a,b</sup>, J.T. Brack<sup>c,</sup><sup>2</sup><sup>2</sup>2Present address: University of Colorado, Boulder CO 80309-0446, USA, L. Felawka<sup>c</sup>, E.F. Gibson<sup>f</sup>, G. Hofman<sup>d,2</sup>, M. Kermani<sup>d,</sup>, E.L. Mathie<sup>e</sup>, R. Meier<sup>c,</sup><sup>3</sup><sup>3</sup>3Permanent address: Physikalisches Institut, Universität Tübingen, 72076 Tübingen, Germany, D. Ottewell<sup>c</sup>, K. Raywood<sup>c</sup>, M.E. Sevior<sup>g</sup>, G.R. Smith<sup>c</sup> and R. Tacik<sup>e</sup>.
<sup>a</sup> Istituto Nazionale di Fisica Nucleare, 34127 Trieste, Italy
<sup>b</sup> Dipartimento di Fisica dell’Universita’ di Trieste, 34127 Trieste, Italy
<sup>c</sup> TRIUMF, Vancouver, B.C., Canada V6T 2A3
<sup>d</sup> Department of Physics and Astronomy, University of British Columbia, Vancouver, B.C., Canada V6T 2A6
<sup>e</sup> University of Regina, Regina, Saskatchewan, Canada S4S 0A2
<sup>f</sup> California State University, Sacramento CA 95819, USA
<sup>g</sup> School of Physics, University of Melbourne, Parkville, Vic., 3052, Australia
The CHAOS Collaboration
The pion-production reactions $`\pi ^+A\pi ^+\pi ^\pm A^{}`$ were studied on $`{}_{}{}^{2}H`$, $`{}_{}{}^{12}C`$, $`{}_{}{}^{40}Ca`$, and $`{}_{}{}^{208}Pb`$ nuclei at an incident pion energy of $`T_{\pi ^+}`$=283 MeV. Pions were detected in coincidence using the CHAOS spectrometer. The experimental results are reduced to differential cross sections and compared to both theoretical predictions and the reaction phase space. The composite ratio $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ between the $`\pi ^+\pi ^\pm `$ invariant masses on nuclei and on the nucleon is also presented. Near the $`2m_\pi `$ threshold pion pairs couple to $`(\pi \pi )_{I=J=0}`$ when produced in the $`\pi ^+\pi ^+\pi ^{}`$ reaction channel. There is a marked near-threshold enhancement of $`𝒞`$$`{}_{}{}^{A}{}_{\pi ^+\pi ^{}}{}^{}`$ which is consistent with theoretical predictions addressing the partial restoration of chiral symmetry in nuclear matter. Furthermore, the behaviour of $`𝒞`$$`{}_{}{}^{A}{}_{\pi ^+\pi ^{}}{}^{}`$ is well described when the restoration of chiral symmetry is combined with standard P-wave renormalization of pions in nuclear matter. On the other hand, nuclear matter only weakly influences $`𝒞`$$`{}_{}{}^{A}{}_{\pi ^+\pi ^+}{}^{}`$, which displays a flat behaviour throughout the energy range regardless of $`A`$.
PACS:25.80 Hp
1. Introduction
The influence of the nuclear medium on the $`\pi \pi `$ interaction was investigated at TRIUMF using the pion induced pion-production reactions $`\pi ^+A\pi ^+\pi ^\pm A^{}`$ (henceforth labelled $`\pi 2\pi `$). The initial study was directed to the reactions in deuterium, that is, to $`\pi ^+n(p)\pi ^+\pi ^{}p(p)`$ and $`\pi ^+p(n)\pi ^+\pi ^+n(n)`$, in order to understand the $`\pi 2\pi `$ behaviour on both a neutron and a proton through quasi-free reactions. The $`\pi 2\pi `$ process was then examined on the complex nuclei $`{}_{}{}^{12}C`$, $`{}_{}{}^{40}Ca`$ and $`{}_{}{}^{208}Pb`$ in order to study possible $`\pi \pi `$ medium modifications by direct comparison of the $`\pi 2\pi `$ data. To ensure the validity of this approach, the following experimental method was applied. All the $`\pi 2\pi `$ data were taken under the same kinematical conditions. The final $`\pi ^+\pi ^\pm `$ pairs were detected in coincidence to ensure a reliable identification of $`\pi 2\pi `$ events. Pion pairs were analysed down to $`0^{}`$ $`\pi \pi `$ opening angles to determine the $`\pi \pi `$ invariant mass at the $`2m_\pi `$ threshold. The energy of the incident pion beam, $`T_{\pi ^+}`$=283 MeV, was chosen to enable the investigation of the $`\pi A\pi \pi A^{}`$ reaction in the near-threshold region.
In early experimental works on $`\pi 2\pi `$ in nuclei it was observed that (i) pions preferentially populate the low-energy part of the kinetic energy spectra , and (ii) the $`\pi ^+\pi ^{}`$ invariant mass $`M_{\pi ^+\pi ^{}}^A`$ peaks increasingly toward the $`2m_\pi `$ threshold as the nucleus mass number increases . The property (i) could not be explained in terms of the $`\pi ^+A\pi ^+\pi ^{}p[A1]`$ phase space . Nor was a detailed model of the $`A(\pi ^+,\pi ^+\pi ^{}`$) reaction able to reproduce the low-energy pion yield, although it embodied the kinematical limits of the experimental apparatus. The same model, however, could correctly predict the total cross section of the $`\pi 2\pi `$ reaction in nuclei , as well as many-fold differential cross sections . A novel approach was then employed to explain the near threshold behaviour of $`M_{\pi ^+\pi ^{}}^A`$ , i.e. the property (ii). In this approach a $`\pi \pi `$ pair is considered as a strongly interacting system when the two pions couple to the I=J=0 quantum numbers. The $`(\pi \pi )_{I=J=0}`$ properties in nuclear matter are studied by dressing the single-pion propagator to account for the $`P`$wave coupling of pions to $`particlehole`$ and $`\mathrm{\Delta }hole`$ configurations. The model is able to explain the general features of the $`M_{\pi ^+\pi ^{}}^A`$ distributions, which are predicted to increasingly accumulate strength near the $`2m_\pi `$ threshold for $`\rho `$, the nuclear medium density, approaching $`\rho _0`$, the saturation density. However, within the same theoretical framework, the near-threshold strength of the $`\pi \pi `$ $`T`$matrix is considerably reduced when the $`\pi \pi `$ interaction is constrained to be chiral symmetric . This may indicate that effects other than the in-medium $`(\pi \pi )_{I=J=0}`$ interaction contribute to the observed strength. Conversely, the absence of any in-medium modification of the $`(\pi \pi )_{I=J=0}`$ interaction leads to $`M_{\pi ^+\pi ^{}}^A`$ distributions which lack of strength at threshold. This is shown in the work of , which compares the experimental results with the model predictions of .
Some $`\pi 2\pi `$ articles were recently published by the CHAOS collaboration at TRIUMF. The novel data highlighted some properties of the in-vacuum $`\pi \pi `$ interaction , as well as the in-medium modifications of the $`\pi \pi `$ interaction . The appearance of the CHAOS results have renewed theoretical interest, which now bases the interpretation of the $`\pi 2\pi `$ and $`\pi \pi `$ data on some common features:
1. The $`\pi \pi `$ interaction is strongly influenced and modified by the presence of the nuclear medium when the $`\pi \pi `$ interaction occurs in the $`(\pi \pi )_{I=J=0}`$ channel, conventionally called the $`\sigma `$channel .
2. Nuclear matter affects the $`(\pi \pi )_{I=2,J=0}`$ interaction only weakly .
3. Models which only include standard many-body correlations, i.e. the $`P`$wave coupling of $`\pi ^{}s`$ to $`ph`$ and $`\mathrm{\Delta }h`$ configurations, are able to explain part of the observed $`M_{\pi ^+\pi ^{}}^A`$ yield near the $`2m_\pi `$ threshold .
In recent theoretical works on the $`(\pi \pi )_{I=J=0}`$ interaction in nuclear matter , the effects of standard many-body correlations have been combined with those derived from the restoration of chiral symmetry in nuclear matter. The resulting $`M_{\pi ^+\pi ^{}}^A`$ distributions regain strength near the $`2m_\pi `$ threshold as $`A`$ (thus the average $`\rho `$) increases. Such behaviour was outlined in some earlier theoretical works, which demonstrated that the $`M_{\pi ^+\pi ^{}}^A`$ enhancement near threshold is a distinct consequence of the partial restoration of the chiral symmetry at $`\rho <\rho _0`$ . The purpose of this article is to present a comprehensive set of $`\pi 2\pi `$ data and to discuss them in the light of the most recent theoretical findings.
Some of the characteristics of the $`\pi 2\pi `$ process in nuclei were presented in Refs. and some observables, those most sensitive to the $`\pi \pi `$ interaction, were discussed in Ref.. These TRIUMF results are presented in the form of many-fold differential cross sections, which are of interest in the current work. Most other $`\pi 2\pi `$ measurements deal only with total cross sections. These include old data taken with emulsion techniques , and more recently, with a magnetic spectrometer similar to CHAOS . Since $`\pi 2\pi `$ total cross sections are not part of the present discussion, these data will not be considered in this work. The article is organised as follows: some features of the experiment are reported in Sec. 2. Sec. 3 deals with the data analysis. The available $`\pi 2\pi `$ models are presented in Sec. 4. The $`\pi 2\pi `$ data are discussed in Sec. 5. Other existing results are mentioned in Sec. 6. Finally, conclusions are summarised in Sec. 7.
2. The experiment
The experiment was carried out at the TRIUMF Meson Facility. Incident pions were produced by the collision of 480 MeV protons on a 10 mm thick graphite target. The M11 pion beam line transported the 282.7 MeV pions to the final focus which was located at about 14 m from the production target. M11 was tuned to focus the pion beam on the target, which was placed at the centre of the CHAOS spectrometer. Pion pairs were detected in coincidence to ensure a unique identification of the pion production process.
2.1 The pion beam, the beam monitor and the beam counting
The M11 pion channel was set to deliver positive pions with a central momentum p=398.5 MeV/c. The beam momentum spread ($`\mathrm{\Delta }`$p/p ) was defined by slits which were located at the intermediate dispersed focus of the channel. The particle beam was mainly composed of pions and protons, $`\pi :p1:10`$. The fraction of positive electrons and muons was negligible, typically $`\pi :\mu :e=1.0:0.003:0.002`$ . The proton contamination was moderated by a CH2 absorber placed at the mid-plane of the M11 channel: the energy deposited by protons in the absorber degraded the central momentum of the proton distribution so that most of them were intercepted by the slits of the channel. Pions were finally discriminated from protons with a plastic scintillator counter placed transversally to the particle beam at about 194 cm upstream of the target. The pulse height ($`PH`$) associated with the passage of a proton through the plastic scintillator was over 3 times higher than the pulse height of a pion of the same momentum, i.e. $`PH_p/PH_\pi 3`$ at 400 MeV/c. This allowed for a fast discrimination between the two beam components and, ultimately, the determination of the beam composition, which was monitored throughout the experiment. The systematic uncertainty in determining the fraction of pion $`f_\pi `$ in the particle beam was about 1%.
The beam monitor ($`M_b`$) consisted of 4 slabs of Pilot U plastic scintillator 3 mm thick, with a combined cross-sectional area of 9$`\times `$26 cm<sup>2</sup>. This area was wide enough to contain the full particle beam. Each segment of the monitor was designed to count roughly the same beam flux, i.e. 1.0-1.5 MHz for $`{}_{}{}^{12}C`$,$`{}_{}{}^{40}Ca`$ and $`{}_{}{}^{208}Pb`$ targets, and $``$0.5 MHz for $`{}_{}{}^{2}H`$ (see Table 1). The four segments were coupled through lucite light guides to four fast-counting tubes in order to avoid piling-up of the output pulses. The beam monitor was placed at about 161 cm from the CHAOS dipole tip edge to prevent any background arising from the interaction of pions and protons with the monitor medium from entering the spectrometer. However, with such an arrangement of the beam monitor, an appreciable fraction of the pions crossing it decayed before reaching the target. The decay rate was $`f_d=8\%`$ and the uncertainty in assessing $`f_d`$ is negligible. The beam momentum spread ($`\mathrm{\Delta }`$p/p at FWHM, \[%\]), total intensity over all segments ($`N`$, \[$`10^6`$particles/s\]) and composition ($`\pi :p:\mu :e`$) at the monitor location are summarised in Table 1.
The beam monitor was only capable of counting the number of incident particle bursts $`N_i`$. Depending on its intensity, a $``$5ns FWHM beam burst may have contained one or more particles, i.e. pions and protons. Therefore the effective number of particles traversing the monitor is $`N_m=\mu _pN_i`$, where $`\mu _p`$ is the average number of particles per beam burst. Since particle counting follows Poisson statistics, the average number of particles is $`\mu _p=ln(1r)/r`$, where $`r=N_i/N_p`$ and $`N_p`$ is the number of primary proton bursts per second 23.06$`\times 10^6`$. In the present measurement the maximum number of particles traversing a monitor single slab did not exceed $`1.5\times 10^6`$/s. Thus, the percentage of particles not counted by the monitor was $`100(\mu _p1)<3.5\%`$, and the uncertainty associated with this correction is negligible.
The number of pions at the target location was finally given by the equation $`N_\pi =\mu _p(1f_d)f_\pi N`$, and the uncertainty in assessing $`N_\pi `$ was slightly above 1%.
2.2 The targets
The targets used in this experiment were either solid self-supporting plates of carbon, calcium, and lead, or a vessel for the liquid deuterium. The targets were located at the centre of the dipole, and their characteristics are summarised in Table 2.
The deuterium target consisted of a low-mass cylindrical vessel filled with liquid deuterium. It was continuously cooled by a refrigerator, and then warmed to its operating temperature, which was monitored both by an external helium gas thermometer and direct measurement of the vapour pressure of the deuterium. The variation of density due to the uncertainty in temperature (23.0$`\pm `$0.2K) was negligible. The variation of density due to the formation of bubbles throughout the liquid was larger, and most significant near the top of the target cell. The variation of thickness of the target cell due to bulging outward of the thin windows contributes a further uncertainty to the target thickness. Both of these effects were convoluted with measured beam profiles to determine the areal target thickness. The total number of scattering centres was evaluated by integrating the beam profile over that of the target, which was determined by tracing elastically scattered pions during a calibration run. The uncertainty in the areal target thickness was estimated to be 3%. In order to evaluate the contribution of the background from the target vessel, a dedicated run was performed with the deuterium target emptied. Then, the rejection of events with the reaction vertex on the vessel reduced the background to about 2%. A reconstructed interaction vertex for a vessel of the same dimensions as the one used in the present experiment is shown in Fig. 2 of Ref..
The solid targets were made of self-supporting plates surrounded by a frame for handling. The frame was made of lucite for the $`{}_{}{}^{12}C`$ and $`{}_{}{}^{208}Pb`$ targets, and alluminum for the $`{}_{}{}^{40}Ca`$ target. The target thicknesses and areas are listed in Table 2. The areal size of the plates was large enough to entirely contain the beam spot. The beam envelope at the target location was measured previously , and was found to be entirely confined to an area of $`40\times 40mm^2`$. This can be seen in Fig. 1,
which illustrates the vertex reconstruction provided by the CHAOS spectrometer from $`\pi ^+\pi ^+\pi ^{}`$ events on $`{}_{}{}^{12}C`$. Analogous results were obtained for the other solid targets.
2.3 The CHAOS spectrometer and triggering system
CHAOS is a magnetic spectrometer which was designed for the detection of multi-particle events in the medium-energy range . The magnetic field is generated by a dipole whose pole tip is 66 cm in diameter. The magnet is capable of producing a field intensity up to 1.6 T with an uniformity of about 1%. There is a 12 cm bore at the centre for the insertion of targets. Fig. 2 illustrates
a reconstructed $`\pi _i^+\pi ^+\pi ^{}p`$ event on $`{}_{}{}^{12}C`$, and shows the geometrical disposition of the wire chambers (WC), the CHAOS first level trigger hardware (CFT), and the magnet return yokes in the corners. WC1 and WC2 are multiwire proportional chambers which are capable of handling rates exceeding 5$`\times 10^6`$ particles/s for extended periods of time with high efficiency ($``$95%). They are cylindrical in shape, with diameters of 22.8 cm and 45.8 cm, respectively. WC3 is a cylindrical drift chamber designed to operate in a magnetic field . Its diameter is 68.6 cm. The outermost chamber, WC4, is a vector drift chamber 122.6 cm in diameter, which operates in the tail of the magnetic field of CHAOS. The segments of WC3 and WC4 which were crossed by the incident particle beam were turned off. The CFT hardware consists of three adjacent cylindrical layers of fast-counting detectors . The first two layers ($`\mathrm{\Delta }E1`$ and $`\mathrm{\Delta }E2`$) are NE110 plastic scintillators 0.3 cm and 1.2 cm thick, respectively. $`\mathrm{\Delta }E1`$ is 72 cm from the magnet centre and spans a zenithal angle of $`\pm 7^{}`$; thus, it defines the geometrical solid angle of CHAOS $`\mathrm{\Omega }`$=1.5 sr. The third layer is a SF5 lead-glass 12.5 cm thick, about 5 radiation lengths, which is used as a Cerenkov counter. The three layers were segmented in order to provide an efficient triggering system for multi-particle events. Each segment covered an azimuthal angle (i.e., in-the-reaction plane) of 18. During data taking two segments were removed in order to allow the passage of the particle beam (see Fig. 2).
The experiment relied on two on-line triggers: the first ($`1^{st}LT`$) and second ($`2^{nd}LT`$) level trigger. The logical signals which were issued by the CFT following $`\pi ^+\pi ^+\pi ^\pm `$ events were initially pipe-lined and then filtered by requiring a coincidence $`M_b\times (\mathrm{\Delta }E1\times \mathrm{\Delta }E2)_i\times (\mathrm{\Delta }E1\times \mathrm{\Delta }E2)_j`$ with $`i,j=1,..,18`$ and $`ij`$. Such a logic analysis was performed at a rate of about 30 MHz. Despite of the fast filtering, the $`1^{st}LT`$ trigger was overwhelmed ($`>`$99.99%) by multi-particle events from competing reactions: quasi-elastic scattering (qes) $`\pi ^+\pi ^+p(d)`$, and pion absorption (abs) $`\pi ^+pp`$. This called for a second and more selective triggering system, which was based on the information delivered by the innermost wire chambers. For each event passed by the $`1^{st}LT`$, the $`2^{nd}LT`$ trigger calculated the momentum of each track, its polarity and the reaction vertex topology. It then compared a combination of these parameters with predefined criteria. When the criteria were fulfilled the event was accepted and recorded on tape, otherwise a fast clear signal was sent to the readout electronics, and the two triggers were reset. An entire $`2^{nd}LT`$ cycle was accomplished in $`<10\mu `$s, depending on the number of reconstructed tracks. The event rate recorded on tape ranged from 30 to 100 Hz depending on the nucleus studied. Roughly 0.1% of the recorded events were $`\pi \pi \pi `$ as determined by detailed off-line analysis.
2.4 The CHAOS acceptance and performance
In the present measurement CHAOS was operated with a magnetic field of 0.5 T. The field and the energy deposited by the particles emerging from the target in the WC’s and $`\mathrm{\Delta }E1`$ media determined the CHAOS threshold, which was 57 MeV/c for pions. The pion threshold can be observed in Fig. 3, which shows a diffusion plot of the $`\pi ^\pm `$ momenta ($`p_\pi `$) versus their azimuthal angles ($`\mathrm{\Theta }_\pi `$), for the $`\pi ^+\pi ^+\pi ^{}`$ reaction channel on $`{}_{}{}^{12}C`$.
The figure also depicts the overall CHAOS acceptance for pions. The removal of the two CFT segments intercepting the particle beam defines strips in the ($`p_\pi ,\mathrm{\Theta }_\pi `$) plane which are unreachable by pions. These strips are confined within bounds (continuous lines), which are the results of GEANT Monte Carlo simulations. The strips are 14 wide, 4 narrower than a CFT segment. The narrowing is due to the circular trajectories that pions acquire by moving in the magnetic field of CHAOS, which ultimately defines the exit direction of pions before impinging upon $`\mathrm{\Delta }E1`$. For both positive and negative pions the removal of the CFT segment which allows the particle beam to enter CHAOS (strip at around 180) negligibly affects the reaction phase space. The missing CFT segment at the exit (strip at around 0) somewhat restricts the available $`\pi \pi \pi `$ phase space. However, the limitation is only for an angular interval of 14.
The acceptance of CHAOS is irregular in the ($`p_\pi `$,$`\mathrm{\Theta }_\pi `$) plane due to the finite segmentation of the CFT, the two missing CFT segments, and the pion decays inside CHAOS. All these sources of non-uniformity were accounted for by assigning a weight to each $`\pi \pi \pi `$ event. The weights were determined using GEANT Monte Carlo simulations, as described in more detail in Ref.. Weight distributions are represented in Fig. 4 for the two $`\pi ^+\pi ^+\pi ^\pm `$ reaction
channels on $`{}_{}{}^{12}C`$. These distributions show a similar behaviour: they are peaked around 1.5 and have a tail which extends up to 8. Since each $`\pi \pi `$ event was binned with its weight, to avoid large corrections in the distributions, weights were restricted to vary from $`>`$0 to $`\mu +2\sigma `$, where $`\mu `$ is the mean value and $`\sigma `$ is the standard deviation of the weight distribution. For $`\pi ^+\pi ^+`$ events, $`\mu (\sigma )`$ = 2.43(1.42), while for $`\pi ^+\pi ^{}`$ events, $`\mu (\sigma )`$ = 2.52(1.31).
As an example, in Fig. 5 are reported the weight distributions used to correct the measured invariant mass spectra for the two $`\pi ^+\pi ^+\pi ^\pm `$ reaction channels on $`{}_{}{}^{12}C`$. The average uncertainty in the weight is $``$6% . The distributions are represented over an invariant mass interval where the yields of the $`M_{\pi \pi }`$ distributions are different from zero (see Fig. 9). The weight distributions appear rather similar and flat, especially over the shorter interval from 280 MeV to 395 MeV, where the $`M_{\pi ^+\pi ^\pm }^C`$ yields are not negligible. The flatness of the weight distributions implies that the shape of the measured $`M_{\pi ^+\pi ^\pm }^C`$ distributions changes only moderately when corrected for the CHAOS acceptance. The same conclusion applies to the other nuclei.
The angular and momentum resolutions were measured during the CHAOS commissioning, which just preceded the present measurement. The angular resolution was $`<0.5^{}`$, and the momentum resolution was 1% ($`\sigma `$) for a 1 T magnetic field and 225 MeV/c pions . However, for the pion-production experiment, the pion mean momentum was $``$130 MeV/c (see Table 5), and the CHAOS field was set at 0.5 T. Simulations show that under these conditions, the angular resolution, which was dominated by the multiple scattering, is below 2, and the momentum resolution is $``$4% ($`\sigma `$) for 130 MeV/c pions.
2.5 Particle identification
In the case of pion-production studies, pions must be selectively separated from other reaction byproducts, essentially from $`e,p`$ and $`d`$, because the $`\pi \pi \pi `$ reaction cross section in nuclei is from 2 to 3 orders of magnitude lower than the cross section of other competing reactions. For instance, at $`T_{\pi ^+}`$ 280 MeV and for $`{}_{}{}^{40}Ca`$, the quasi-elastic scattering process $`\pi ^+\pi ^+p`$ has a total cross section $`\sigma _{qes}800\sigma _{\pi ^+\pi ^+}`$, while the pion absorption process $`\pi ^+pp`$ has a total cross section $`\sigma _{abs}600\sigma _{\pi ^+\pi ^+}`$. For other processes like single charge exchange (scx) $`\pi ^+\pi ^{}`$, the background arises from the $`\pi ^{}\gamma \gamma `$ decay followed by the $`\gamma `$ conversion $`\gamma e^+e^{}`$ in the target or in the material which surrounds the target. In this case the $`e^+e^{}`$ pair may emulate a $`\pi ^+\pi ^{}`$ pair.
Particles were identified by combining the information provided by WC’s (particle momentum and polarity) and CFT’s (pulse heights). For reactions with protons (deuterons) in the final state, i.e. qes and abs, the particle separation was easily achieved through the different response function of $`\mathrm{\Delta }E1`$ and $`\mathrm{\Delta }E2`$ to $`\pi `$’s and p’s (d’s). Pions from the pion-production process have a momentum distribution which does not exceed 210 MeV/c (see Fig. 14). At this momentum, $`PH_p/PH_\pi 6`$, and increases as the momentum decreases, see Fig. 8 in Ref. and Fig. 3 in Ref.. However, the same method of particle separation could not successfully be applied to $`\pi `$’s and e’s for momenta exceeding 140 MeV/c. Thus, Cerenkov counters were used. These counters were capable of separating $`\pi `$’s from e’s with an efficiency of about 95% for particle momenta up to about 230 MeV/c Fig. 11. This value increased slightly ($``$98%) when the Cerenkov pulse heights were combined with the $`\mathrm{\Delta }E1`$ and $`\mathrm{\Delta }E2`$ pulse heights. Finally, the $`\pi `$ to p and $`e`$ discrimination turned out to be $``$100% when the following soft kinematic cuts were applied to particle momenta and angles:
Protons which passed the pion $`PH`$ test were definitively rejected by restricting their momenta below 210 MeV/c, the maximum momentum available to pions. Such a momentum was slightly above the CHAOS threshold to protons, 185 MeV/c.
Electron pairs, which might emulate $`\pi ^+\pi ^{}`$ pairs, appeared with opening angles $`<2^{}`$. This is in accord with simulations of the $`\pi ^{}\gamma \gamma `$ decay followed by the $`\gamma `$ conversion $`\gamma e^+e^{}`$ in CHAOS. Therefore a 2 cut was applied during the analysis, which resulted in a $``$0.3% reduction of the $`\pi 2\pi `$ in-the-reaction plane phase space.
3. Analysis
The data reduction was based on fully reconstructed $`\pi ^+\pi ^+\pi ^{}`$ and $`\pi ^+\pi ^+\pi ^+`$ events. Table 3 reports the yields for each channel and nucleus studied. In order to form differential cross sections these events were binned with their weights, which include the $`\pi ^+\pi ^\pm `$ decay rates inside CHAOS. Some particularly useful distributions are represented included those of the canonical observables such as invariant mass $`M_{\pi \pi }`$, angles $`\mathrm{\Theta }_\pi `$, $`\mathrm{\Theta }_{\pi \pi }`$, $`cos\mathrm{\Theta }_{\pi \pi }^{CM}`$, and kinetic energy $`T_\pi `$.
The ability to measure the kinetic energies ($`T`$) and laboratory angles ($`\mathrm{\Theta }`$) of both final state pions with the CHAOS spectrometer permitted the determination of the five-fold differential cross section $`^5\sigma /(T\mathrm{\Theta })_{\pi _1}(T\mathrm{\Theta })_{\pi _2}\mathrm{\Phi }_{\pi _1\pi _2}`$. Here, $`\mathrm{\Phi }_{\pi _1\pi _2}`$ is the zenithal angle between $`\pi _1`$ and $`\pi _2`$, which, in the CHAOS detector, can be measured at either $`180^{}\pm 7^{}`$, or $`0^{}\pm 7^{}`$. Four-fold differential cross sections were then obtained by integrating out the $`\mathrm{\Phi }_{\pi \pi }`$ dependence. This was accomplished by using a linear function joining the two measured data points, and the result was the factor $`f_\mathrm{\Phi }`$. This method of assessing $`f_\mathrm{\Phi }`$ leads to a systematic uncertainty of 8% ($`\sigma `$) for deuterium and 6-7% ($`\sigma `$) for nuclei. The data were also reduced to single differential cross sections $`d\sigma /d𝒪_\pi `$, where $`𝒪_\pi `$ represents $`(T_\pi `$ or $`\mathrm{\Theta }_\pi )_{1,2}`$. The cross section was then related to measured quantities $`\frac{d\sigma }{d𝒪_\pi }`$= $`f_e\frac{N(𝒪_\pi )}{\mathrm{\Delta }𝒪_\pi }`$, where $`f_e`$ is a parameter which is determined by the experimental conditions, $`N(𝒪_\pi )`$ is the number of weighted events in a given bin and $`\mathrm{\Delta }𝒪_\pi `$ is the bin width for the $`𝒪_\pi `$ observable. The experimental parameter $`f_e`$ is the product of several factors which were individually measured in order to obtain absolute values for the cross sections: $`f_e=f_\mathrm{\Phi }(N_\pi 𝒞_l_{WC}N_{sc})^1`$, where $`N_\pi `$ is the number of pions impinging upon the target (see discussion in Paragraph 2.1), $`𝒞_l`$ is the computer live-time during data acquisition which was 61% for $`{}_{}{}^{2}H`$, 53% for $`{}_{}{}^{12}C`$, 70% for $`{}_{}{}^{40}Ca`$ and 80% for $`{}_{}{}^{208}Pb`$, $`_{WC}`$ is the overall WC efficiency 49.5%, and $`N_{sc}`$ is the number of target scattering centres per unit area (see discussion in Paragraph 2.2). The uncertainty in $`f_e`$ $``$11% ($`\sigma `$) comes mostly from the uncertainties in $`_{WC}`$ and $`f_\mathrm{\Phi }`$. In fact, $`_{WC}=_{WC1}_{WC2}(_{WC1}_{WC2}_{WC3})^2`$, and the average uncertainty in measuring the efficiency of a wire chamber was 0.93% thus yielding an overall $`_{WC}`$ uncertainty of 8.3%. The WC4 efficiency uncertainty was 0.0% since only three out of eight independent anodes of one cell were required to fire in order to determine the direction of an outgoing pion. Finally, the overall uncertainty was obtained by summing the uncertainty due to $`f_\mathrm{\Phi }`$ and $`_{WC}`$ in quadrature.
4. Models of the $`\pi 2\pi `$ reaction and the $`\pi \pi `$ interaction in the nuclear medium
This experiment employed the $`\pi ^+A\pi ^+\pi ^\pm A^{}`$ reaction ($`\pi 2\pi `$) to study the near-threshold $`\pi \pi `$ interaction as a function of the (average) nuclear density. Other reactions can, in principle, be used for similar studies, eg. the $`p\pi \pi `$ and $`\gamma \pi \pi `$ reactions. For any reaction and for a reliable interpretation of the $`\pi \pi `$ dynamics in nuclear matter, theoretical models should account for both the reaction mechanism, and nuclear distortions (i.e. pion absorption and others). In the case of $`\pi 2\pi `$, this approach was followed in the work of Refs. , which also included the kinematical limits of the experimental apparatus. For those theoretical works which addressed only the dynamics of a $`(\pi \pi )_{I=J=0}`$ interacting system in nuclear matter, it was useful to define the observable $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ (see Paragraph 5.5) since it is only weakly dependent on the reaction mechanism and nuclear distortions .
Before discussing the data and the comparison with the available theoretical calculations for both the $`\pi 2\pi `$ reaction and the behaviour of a $`\pi \pi `$ interacting system in nuclear matter, some details of the models are presented below.
In order to understand the main features of medium modification on the $`\pi \pi `$ interacting system , the pion-production reaction on nuclei is modelled by means of the two leading reactions: 1) the one-pion exchange (OPE) reaction, $`\pi ^+N\pi ^+\pi ^\pm N`$, which contributes to both the isoscalar and isotensor channels, and 2) the $`N^{}(1440)`$ resonance excitation which is restricted to decay only to isoscalar $`\pi \pi `$ states, $`\pi ^+NN^{}ϵNN(\pi \pi )_{I=J=0}`$. The in-vacuum $`\pi \pi `$ interaction is accounted for by the chirally-improved Jülich model, which in the GeV region is able to describe the $`\pi \pi \pi \pi `$ scattering and to reproduce the $`I=J=0`$ $`\pi \pi `$ phase shifts. The $`S`$wave $`\pi \pi `$ (final state) interaction is modified in the nuclear medium through standard renormalization of the pion progagator (see Fig. 6) which determines the following: in the $`\pi ^+\pi ^{}`$ channel, the isoscalar $`\pi \pi `$ amplitude is strongly reshaped which finally provides the near-threshold $`M_{\pi ^+\pi ^{}}`$ enhancement observed in the CHAOS data; in the $`\pi ^+\pi ^+`$ channel, the nuclear density has little effect on the pure isotensor $`\pi \pi `$ amplitude. In order to have a realistic comparison with the CHAOS results, the model deals with common nuclear effects like Fermi motion and pion absorption although using approximate approaches, and calculates the many-fold differential cross sections by accounting for the CHAOS acceptance.
The effort in Ref. is to model the $`\pi ^+A\pi ^+\pi ^\pm A^{}`$ reaction through a microscopic description of the elementary $`\pi N\pi \pi N`$ reaction and a detailed study of the production reaction in nuclei.
The elementary pion-production reaction relies on five Feynmann diagrams, which have $`N`$’s, $`\mathrm{\Delta }`$’s and $`N^{}`$’s as intermediate isobars. The relative scattering amplitudes are derived from chiral lagrangian. The purely mesonic $`\pi \pi `$ amplitude is calculated by means of the coupled-channel Bethe-Salpeter equation, and the resulting amplitude is capable of predicting the experimental phase shifts of the scalar-isoscalar channel up to 1.2 GeV. With this detailed approach to the $`\pi N\pi \pi N`$ reaction, the model is able to predict the total cross sections for the $`\pi ^+\pi ^{}`$, $`\pi ^{}\pi ^{}`$ and $`\pi ^+\pi ^+`$ channels from threshold up to 300 MeV. In the nuclear medium, the $`(\pi \pi )_{I=J=0}`$ interaction is strongly modified due to the coupling of pions to $`ph`$ and $`\mathrm{\Delta }h`$ excitations, which are represented by the diagrams in Fig. 6. This results in a displacement of the $`ImT_{\pi \pi }`$ strength toward the $`2m_\pi `$ threshold as the nuclear density increases. Conversely, the nuclear medium weakly affects the $`(\pi \pi )_{I=2,J=0}`$ interaction.
The model includes several nuclear effects: Fermi motion, Pauli blocking, pion absorption and quasi-elastic scattering. The last two interactions are found to largely remove $`\pi 2\pi `$ pions from the incident flux, which ultimately confines the pion-production process on the nuclear skin. Fig. 7 shows the probability of a $`\pi 2\pi `$ event to occur ($`\mathrm{\Pi }_{\pi 2\pi }`$, dotted curve) along with the nuclear density distribution ($`\rho `$, solid curve) for $`Ca`$ as a function of the nuclear radius. The two distributions yield a weighted mean nuclear density of $`\rho `$ =0.24$`\rho _n`$, which qualitatively agrees with the experimental result reported in, where the interacting nucleon of the $`\pi N[A1]\pi \pi N[A^{}1]`$ reaction is found to preferentially lie in an external nuclear orbit. The above value of $`\rho `$ can also be used for comparison with the nuclear densities quoted in Ref. and discussed in the Paragraph 5.1. Finally, the model employs a Monte Carlo technique to generate and propagate $`\pi 2\pi `$ events, thus their phase space can easily accommodate the CHAOS acceptance.
The work of Ref. aims at studing the possibility of the $`\sigma `$meson identification in nuclear matter, where the elusive particle might be detected as a consequence of the partial restoration of chiral symmetry. Due to the strong interaction of $`\sigma `$ with the nuclear medium, the description of the $`\sigma `$ properties in terms of parameters like $`m_\sigma `$ and $`\mathrm{\Gamma }_\sigma `$ appears inadequate, and a proper observable becomes the $`\sigma `$ spectral function $`\rho _\sigma `$. By using a simple but general approach, proves that in the proximity of the $`2m_\pi `$ threshold the partial restoration of the chiral symmetry implies $`\rho _\sigma (\omega 2m_\pi )=\pi ^1/Im\mathrm{\Sigma }_\sigma (\omega ,\rho )`$, where $`\mathrm{\Sigma }_\sigma `$ is the $`\sigma `$ self-energy, $`\omega `$ is the total energy and $`\rho `$ the nuclear density. Near $`2m_\pi `$ threshold, the $`Im\mathrm{\Sigma }_\sigma `$ is proportional to the phase space factor $`(1\frac{2m_\pi ^2}{\omega ^2})^{1/2}`$, thus $`\rho _\sigma (\omega 2m_\pi )\theta (\omega 2m_\pi )/(1\frac{2m_\pi ^2}{\omega ^2})^{1/2}`$. A direct consequence is that the $`\sigma `$ spectral function strength enhances as the $`\sigma `$ total energy approaches $`2m_\pi `$. In order to render this general finding more quantitative, the theory uses the SU(2) linear sigma model. The mean field correction in nuclear matter of $`\mathrm{\Sigma }_\sigma (\rho )`$ is accounted for by the leading diagram sketched in Fig. 8. Furthermore, the model does not include collective pionic modes to derive a possible source of near-threshold strength from the pion coupling to $`ph`$ and $`\mathrm{\Delta }h`$ correlated states.
This work studies the $`(\pi \pi )_{I=J=0}`$ interaction in nuclear matter by constructing a microscopic theory for both the $`\sigma `$meson propagator ($`D_\sigma `$) and the $`\pi \pi `$ $`T`$matrices in the framework of the linear sigma model. In this approach the basic chiral symmetry constraints, i.e. the vanishing of the scattering length in the chiral limit and the Ward’s identities are satisfied, and in addition, the theory is capable of reproducing the experimental phase shift in the scalar-isoscalar channel. In-medium pions are renormalized through their $`P`$wave coupling to correlated $`ph`$ and $`\mathrm{\Delta }h`$ states, as depicted in Fig. 6. Bare sigmas, of the linear sigma model, are coupled to the nuclear medium via the tad-pole diagram, which is sketched in Fig. 8. The $`P`$wave renormalization of sigmas leads to a downward shifts of the $`\sigma `$ mass, which ultimately produces a strong enhancement of the $`\sigma `$meson mass distribution (i.e. $`ImD_\sigma `$) around the $`2m_\pi `$ threshold. This enhancement is further increased by a factor of 2 to 3 by the $`S`$wave renormalization of the bare $`\sigma `$meson mass (i.e. by the partial restoration of chiral symmetry), which presents a similar behavior near the $`2m_\pi `$ threshold.
5. Results of the $`\pi \pi \pi `$ reaction in nuclei
A general property of the $`\pi 2\pi `$ process on nuclei in the low-energy $`M_{\pi \pi }`$ regime was outlined by previous experimental works: it is a quasi-free process both when it occurs on deuterium and on complex nuclei . Furthermore, a common reaction mechanism underlies the process whether it occurs on a nucleon or a nucleus . Thus the study of the $`\pi ^+`$ $`{}_{}{}^{2}H\pi ^+\pi ^\pm NN`$ reaction is dynamically equivalent to studying the elementary $`\pi ^+n\pi ^+\pi ^{}p`$ and $`\pi ^+p\pi ^+\pi ^+n`$ reactions separately.
In the present measurement, the $`\pi ^+A\pi ^+\pi ^\pm A^{}`$ reactions were studied under the same experimental conditions. Thus for a given observable the distributions are directly comparable. The error bars explicitly shown on the spectra represent the statistical uncertainties, which must be added in quadrature to the systematic one $``$11% ($`\sigma `$) to obtain the overall uncertainty associated with the data points. The only exception is for the $`cos\mathrm{\Theta }_{\pi \pi }`$ distributions in Fig. 10, where the overall uncertainties are shown.
5.1 The $`\pi \pi `$ invariant mass
Fig. 9 shows the single differential cross sections (diamonds) as a function of the $`\pi \pi `$ invariant mass ($`M_{\pi \pi }`$, MeV) for the two reaction channels $`\pi ^+\pi ^+\pi ^{}`$ and $`\pi ^+\pi ^+\pi ^+`$ . Horizontal error bars are not indicated since they lie within symbols. The distributions span the total energy interval available to the $`\pi 2\pi `$ reaction, which ranges from $`2m_\pi `$, the low-energy threshold, up the 420 MeV, the maximum allowed by the reaction. The $`\pi A\pi \pi N[A1]`$ phase space simulations (dotted histograms for $`A:`$ $`{}_{}{}^{2}H`$, $`{}_{}{}^{12}C`$, $`{}_{}{}^{40}Ca`$ and $`{}_{}{}^{208}Pb`$) are also shown, and are normalized to the same area as the experimental distributions.
Regardless of the nuclear mass number, the invariant mass distributions for the $`\pi ^+\pi ^+\pi ^+`$ channel closely follow phase space, and the energy maximum increases with increasing $`A`$, that is, with increasing nuclear Fermi momentum. The $`\pi ^+\pi ^+\pi ^{}`$ channel displays a different behaviour. Compared to phase space, the $`{}_{}{}^{2}H`$ invariant mass displays little strength from $`2m_\pi `$ to 310 MeV, while, in the same energy interval, the $`{}_{}{}^{12}C`$, $`{}_{}{}^{40}Ca`$ and $`{}_{}{}^{208}Pb`$ $`\pi ^+\pi ^{}`$ invariant mass distributions increasingly peak as $`A`$ increases.
In order to explain the nature of the reaction mechanism contributing to the peak structure, it is useful to examine cos$`\mathrm{\Theta }_{\pi \pi }^{CM}`$ distributions and $`(p_{\pi ^+},p_\pi ^{})`$ diffusion plots for those events with invariant masses in the region of the peak. $`\mathrm{\Theta }_{\pi \pi }^{CM}`$ is the angle between the direction of a final pion and the direction of the incoming pion beam in the $`\pi ^+\pi ^{}`$ rest frame. Fig. 10 shows the cos$`\mathrm{\Theta }_{\pi \pi }^{CM}`$ distributions (diamonds) for $`2m_\pi M_{\pi ^+\pi ^{}}310`$ MeV and $`310<M_{\pi ^+\pi ^{}}420`$ MeV, the latter being shown for comparison. The vertical error bars represent the overall uncertainties, which are the systematic and statistical uncertainties summed in quadrature. The solid lines in Fig. 10 represent the best-fit to the differential distributions, which was obtained with a partial wave expansion limited to the three lowest waves, i.e. S, P and D. The best-fit results are reported in Table 4, along with $`\chi _\nu ^2`$, which was evaluated using the overall uncertainties. For all the nuclei studied $`\chi _\nu ^21`$ which indicates that a proper number of waves was used in the expansion. In the case of heavier nuclei, the $`\pi ^+\pi ^{}`$ system predominantly couples in $`S`$wave $``$ 95% (or $`\mathrm{}`$=0 relative angular momentum) and a remaining 5% is spent in a $`D`$wave ($`\mathrm{}`$=2) state. Furthermore, within the sensitivity of the $`\chi _\nu ^2`$method, any $`P`$wave ($`\mathrm{}`$=1) coupling of the two pions is excluded.
Fig. 11 displays the diffusion plots of final pion momenta for the $`\pi ^+\pi ^+\pi ^{}`$ reaction channel. For the plots on the left side, $`M_{\pi ^+\pi ^{}}`$ is constrained to the peak region, i.e. $`2m_\pi M_{\pi \pi }`$ 310 MeV. The plots on right side, reported for comparison, are for invariant masses in the higher energy region, 310$`<M_{\pi \pi }`$420 MeV. In the peak region, $`{}_{}{}^{2}H`$ has the pion momenta flatly diffused over the available phase space, while for the more complex nuclei, the choice of $`M_{\pi ^+\pi ^{}}`$310 MeV highlights a bump structure in the momentum range 90-140 MeV/c. Generally, the ($`p_{\pi ^+},p_\pi ^{}`$) diffusion plots on the right side have a peak-structure which is somewhat wider and centered at higher pion momenta.
The results of two recent theoretical works, $`1`$ and $`2`$ , which have modelled the $`\pi 2\pi `$ reaction on nuclei, are reported in Fig. 12. The (short and long) dashed lines denote the $`1`$ calculations while the $`2`$ predictions are shown as full lines. In the case of $`Ca`$, $`R1`$ ($`R2`$) indicates the $`1`$ predictions for $`\rho `$=0.7$`\rho _n`$ ($`\rho `$=0.5$`\rho _n`$), where $`\rho _n`$ is the nuclear saturation density, while $`V1`$ is the result of the $`2`$ calculations for a mean $`\rho `$=0.24$`\rho _n`$. For the purpose of the present discussion, the curves in Fig. 12 are normalized to the experimental data.
For the $`\pi ^+\pi ^+\pi ^+`$ channel, the $`R1`$ and $`R2`$ distributions differ slightly from each other. The $`R1`$ distribution is broader, due to the larger nuclear Fermi momentum, which is a consequence of the higher nuclear density used . The lower $`\rho `$ used for $`R1`$ seems to describe the measured distribution better. This trend is also supported by $`V1`$ whose $`2`$ predictions agree well with the invariant mass distributions. Therefore, both $`1`$ and $`2`$ indicate that the average nuclear density of $`{}_{}{}^{40}Ca`$ for which the production reaction takes place cannot exceed 0.5$`\rho _n`$.
In the $`\pi ^+\pi ^+\pi ^{}`$ channel the distribution predicted by $`2`$ for $`{}_{}{}^{2}H`$ is able to describe the data, although it tends to overestimate the low-energy $`M_{\pi \pi }`$ yield. The present version of the model is considerably improved compared to previous versions , thus placing the construction of nuclear medium effects on more reliable grounds. The $`2`$ approach to $`{}_{}{}^{40}Ca`$ includes several medium effects: Fermi motion, pion absorption, pion quasi-elastic scattering, and $`(\pi \pi )_{I=J=0}`$ medium modifications. Nevertheless, the model is unable to reproduce the observed $`M_{\pi \pi }`$ strength in the near threshold region. An increase of $`\rho `$ from 0.24$`\rho _n`$ to 0.5$`\rho _n`$ and to 0.7$`\rho _n`$ is unlikely to improve the agreement with the experimental cross section, see also Fig. 9 of Ref.. In the case of $`1`$, the $`R2`$ prediction (0.5$`\rho _n`$) seems to reproduce the $`M_{\pi \pi }`$ distribution better, although some of the near-threshold yield already comes from the $`\pi ^+`$ $`{}_{}{}^{2}H\pi ^+\pi ^{}pp`$ production. The model, in fact, overestimates the cross section at low invariant masses whose amount is shown in Fig. 12. The $`R1`$ solution (0.7$`\rho _n`$) provides the correct $`M_{\pi \pi }`$ near-threshold intensity but does not describe the remaining part of the distribution very well. Therefore, $`V1`$, $`R1`$ and $`R2`$ suggest that the missing $`M_{\pi ^+\pi ^{}}`$ strength near the 2m<sub>π</sub> threshold should be searched for in a stronger $`\rho `$-dependence of the $`(\pi \pi )_{I=J=0}`$ interaction, rather than by requiring the pion-production process to occur in an unlikely high-density nuclear environment.
5.2 The $`\pi \pi `$ opening angle
For the $`\pi A\pi _1\pi _2A^{}`$ production reaction, the relation between the $`\pi _1\pi _2`$ invariant mass and the directly measured quantities is: $`M_{\pi _1\pi _2}^2`$ = 4$`m_\pi ^2+2T_{\pi _1}T_{\pi _2}+2m_\pi (T_{\pi _1}+T_{\pi _2})p_{\pi _1}p_{\pi _2}cos\mathrm{\Theta }_{\pi _1\pi _2}`$, where $`T_\pi (p_\pi )`$ is the pion kinetic energy (momentum) and $`\mathrm{\Theta }_{\pi _1\pi _2}`$ is the opening angle of the pion pair. The behaviour of $`\mathrm{\Theta }_{\pi _1\pi _2}`$ as a function of $`M_{\pi _1\pi _2}`$ is studied by directly applying soft-cuts on $`M_{\pi _1\pi _2}`$. Fig. 13 illustrates the differential cross section as a function of $`\mathrm{\Theta }_{\pi \pi }`$ for $`2m_\pi M_{\pi \pi }`$310 MeV (crosses), and $`2m_\pi M_{\pi \pi }`$ 420 MeV (diamonds). The differential cross sections nearly coincide for all nuclei from 0 up to 60, which indicates that the near-threshold pion pairs are preferentially emitted at small opening angles, regardless of both the pion energy and the reaction channel. In general, the $`\mathrm{\Theta }_{\pi \pi }`$ distributions for the two reaction channels vary markedly with $`A`$, making this observable suitable for selective tests of various $`\pi 2\pi `$ models.
5.3 The $`\pi `$ energy
Fig. 14 shows the differential cross sections for the two $`\pi ^+A\pi ^+\pi ^\pm A^{}`$ reactions as a function of the kinetic energy of each final pion. The weighted averages of the distributions are reported in Table 5. $`𝒞`$<sub>+</sub> and $`𝒞`$<sub>-</sub> refer to the centroids of the distributions, in MeV, for positive and negative pions, respectively. The $`𝒞`$<sub>+</sub>’s display a weak $`A`$dependence. For the complex nuclei the mean energy is 49.0$`\pm `$1.0 MeV which is similar to the mean energy of 51.6$`\pm `$0.9 MeV for $`{}_{}{}^{2}H`$. This can be explained by observing that the nuclear Coulomb repulsion on final pions is in part compensated by the same repulsion on initial pions. In the $`\pi ^+\pi ^+\pi ^{}`$ channel, centroids follow a Coulomb-like behaviour, although the intensity is smaller than expected. In fact, $`𝒞`$<sub>-</sub> decreases by 2.2 MeV from $`{}_{}{}^{12}C`$ to $`{}_{}{}^{40}Ca`$, and by 1.4 MeV from $`{}_{}{}^{40}Ca`$ to $`{}_{}{}^{208}Pb`$, while a semi-classical treatment of the nuclear Coulomb attraction on point-like particles yields a decrease of 2.2 MeV and 4.6 MeV, respectively.
In Fig. 14 the experimental $`T_\pi `$ distributions are shown, along with results from $`\pi ^+A\pi ^+\pi ^+n[A1]`$ phase space simulations. Except for $`{}_{}{}^{2}H`$, the experimental distributions are well reproduced by the simulations, especially at higher energies, thus giving confidence in the reliability of the data. The result of the simulations cannot be applied to negative pions, since the energy losses of $`\pi ^{}`$’s due to Coulomb interactions is not exactly known. The kinetic energy distributions are not compared to the results of any theoretical predictions because none are available at this time.
5.4 The $`\pi `$ angles
In order to complete the discussion of the canonical observables, in Fig. 15 the cos$`\mathrm{\Theta }_\pi `$ differential cross sections are reported, where pion angles are in the laboratory frame and cos$`\mathrm{\Theta }_\pi `$=1 coincides with the direction of the incident pion. Positive and negative pion distributions are denoted with full and open diamonds, respectively. The interpretation of these distributions awaits microscopic $`\pi 2\pi `$ calculations, since phase space simulations cannot describe the observed bump structures.
5.5 The $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ ratio
In this paragraph the observable $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ is presented in comparison with recent theoretical predictions. $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ is defined as the composite ratio $`\frac{M_{\pi \pi }^A}{\sigma _T^A}/\frac{M_{\pi \pi }^N}{\sigma _T^N}`$, where $`\sigma _T^A`$ ($`\sigma _T^N`$) is the measured total cross section of the $`\pi 2\pi `$ process in nuclei (nucleon). This observable has the property of yielding the net effect of nuclear matter on the $`(\pi \pi )_{I=J=0}`$ interacting system regardless of the $`\pi 2\pi `$ reaction mechanism used to produce the pion pair . Therefore, $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ can be compared not only with the $`1`$ and $`2`$ predictions which explicitly calculate both $`M_{\pi \pi }^{Ca}`$ and $`M_{\pi \pi }^{{}_{}{}^{2}H}`$, but also with the theories described in $`3`$ and $`4`$ because they calculate the mass distribution of an interacting $`(\pi \pi )_{I=J=0}`$ system (i.e. $`ImD_\sigma `$) both in vacuum and in nuclear matter. Since the calculations are reported either in arbitrary units or in units which are complex to scale , theoretical predictions are normalized to the experimental distributions at $`M_{\pi \pi }`$=350$`\pm `$10 MeV, where the $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ distribution is flat.
For both reaction channels, the full and dotted curves in Fig. 16 are obtained by simply dividing $`M_{\pi \pi }^{Ca}`$/ $`M_{\pi \pi }^{{}_{}{}^{2}H}`$. It is worthwhile recalling that for the option $`\rho `$=0.5$`\rho _n`$ is used while for the mean density is $`\rho `$=0.24$`\rho _n`$. Furthermore, for both approaches the underlying medium effect is the $`P`$wave coupling of $`\pi `$’s to $`ph`$ and $`\mathrm{\Delta }h`$ configurations, which accounts for the near-threshold enhancement. When applied to the $`𝒞`$$`{}_{}{}^{Ca}{}_{\pi \pi }{}^{}`$, both $`1`$ and $`2`$ predict the same result, which describes the behaviour of $`𝒞`$$`{}_{}{}^{Ca}{}_{++}{}^{}`$ fairly well throughout most of the $`M_{\pi \pi }`$ energy range, but only reproduces a small part of the near-threshold strength for $`𝒞`$$`{}_{}{}^{Ca}{}_{+}{}^{}`$.
$`3`$ and $`4`$ examine the medium modifications on the scalar-isoscalar meson, the $`\sigma `$meson. Nuclear matter is assumed to partially restore chiral symmetry, and consequently $`m_\sigma `$ is assumed to vary with $`\rho `$. The variation is parametrized as $`1p\frac{\rho }{\rho _n}`$, where $`p`$ is 0.1$`p`$0.3 for $`3`$, and 0.2$`p`$0.3 for $`4`$, and for both is $`\rho `$=$`\rho _n`$. Both models are capable of yielding large strength near the $`2m_\pi `$ threshold, therefore the results for $`𝒞`$$`{}_{}{}^{A}{}_{+}{}^{}`$ are compared for a common minimum value of the parameter $`p`$=0.2. In addition, the choice $`\rho `$=$`\rho _n`$ may result appropriate for $`{}_{}{}^{208}Pb`$ but it is not for medium (i.e. $`{}_{}{}^{40}Ca`$) and light (i.e. $`{}_{}{}^{12}C`$) nuclei. In Fig. 16 the predictions of $`3`$ and $`4`$ are shown as the dash-dotted curves and dashed curves, respectively. $`4`$ predicts a larger near-threshold strength, which is due to the combined contributions of the in-medium $`P`$wave coupling of pions to $`ph`$ and $`\mathrm{\Delta }h`$ configurations, and to the partial restoration of chiral symmetry in nuclear matter. These two models, however, are still too schematic for a conclusive comparison to the present data, and full theoretical calculations are called for.
6. Existing $`\pi 2\pi `$ results
Novel $`\pi 2\pi `$ results were recently presented by the Crystal Ball (CB) Collaboration at the AGS. The $`\pi ^{}A\pi ^{}\pi ^{}A^{}`$ production reaction was studied on several nuclei $`H`$, $`{}_{}{}^{12}C`$, $`{}_{}{}^{27}Al`$ and $`{}_{}{}^{64}Cu`$ at incident pion momenta (energies) of $`p_\pi ^{}`$=408 and 750 MeV/c ($`T_\pi ^{}`$=291.6 and 623.3 MeV). The figures presented in Ref., however, are only for the results obtained at the higher momentum. In the near-threshold region the $`\pi ^+\pi ^{}`$ and $`\pi ^{}\pi ^{}`$ channels can be directly compared since they share similar dynamical aspects. The $`\pi ^{}\pi ^{}`$ system cannot be in a $`P`$state, and $`\pi ^+\pi ^{}`$ does not couple to $`P`$waves (see Paragraph 5.1). Also, the scattering amplitudes of the two reaction channels, $`\pi ^{}\pi ^{}\pi ^{}`$ and $`\pi ^+\pi ^+\pi ^{}`$, have the two leading isospin amplitudes $`T_{3,2}`$ and $`T_{1,0}`$ in common, where the first (second) index is the total isospin (dipion isospin) . The $`M_{\pi ^{}\pi ^{}}^A/M_{\pi ^{}\pi ^{}}^C`$ ratio in Ref. has marked similarities to the observable $`𝒞`$$`{}_{}{}^{A}{}_{\pi ^+\pi ^{}}{}^{}`$. Namely, the two ratios are flat at around 400 MeV invariant mass but increase as $`M_{\pi \pi }`$ approaches the $`2m_\pi `$ threshold, and, near threshold, both ratios increase as $`A`$ increases. Unlike the CHAOS data, the invariant mass distributions in show no evidence of a peak at the $`2m_\pi `$ threshold. However, the CB and CHAOS spectra do share relevant features: they both have meager $`M_{\pi \pi }`$ strength near threshold for $`H`$ ($`{}_{}{}^{2}H`$ for CHAOS), and show a dramatic increase of strength with increasing $`A`$. As outlined earlier in the present work, the CHAOS peak may be attributed to the finite out-of-plane acceptance of the spectrometer, but the pion distributions are relatively unaffected by nuclear distortions because of the low kinetic energies of the two final pions . The mean kinetic energy of the CB $`\pi ^{}`$’s, however, may exceed 200 MeV, and at this energy pions may undergo large distortions while leaving the nucleus . Pions may be absorbed and the absorption rate depends on both the pion kinetic energy and its initial position inside a nucleus. In addition, pions scatter via the $`\pi A\pi ^{^{}}A`$ reaction and change direction. This smears out the $`\pi ^{}\pi ^{}`$ opening angle distributions, and ultimately the $`M_{\pi ^{}\pi ^{}}`$ ones (see discussion in Paragraph 5.2). Such distortions are difficult to model since they depend on both the pion energy and the position inside the nucleus. Thus, in order to study $`\pi \pi `$ medium modifications, it is advisable to deal with low-energy final pions .
7. Conclusions
In this article the results of an exclusive measurement of the pion-production $`\pi ^+A\pi ^+\pi ^\pm A^{}`$ reactions on $`{}_{}{}^{2}H`$, $`{}_{}{}^{12}C`$, $`{}_{}{}^{40}Ca`$, and $`{}_{}{}^{208}Pb`$, at an incident pion energy $`T_{\pi ^+}`$=283 MeV, were presented. The primary interest was directed to the study of the $`\pi \pi `$dynamics in nuclear matter. The reaction was initially examined on deuterium in order to understand the elementary pion-production mechanism, then on nuclei in order to determine the effects of medium modifications on the $`\pi \pi `$system. Some of latter effects could be obtained by direct comparison of the $`\pi 2\pi `$ distributions, since the data were taken under the same kinematical conditions. The $`\pi \pi `$ properties were highlighted by means of $`M_{\pi \pi }`$, $`cos\mathrm{\Theta }_{\pi \pi }^{CM}`$, $`\mathrm{\Theta }_{\pi \pi }`$, etc. differential cross sections as well as the composite $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ ratios.
$`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ was found to yield the net effect of nuclear matter on the $`\pi \pi `$ system regardless of the $`\pi 2\pi `$ reaction mechanism used to produce the pion pair. These distributions display a marked dependence on the charge state of the final pions. (i) The $`𝒞`$$`{}_{}{}^{A}{}_{\pi ^+\pi ^{}}{}^{}`$ distributions peak at the $`2m_\pi `$ threshold and their yield increases as $`A`$ increases, thus indicating that pion pairs form a strongly interacting system. Furthermore, the $`\pi \pi `$ system couples to $`I=J=0`$, the $`\sigma `$meson channel. (ii) In the $`\pi ^+\pi ^+\pi ^+`$ channel, the $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ distributions show little dependence on either $`A`$ or $`T`$, thus indicating that nuclear matter only weakly affects the $`(\pi \pi )_{I,J=2,0}`$ interaction.
The $`𝒞`$$`{}_{}{}^{A}{}_{\pi ^+\pi ^{}}{}^{}`$ observable was compared with theories which included the $`(\pi \pi )_{I=J=0}`$ in-medium modifications associated with the partial restoration of chiral symmetry in nuclear matter, and with model calculations which only included standard many-body correlations, i.e. the $`P`$wave coupling of $`\pi ^{}s`$ to $`ph`$ and $`\mathrm{\Delta }h`$ configurations. It was found that both mechanisms are necessary to interpret the data, although chiral symmetry restoration yields the larger near-threshold contribution. If this conclusion is correct, the $`\pi 2\pi `$ CHAOS data would provide an example of a distinct $`QCD`$ effect in low-energy nuclear physics.
Monte Carlo simulations of the $`\pi ^+A\pi ^+\pi ^\pm N[A1]`$ reaction phase space proved useful in the interpretation of some of the $`\pi 2\pi `$ data. In the case of $`M_{\pi ^+\pi ^+}`$, the distribution of $`\pi ^+\pi ^+`$ pairs follows phase space. In addition, simulations are able to describe the high-energy part of the distributions, which are sensitive to the nuclear Fermi momentum of the interacting $`\pi ^+p[A1]\pi ^+\pi ^+n[A1]^{}`$ proton. The same conclusions apply to the $`T_{\pi ^+}`$ distributions, thus giving confidence in the reliability of the data. For the $`M_{\pi ^+\pi ^{}}`$ distributions, the $`\pi \pi `$ dynamics overwhelms the dipion kinematics. Unlike the phase space simulations, the near-threshold $`\pi ^+\pi ^{}`$ yield is suppressed in the elementary production reaction, $`\pi ^+`$ $`{}_{}{}^{2}H\pi ^+\pi ^{}pp`$ in the present work, while, in the same energy range, medium modifications strongly enhance $`M_{\pi ^+\pi ^{}}`$. Any interpretation of the shape of the $`M_{\pi ^+\pi ^\pm }^A`$ (or $`𝒞`$$`{}_{}{}^{A}{}_{\pi ^+\pi ^\pm }{}^{}`$) distributions should combine the effects of the restoration of chiral symmetry in nuclear matter, and standard many-body correlations. Such an approach would not require a high-density nuclear medium to obtain strong $`\pi \pi `$ medium modification in the proximity of the $`2m_\pi `$ threshold.
The results of the present measurement were compared with results from the Crystal Ball Collaboration at the BNL, which are the only relevant ones available. The two results show remarkable similarities. The difference in shape of the near-threshold $`M_{\pi \pi }^A`$ distributions can be ascribed to the limited CHAOS acceptance and to nuclear pion distortions in the case of the CB data. Most important, however, is that the exclusive $`\pi 2\pi `$ CB measurement independently confirms that the nuclear medium strongly influences the $`\pi \pi `$ interaction in the $`I=J=0`$ channel.
The experiment was performed at the TRIUMF Meson Facility by using the CHAOS spectrometer. The experimental method used in a previous measurement at TRIUMF and the results obtained were similar to those discussed in the present article. However, the first-generation equipment employed and the $`\pi ^+A\pi ^+\pi ^{}A^{}`$ model used to interpret the results were able to disclose only part of the underlying physics.
Acknowledgements
The authors would like to acknowledge the support received from TRIUMF. The present work was made possible by grants from the Istituto Nazionale di Fisica Nucleare (INFN) of Italy, the National Science and Engineering Research Council (NSERC) of Canada, the Australian Research Council. The authors would also like to acknowledge useful discussions with Z. Aouissat, G. Chanfray, T. Hatsuda, E. Oset, R. Rapp, P. Schuck, M. Vicente-Vacas and J. Wambach.
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# 1 Introduction
## 1 Introduction
The self-interactions of the gauge bosons are consequences of the non-Abelian structure of the electroweak sector of the Standard Model; therefore, the strength of trilinear and quartic gauge couplings is predicted as a result of the gauge symmetry of the theory. The study of trilinear gauge boson couplings in two-boson production processes is within the reach of existing accelerators and measurements of WWZ and WW$`\gamma `$ couplings are being performed with increasing precision in $`\mathrm{e}^+`$$`\mathrm{e}^{}`$ and $`\mathrm{p}\overline{\mathrm{p}}`$ collisions. While non-zero values are predicted for the couplings in the triple and quartic vertices involving charged gauge bosons, the tree level vertices Z$`\gamma `$Z, Z$`\gamma \gamma `$ and ZZZ are not generated by the Standard Model Lagrangian; higher order corrections through virtual loops contribute at the level of 10<sup>-4</sup> , well below the current experimental sensitivity. Nevertheless, new phenomena with a characteristic mass scale above the present experimental threshold might lead to tree-level neutral trilinear gauge couplings (NTGC) in the effective Lagrangian parametrising the residual low energy effects from new physics. For example, as suggested in , virtual effects from new heavy fermions having non-standard couplings to the gauge bosons might generate sizeable anomalous contributions.
The most general Z$`\gamma `$V vertex (where V is the intermediate virtual boson, either photon or a Z) compatible with Lorentz invariance and electro-magnetic gauge invariance involves four independent operators, corresponding to the allowed helicity states for the Z$`\gamma `$ pair . Therefore, in a model independent description, there exist eight couplings: four of them ($`\mathrm{h}_i^\mathrm{Z}`$, $`i=1,\mathrm{},4`$) corresponding to V=Z and four ($`\mathrm{h}_i^\gamma `$ ) corresponding to V=$`\gamma `$. The vertex function was first given in . In this analysis, we adopt the most recent convention established in . The lowest dimensional operators associated to the $`\mathrm{h}_i^{\mathrm{Z},\gamma }`$ ($`i=1,3`$) couplings are of dimension six, while dimension eight operators are associated to $`\mathrm{h}_i^{\mathrm{Z},\gamma }`$ ($`i=2,4`$). As discussed in , no additional symmetry constraints, such as the SU(2)$`\times `$U(1) gauge invariance usually assumed in the case of WWV anomalous couplings , can help in reducing the number of free parameters in the neutral gauge vertex, since operators of even higher dimensionality would be required. Therefore a model independent approach is used which retains all eight couplings.
In this paper, the process $`\mathrm{e}^+`$$`\mathrm{e}^{}`$$`\mathrm{Z}\gamma `$ at $`\sqrt{s}=189`$ GeV is investigated through the final states $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$and $`\nu \overline{\nu }\gamma `$, the dominant decay modes of the $`\mathrm{Z}`$ boson, with the aim of searching for Z$`\gamma `$Z and Z$`\gamma \gamma `$ couplings. Experimental constraints on these couplings have been produced in the past from the analysis of LEP data at lower energies and of the Tevatron data . The analysis presented here has higher sensitivity due to the increased centre-of-mass energy and due to the large data sample collected during the 1998 operation of LEP. It also benefits from the recent clarification of the theoretical framework in which neutral gauge boson self-interactions can be described. A recent analysis based on data collected at $`\sqrt{\mathrm{s}}=189\mathrm{GeV}`$ by the L3 collaboration adopts the same convention as in this paper. On the other hand, due to this different convention, the comparison of the results presented here and in with previous published results is not straightforward.
In $`\mathrm{e}^+`$$`\mathrm{e}^{}`$ collisions, the production of Z$`\gamma `$ final states via anomalous neutral gauge couplings has a large irreducible Standard Model background from $`\mathrm{Z}^0`$ production with hard initial state radiation (ISR). Small contributions to $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$production also arise from $`\mathrm{e}^+`$$`\mathrm{e}^{}`$$`\gamma \gamma ^{}\gamma \mathrm{q}\overline{\mathrm{q}}`$ and from final state radiation in $`\mathrm{q}`$$`\overline{\mathrm{q}}`$ production. In the $`\nu `$$`\overline{\nu }`$$`\gamma `$channel, final states arising from the exchange of a W boson in the t-channel also contribute to the cross-section, although their contribution is small within the Z$`\gamma `$ signal acceptance used in this analysis. As a general property, the Z and $`\gamma `$ produced at the anomalous vertices are more isotropically distributed than the dominant Standard Model background, which is characterized by the strongly forward peaked angular distribution of initial state radiation. Therefore, deviations from the Standard Model predictions due to Z$`\gamma `$V couplings would be more pronounced for large angles between the beam direction and the photon. In the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$final state, the angular distribution of the jets can also be exploited in order to gain sensitivity to anomalous couplings, due to the resulting enhancement of the longitudinal polarisation of the Z boson affecting the fermion decay angle. On the other hand, the photon energy spectrum has a marginal sensitivity, due to the kinematic constraints from the fixed centre-of-mass energy and the narrow Z resonance. Since all the terms in the Z$`\gamma `$V vertex are proportional to the momenta of the gauge bosons involved, this results in an enhancement of the sensitivity to NTGC as the centre-of-mass energy increases. Finally, the experimental signature of the anomalous $`\mathrm{h}_i^{\mathrm{Z},\gamma }`$ couplings depends on the CP parity of the associated operators. The vertex terms proportional to $`\mathrm{h}_1^{\mathrm{Z},\gamma }`$ and $`\mathrm{h}_2^{\mathrm{Z},\gamma }`$ violate CP and, hence, do not interfere with the CP conserving Standard Model amplitudes. As a result the total and differential cross-sections receive only additive contributions from the anomalous processes. The remaining couplings, associated to CP even terms, lead to amplitudes which interfere with the Standard Model; therefore the differential and total cross-sections are enhanced or suppressed depending on the sign and the size of the $`\mathrm{h}_i^{\mathrm{Z},\gamma }`$ ($`i=3,4`$) couplings.
## 2 Detector and Monte Carlo Simulation
The OPAL detector, described in detail in , consists of a central tracking system inside a solenoid providing a magnetic field of 0.435 T, and of an electromagnetic calorimeter, complemented by a presampling system and an array of scintillation counters for time-of-flight measurements; hadron calorimetry is obtained by instrumenting the magnet return yoke which is surrounded by muon chambers. A system of forward calorimeters extends the angular coverage of the detector down to a polar angle<sup>1</sup><sup>1</sup>1In the OPAL coordinate system, $`\theta `$ is the polar angle defined with respect to the electron beam direction and $`\varphi `$ is the azimuthal angle. of 24 mrad. However, due to the installation in 1996 of a thick tungsten shield designed to protect the tracking chambers from synchrotron radiation background, the effective limit of electromagnetic hermeticity is around 33 mrad. The integrated luminosity of the data samples is determined from the rate of small angle Bhabha scattering events observed in the silicon-tungsten calorimeter with a precision of 0.22%.
Track reconstruction is performed by combining the information from a silicon microvertex detector, a vertex drift chamber, a large volume jet drift chamber and an outer layer of drift chambers for the measurement of the $`z`$ coordinate. The most relevant subdetector for the event topologies used in the analysis presented here is the electromagnetic calorimeter. It consists of an array of 9440 lead-glass blocks in the barrel ($`|\mathrm{cos}\theta |<0.82`$) arranged in an almost-pointing geometry and two dome-shaped end caps, each of 1132 longitudinally aligned lead-glass blocks, covering the polar angle range $`0.81<|\mathrm{cos}\theta |<0.984`$. Trigger signals , based on energy deposits in the lead-glass blocks and also on a coincidence of energy in the barrel electromagnetic calorimeter and a hit in the time-of-flight system, guarantee full trigger efficiency for both the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$ and the $`\nu `$$`\overline{\nu }`$$`\gamma `$events falling within the signal definition criteria used in this analysis.
The Standard Model processes leading to the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$and $`\nu `$$`\overline{\nu }`$$`\gamma `$ final states have been simulated using the Monte Carlo generators KK2f for the hadronic channel and Koralz and Nunugpv98 for the missing energy channel. For the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$channel, the fragmentation, which includes photon radiation from the quarks, and the hadronization are simulated with the Jetset package tuned on the basis of extensive studies of hadronic events at the Z resonance as described in . The grc4f generator has been used to estimate the background to the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$channel from four-fermion production. The contribution to the background from two-photon interactions has been studied with the Monte Carlo generators Phojet , for the untagged and double-tagged events, and Herwig , for the single-tagged events and charged current deep inelastic scattering events. In the $`\nu `$$`\overline{\nu }`$$`\gamma `$channel, the contamination from Bhabha events has been estimated using the Bhwide and Teegg generators. The contamination from four-fermion production has been studied using grc4f and Koralw , while to determine the background from two-photon interactions the Vermaseren Monte Carlo generator has been used. The Radcor Monte Carlo has been used to study the $`\mathrm{e}^+\mathrm{e}^{}\gamma \gamma `$ background and the energy response of the calorimeter to photons. All the Monte Carlo samples described above were processed through the OPAL detector simulation . For the interpretation of the data, as will be described in more detail in section 4, a Monte Carlo generator , based on the matrix element for $`\mathrm{f}\overline{\mathrm{f}}\gamma `$ production in $`\mathrm{e}^+`$$`\mathrm{e}^{}`$ collisions and including the contributions from NTGC, has been used.
## 3 Event Selection and Cross-Section Measurements
### 3.1 The selection of $`𝐪`$$`\overline{𝐪}`$$`𝜸`$ events
The selection of hadronic events with isolated high energy photons is performed on events preselected as high multiplicity hadronic events in a data sample corresponding to an integrated luminosity of 176.2 $`\mathrm{pb}^1`$ with an average centre-of-mass energy of 188.6 $`\mathrm{GeV}`$. The preselection criteria, described in , are based on the track and cluster multiplicity, on the visible energy and on the longitudinal imbalance of the energy measured in the electromagnetic calorimeter. The events satisfying the preselection requirements are processed by a photon search algorithm.
The photon identification is based on an algorithm optimised for photon search in hadronic events described in . Electromagnetic clusters without associated tracks in the central detector are accepted as photon candidates if their energy is higher than 5% of the beam energy and their polar angle lies in the acceptance region of the lead-glass calorimeter. The number of lead-glass blocks involved and the energy sharing among them are required to correspond to typical patterns defined for photon identification in the OPAL calorimeter. An isolation criterion is then applied in order to reject electromagnetic clusters associated with jets. The total energy deposition in the electromagnetic calorimeter (not associated to the photon candidate) within a cone of $`15^{}`$ around the photon flight direction is required to be less than 2 $`\mathrm{GeV}`$. In addition, the sum of the momenta of tracks which, extrapolated to the calorimeter surface, fall inside a $`15^{}`$ cone around the photon impact point is also required to be lower than 2 $`\mathrm{GeV}`$. Finally, systems of one or two well reconstructed tracks associated with electromagnetic clusters, consistent with a photon conversion according to the criteria described in , are included if the reconstructed photon satisfies the criteria listed above. The photon identification algorithm has been extensively studied in order to assess the level of accuracy of the estimate of the efficiency obtained from Monte Carlo $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$ events. In particular, the selection criteria of the algorithm have been adjusted in order to minimise their sensitivity to unsatisfactory modelling in the simulation. A residual discrepancy in the identification efficiency for converted photons has been observed and taken into account as a correction of 1.25% to the overall selection efficiency. As an example of the quality of the modelling of the photon identification algorithm, figure 1 shows the angle between the photon candidate and the closest track in the event and the total charged energy in the isolation cone.
After the photon search, all the clusters and tracks in the event which are not associated to the most energetic photon candidate are grouped into jets according to the Durham $`k_T`$ scheme with resolution parameter $`y=0.02`$. If more than four jets are reconstructed, the event is forced to have four jets in addition to the isolated high energy photon.
In $`\mathrm{e}^+`$$`\mathrm{e}^{}`$$`\mathrm{Z}\gamma `$ at $`\sqrt{s}=189`$ GeV the photon energy spectrum is peaked at approximately 72 GeV, reflecting the sharp Z resonance. In order to select the topology corresponding to a high energy photon recoiling against a hadronic system of invariant mass equal to the Z boson mass, the signal definition is based on kinematic cuts applied to the most energetic photon in the event:
* 50 $`\mathrm{GeV}`$ $`<\mathrm{E}_\gamma <`$ 90 $`\mathrm{GeV}`$ ;
* 15 $`<\theta _\gamma <`$ 165,
where $`\mathrm{E}_\gamma `$ and $`\theta _\gamma `$ are the photon energy and polar angle, respectively.
To improve the photon energy resolution and suppress further the background from non-$`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$events and the feedthrough from events outside the signal definition, a kinematic fit is applied to all the events with at least one photon of energy larger than 30 GeV. The fit imposes energy and momentum conservation using as input the photon and the jet momenta. Undetected ISR is allowed to compensate for missing longitudinal momentum in the beam pipe region if the $`\chi ^2`$ probability of the fit is smaller than 1%.
From a study of Monte Carlo $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$events, the kinematic fit improves the photon energy resolution by a factor two. The events for which the fit converges (99.5%) are finally selected if the fitted values of $`E_\gamma `$ and $`\theta _\gamma `$ satisfy the signal definition cuts. In order to suppress contamination from photons originating from the jets, only the sub-sample in which $`\alpha _{\gamma \mathrm{jet}}>30^{}`$, where $`\alpha _{\gamma \mathrm{jet}}`$ is the angle between the photon and the closest jet, is retained for the analysis.
The selection efficiency and the feedthrough in the kinematic acceptance are listed in table 1. They have been estimated from fully simulated $`\mathrm{e}^+`$$`\mathrm{e}^{}`$$``$$`\mathrm{q}`$$`\overline{\mathrm{q}}`$ Monte Carlo events, where the signal-like topology arises from radiative return to the Z resonance. The efficiency has been calculated with respect to the kinematic signal acceptance defined above. The feedthrough represents the fraction of selected events which do not belong to the kinematic acceptance; it is mainly due to resolution effects, but it includes also a small contamination (0.35%) due to non-ISR photon candidates in $`\mathrm{e}^+`$$`\mathrm{e}^{}`$$``$$`\mathrm{q}`$$`\overline{\mathrm{q}}`$ events falling in the signal acceptance after reconstruction. Both the efficiency and the feedthrough are corrected for the aforementioned residual disagreement between the performance of the photon search algorithm in the data and in the Monte Carlo.
The numbers of events selected in the data and in the Monte Carlo samples are listed in table 2. Figure 2 shows the distributions of the photon energy and polar angle and angular separation with respect to the closest jet for the events selected in the data, compared with the Standard Model expectation from Monte Carlo. The total background (2.5% of the selected events) comes from four-fermion production (1.59%) two-photon interactions (0.59%) and tau-pair production (0.34%). The agreement between the data and the Monte Carlo is in general satisfactory, except perhaps in the photon energy distribution where a slight deficit of events is observed in the radiative return peak.
### 3.2 The selection of $`𝝂`$$`\overline{𝝂}`$$`𝜸`$ events
The selection of events with an isolated high energy photon accompanied by missing energy and low activity in the detector follows the single-photon analysis described in . The data sample used in the analysis corresponds to an integrated luminosity of 177.3 $`\mathrm{pb}^1`$, with an average centre-of-mass energy of 188.6 $`\mathrm{GeV}`$. After the single photon selection, 643 events are retained in the data. The same additional conditions which define the kinematic acceptance in the $`\mathrm{q}\overline{\mathrm{q}}\gamma `$ channel, $`50\mathrm{GeV}<\mathrm{E}_\gamma <90\mathrm{GeV}`$ and $`15^{}<\theta _\gamma <165^{}`$, are then applied to the most energetic photon in the electromagnetic calorimeter.
The predicted efficiency and the feedthrough in the kinematic acceptance are listed in table 1. They have been estimated averaging the predictions of the Koralz and the Nunugpv98 Monte Carlo, which in the signal acceptance agree within (0.2$`\pm `$0.5)% in the efficiency and (1$`\pm `$9)% in the feedthrough.
The number of events selected in the data are listed in table 3, together with those expected from all the relevant physics processes and from instrumental backgrounds. As potential sources of physics background, four-fermion processes, radiative Bhabha and $`\mathrm{e}^+\mathrm{e}^{}\gamma \gamma `$ events have been considered, while the residual cosmic ray and beam-related contamination have been estimated using control samples enriched in these backgrounds. Both these sources result in an overall negligible (0.24%) contribution to the selected events. Figure 3 shows the distributions of the energy and of the polar angle of the photon for the selected events, compared with the Standard Model expectation.
### 3.3 Cross-section Measurements
From the number of observed events in each channel and from the predicted efficiency, feedthrough and backgrounds, as presented in sections 3.1 and 3.2, the cross-sections within the kinematic signal acceptance are measured to be:
$`\sigma _{\mathrm{q}\overline{\mathrm{q}}\gamma }`$ $`=`$ $`9.42\pm 0.25(\mathrm{stat}.)\pm 0.15(\mathrm{syst}.)\mathrm{pb}`$
$`\sigma _{\nu \overline{\nu }\gamma }`$ $`=`$ $`2.52\pm 0.13(\mathrm{stat}.)\pm 0.05(\mathrm{syst}.)\mathrm{pb}.`$
These measurements are in reasonable agreement, within the errors which are dominated by the statistical uncertainty, with the predictions from the Standard Model, which are respectively $`\sigma _{\mathrm{q}\overline{\mathrm{q}}\gamma }^{\mathrm{SM}}=9.75\pm 0.03`$ pb and $`\sigma _{\nu \overline{\nu }\gamma }^{\mathrm{SM}}=2.81\pm 0.02`$ pb. The predicted $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$ cross-section is obtained from KK2f, while the $`\nu `$$`\overline{\nu }`$$`\gamma `$ cross-section is the average of the predictions of Koralz and Nunugpv98 which are consistent within $`(1.6\pm 1.1)\%`$. The errors associated with the Standard Model predictions come from the Monte Carlo statistics. Taking into account the experimental systematic uncertainties discussed in the following section, the overall discrepancy corresponds to $`1.1`$ standard deviations for the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$channel and $`2.1`$ for the $`\nu `$$`\overline{\nu }`$$`\gamma `$channel.
### 3.4 Systematic Errors
The different sources of systematic uncertainties affecting the cross-section measurements are summarised in table 4 and are discussed in the following:
Systematics specific to the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$channel:
* The uncertainty on the selection efficiency. Contributions to this uncertainty come from imperfect modelling of the material in the detector, affecting the estimate of the photon conversion rate (0.77%), from inaccuracies in the modelling of the identification efficiency for converted photons (0.62%), from uncertainty in the simulation of the track-cluster association at forward angles (0.64%) and from the photon isolation criterion (0.26%). The total uncertainty coming from the modelling of the selection efficiency translates into a relative error on the cross-section of 1.2%. The uncertainties arising from limited Monte Carlo statistics used to evaluate the efficiency and the feedthrough are respectively 0.14% and 0.07%.
* The modelling of the jet reconstruction. An additional smearing of the jet energies and directions and a shift of the jet energy scale has been applied in the Monte Carlo, on the basis of an extensive comparison of two-jets events in calibration data collected at the $`\mathrm{Z}^0`$ peak and in the simulation. The resulting variation (0.36%) in the cross-section has been assigned as a systematic error.
* The sensitivity of the analysis to the jet multiplicity in the event. The results of a modified analysis, where each event has been forced to contain exactly two jets in addition to the isolated photon, have been compared with those from the standard analysis. The difference has been assigned as a systematic error (0.69%).
* The uncertainty related to the models used to simulate the hadronisation process. This uncertainty (0.08%) has been evaluated comparing the results when either the Jetset or the Herwig hadronisation schemes have been used in the Monte Carlo.
* The uncertainty (0.61%) due to the background subtraction; this is dominated by the 100% uncertainty on the normalisation of the background from two-photon interactions as predicted by Herwig, which is expected to give the best description of two-photon interactions in the data, and by Pythia. This large uncertainty is assigned to cover possible mismodelling of the hard-fragmentation processes in the very small fraction of the two-photon cross-section retained in the selection, as suggested by a comparison with the F2GEN generator. Finally, the uncertainty related to the contamination from fake ISR photons in non-radiative $`\mathrm{q}`$$`\overline{\mathrm{q}}`$ events is estimated to be 0.35%.
Systematics specific to the $`\nu `$$`\overline{\nu }`$$`\gamma `$channel:
* The uncertainty on the selection efficiency. Systematic errors come from the estimate of the detector occupancy (1%) and from the imperfect description of the conversion probability and of the conversion tagging efficiency in the Monte Carlo (0.8%). A 0.9% uncertainty is assigned to account for systematic effects in the simulation of the cuts used to reject radiative Bhabha and $`\mathrm{e}^+\mathrm{e}^{}\gamma \gamma `$ events. The uncertainty on the efficiency of the timing cuts used to reject background from cosmic rays translates into a systematic error of 0.5%, while an uncertainty of 0.6% is assigned to cover possible mismodelling of the other cuts complementing the timing requirements in the cosmic and instrumental background rejection. The overall uncertainty on the modelling of the selection efficiency is estimated to be 1.7%. Finally, the statistical uncertainty on the efficiency and feedthrough as estimated from Monte Carlo contributes with a 0.3% uncertainty on the cross-section.
* The uncertainty on the background from processes other than $`\nu `$$`\overline{\nu }`$$`\gamma `$, amounting to 0.25%. This includes a 50% systematic uncertainty on the expected background, which was assigned to cover possible mismodelling of the vetoes used in the rejection.
Systematic uncertainties common to both analyses:
* The uncertainty on the angular acceptance arising from residual biasses in the coordinate reconstruction and from the absolute knowledge of the detector geometry. The overall uncertainty on the position of the electromagnetic showers, estimated in to be 0.001 rad, results in a relative uncertainty of 0.19%.
* The effects of mismodelling of the energy scale and the resolution of the electromagnetic calorimeter. They have been investigated using a control sample of about one thousand strongly collinear $`\mathrm{e}^+\mathrm{e}^{}\gamma \gamma `$ events, whose energies are expected to be very close to the beam energy. The Monte Carlo has been found to reproduce the electromagnetic calorimeter energy scale and resolution in the data within respectively 0.3% and 10%. These differences are in very good agreement with those observed using as a reference process a sample of wide angle Bhabha events . The systematic error has then been estimated by modifying the absolute energy scale and resolution in the Monte Carlo within these limits and the corresponding variation of the cross-section has been taken as systematic error.
* The uncertainty on the beam energy ($`\pm 20`$ MeV). The effect on the cross-section was evaluated by appropriately scaling the energy of the most energetic photon in the Monte Carlo, while leaving unchanged the invariant mass of the system recoiling against the photon.
* The uncertainty on the measurement of the integrated luminosity, 0.22%.
## 4 Data Interpretation
Since the data presented in the previous section show no evidence of deviations with respect to the Standard Model expectation, they are used to derive bounds on the strength of anomalous $`\mathrm{Z}\gamma \mathrm{Z}`$ and $`Z\gamma \gamma `$ couplings.
### 4.1 Analysis Procedure
The analysis is based on a comparison of the measured event rate and of the differential distributions with the theoretical predictions for the Standard Model processes and possible contributions from NTGC. The energy spectrum and the cosine of the polar angle of the photon are used in both channels, while in the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$channel the cosine of the angle between the photon and the closest jet is also used. Due to the almost monochromatic photon energy spectrum, this angle is strongly correlated to the quark emission angle in the Z decay rest frame, which is sensitive to NTGC.
The theoretical predictions as a function of different values of the anomalous couplings $`\mathrm{h}_i^{\mathrm{Z},\gamma }`$ are obtained from a Monte Carlo generator for $`\mathrm{f}\overline{\mathrm{f}}\gamma `$ production in $`\mathrm{e}^+`$$`\mathrm{e}^{}`$ collisions. This generator is based on the full matrix element, in the lowest order approximation, for all the relevant Standard Model processes and for processes generated by anomalous trilinear neutral gauge couplings. The only Standard Model contribution missing in the calculation is the t-channel W boson exchange in the $`\nu `$$`\overline{\nu }`$$`\gamma `$channel, which has been estimated from Nunugpv98 and included in the reweighting procedure described in the following. The effect of higher order QED corrections from initial state radiation, which have been found to reduce the contribution from anomalous couplings by typically 15%, has been incorporated into the calculation using a collinear radiator function from the Excalibur Monte Carlo.
The events selected in the data have been classified in $`5\times 4\times 4`$ unequal bins of the three-dimensional $`(\mathrm{E}_\gamma ,\mathrm{cos}\theta _\gamma ,\mathrm{cos}\alpha _{\gamma \mathrm{jet}})`$ space. In the case of the $`\nu `$$`\overline{\nu }`$$`\gamma `$channel, the angle between the photon and the $`\mathrm{Z}`$ decay products is not experimentally accessible and therefore integrated out. The NTGC-dependent theoretical prediction for the population in each cell, provided by the Monte Carlo calculation, has been modified to allow for the reconstruction efficiency and resolution effects, as determined from a large sample of fully simulated Z$`{}_{}{}^{0}/\gamma \mathrm{q}\overline{\mathrm{q}}`$ and Z$`{}_{}{}^{0}\nu \overline{\nu }`$ Standard Model Monte Carlo events. The total number of expected events and the population of each cell as a function of the anomalous couplings, $`\mathrm{N}(\mathrm{h}_i)`$, is determined by reweighting the number $`\mathrm{N}_{\mathrm{SM}}`$ of accepted events predicted by the fully simulated Standard Model Monte Carlo according to:
$$\mathrm{N}(\mathrm{h}_i)=\mathrm{N}_{\mathrm{SM}}(1+\delta (\mathrm{h}_i)),\delta (\mathrm{h}_i)=\frac{\mathrm{N}_{\mathrm{rec}}(\mathrm{h}_i)\mathrm{N}_{\mathrm{rec}}(\mathrm{h}_i=0)}{\mathrm{N}_{\mathrm{rec}}(\mathrm{h}_i=0)},$$
(1)
where $`\mathrm{N}_{\mathrm{rec}}`$ is the number of reconstructed events from the NTGC-dependent theoretical prediction, modified for efficiency and resolution as explained above. Figures 4 and 5 show how the distributions of the kinematic variables, folded with detector effects, are modified by a particular choice of anomalous couplings in the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$and $`\nu `$$`\overline{\nu }`$$`\gamma `$channels. In order to disentangle the effect of NTGC on the event rate and on the differential distributions the latter have been normalised to the number of events selected in the data.
The theoretical expectations for the event rate and the differential distributions have been fitted to the data independently for the two channels, under the hypothesis that only one coupling at a time is non-zero. The most probable values of the anomalous couplings are determined by minimising the negative log-likelihood defined as follows:
$$\mathrm{LogL}=\mathrm{LogP}(\mathrm{N}^{\mathrm{obs}},\mathrm{N}(\mathrm{h}))\underset{j}{}\mathrm{LogP}(\mathrm{N}_j^{\mathrm{obs}},\mathrm{N}_j(\mathrm{h}))$$
where $`\mathrm{P}(\mathrm{N}^{\mathrm{obs}},\mathrm{N}(\mathrm{h}))`$ is the Poisson probability of observing the number of events $`\mathrm{N}^{\mathrm{obs}}`$ if the expectation is $`\mathrm{N}(\mathrm{h})`$. The index $`j`$ runs over the number of cells defining the multidimensional distributions and the condition $`_j\mathrm{N}_j(\mathrm{h})=\mathrm{N}^{\mathrm{obs}}`$ is imposed to disentangle the contributions to the likelihood from event rate and distributions.
The fit procedure has been tested on Monte Carlo samples of Standard Model $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$and $`\nu `$$`\overline{\nu }`$$`\gamma `$ events reweighted for NTGC effects. The central values of the fit results have been found to correctly reproduce the values of the input NTGC parameters. In order to check the reliability of the errors on the fit results, tests have been performed on several Standard Model Monte Carlo samples of size corresponding to the data luminosity. The distributions of the central values of the couplings determined by the fit are found to be consistent with those expected from the statistical sensitivity.
### 4.2 Results on Trilinear Neutral Gauge Couplings
The values of the anomalous couplings and the statistical errors obtained from the likelihood fit of the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$and $`\nu `$$`\overline{\nu }`$$`\gamma `$ event rate and distributions are listed in table 5, together with the expected statistical sensitivity of the analysis which would be achieved in the case of perfect agreement between data and the Standard Model predictions. In general, the total event rate and the differential distributions have similar sensitivities to NTGC. However, due to the quadratic dependence of the cross-section on the couplings, the fit of the event rate can only determine the value of the couplings with a two-fold ambiguity. In the case of the CP violating couplings, $`\mathrm{h}_i^{\mathrm{Z},\gamma }`$ ($`i=1,2`$), which lead to amplitudes that do not interfere with the Standard Model amplitudes, the sign of the couplings is completely undefined, but both the cross-section and the distributions provide a determination of the absolute value of the couplings. In the case of the CP conserving couplings, $`\mathrm{h}_i^{\mathrm{Z},\gamma }`$ ($`i=3,4`$), the interference with the Standard Model amplitudes results in distributions of the kinematic variables which produce a unique minimum in the associated $`\mathrm{log}\mathrm{L}`$ function, thus removing the two-fold ambiguity arising from the event rate information.
The 95% Confidence Level (C.L.) bounds on the eight anomalous couplings have been obtained by convolving the likelihood function with a Gaussian whose width $`\sigma `$ corresponds to the systematic uncertainty on the individual parameters, as estimated in section 4.3. The central values and the 95% C.L. intervals resulting from the combination of the two channels are given in table 6. Figures 6 and 7 show the corresponding negative log-likelihood curves for the individual channels and their combination. In combining the results, the correlations between the systematic uncertainties have been taken into account.
### 4.3 Systematic Errors
The impact of several sources of systematic uncertainty on the NTGC determination has been assessed. Most of the effects considered have already been discussed in the context of the systematic error on the cross-section measurements. A few more are related to the modelling of NTGC effects on the observables used in the likelihood fit and to the reference Standard Model predictions. For all the sources of systematic uncertainty, a symmetric error is assigned based on the maximum absolute shift between the central value, the lower and the upper edges of the 68% C.L. interval as obtained in the standard analysis (table 5) and in a specific fit to the data performed to simulate the systematic effect. The dominant sources of systematic uncertainty come from:
* The effects of the modelling of the selection efficiency. These have been evaluated using the same methods as discussed in section 3.4, and have been assigned as scale uncertainties on the expected number of events in each channel. Where relevant, the angular dependence of the uncertainties has been taken into account. For the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$channel, this includes also the uncertainties coming from the jet multiplicity, the hadronisation modelling and the jet parameter smearing.
* The theoretical uncertainty on the Standard Model prediction. Considerations about missing higher order corrections in the KK2f Monte Carlo lead to an estimate of the theoretical uncertainty of the order of 1% in the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$channel . In the $`\nu `$$`\overline{\nu }`$$`\gamma `$channel, a 2% theoretical uncertainty has been assigned; this uncertainty covers the observed differences between the kinematic cross-sections and the selection efficiencies estimated by Koralz and Nunugpv98, and is consistent with a recent comparison between different theoretical calculations presented in . These theoretical uncertainties have been assigned as an overall normalisation error on the Standard Model prediction in each channel.
Other minor contributions to the systematic error have been considered:
* The contribution of the background to the fit result. This has been separated into a component due to a scale factor (corresponding to the Monte Carlo statistical uncertainty on the background absolute rate) and a component due to the modelling of the background shape. The error due to the background normalisation is evaluated as the maximum effect observed when increasing or decreasing the total background rate by one standard deviation. The error due to the shape modelling has been conservatively assessed by assuming a flat distribution of the background events in the signal acceptance region. In the $`\nu `$$`\overline{\nu }`$$`\gamma `$channel, due to the very small contamination, only the uncertainty on the background rate has been considered.
* The uncertainty related to the Monte Carlo statistics. This effect has been estimated by applying the likelihood fit to the data using, as Standard Model prediction in each bin, a number of events generated according to a Poisson distribution with average equal to the prediction of the reference Monte Carlo sample and rescaled to the integrated luminosity of the data. The maximum among the r.m.s of the central values and of the lower and upper 68% C.L. limits on each coupling, over a thousand fits, has been assigned as systematic error.
* The uncertainty on the reweighting procedure coming from the limited Monte Carlo statistics used to evaluate the corrections for resolution and selection efficiency, and from the limited Monte Carlo statistics used in the calculation of the generator-level weights.
* The uncertainty on the missing t-channel W exchange contribution, which applies only to the $`\nu `$$`\overline{\nu }`$$`\gamma `$channel. The procedure used to correct for the missing W t-channel in the reweighting procedure accounts for this process only in the pure Standard Model contribution, but not in the interference between the Standard Model process and the process leading to anomalous coupling. Since the W t-channel is calculated to contribute less than 4% to the Standard Model cross-section in the signal acceptance, the effect of it being neglected in the interference is expected to be of order 2%. The interference term has been conservatively varied by $`\pm 4\%`$ and the differences in the fit results have been assigned as a systematic error.
* The systematic errors arising from the calorimeter energy scale and resolution, the modelling of the $`\theta _\gamma `$ angular cut, the uncertainty on the beam energy and on the integrated luminosity. They have been evaluated with the same method as discussed in section 3.4 and have been treated as fully correlated between the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$channel and the $`\nu `$$`\overline{\nu }`$$`\gamma `$ channel.
The systematic errors on the couplings are summarised in table 7 for both the $`\mathrm{q}`$$`\overline{\mathrm{q}}`$$`\gamma `$channel and the $`\nu `$$`\overline{\nu }`$$`\gamma `$channel. For all the couplings the total systematic error is small compared to the statistical uncertainty, except for $`\mathrm{h}_{3,4}^\gamma `$, where the size of the systematic error reaches 60% of the statistical error.
## 5 Conclusions
Using the data collected at $`\sqrt{s}=189`$ $`\mathrm{GeV}`$, the cross-sections and the differential distributions for hadronic events with a high energy isolated photon and for events with an energetic photon and missing energy have been measured to search for possible contributions from anomalous $`\mathrm{Z}\gamma \mathrm{Z}`$ and $`\mathrm{Z}\gamma \gamma `$ couplings. Since no significant evidence of deviations with respect to the Standard Model is observed, 95% C.L. limits on the eight trilinear neutral gauge couplings $`\mathrm{h}_i^{\mathrm{Z},\gamma }`$ have been derived. These limits do not yet allow to place significant constraints on specific models of new physics leading to effective anomalous couplings in the neutral sector. Nevertheless, these results and those presented in , which are compatible and of equivalent sensitivity, represent the best available investigations in $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}\gamma `$ of the neutral gauge boson self-interactions, in the recently revised theoretical framework describing the general $`\mathrm{Z}\gamma \mathrm{V}(\mathrm{V}=\mathrm{Z},\gamma )`$ vertex.
## Acknowledgements
The authors wish to thank U. Baur, G.J. Gounaris and F.M. Renard for helpful discussions and clarifications. We particularly wish to thank the SL Division for the efficient operation of the LEP accelerator at all energies and for their continuing close cooperation with our experimental group. We thank our colleagues from CEA, DAPNIA/SPP, CE-Saclay for their efforts over the years on the time-of-flight and trigger systems which we continue to use. In addition to the support staff at our own institutions we are pleased to acknowledge the
Department of Energy, USA,
National Science Foundation, USA,
Particle Physics and Astronomy Research Council, UK,
Natural Sciences and Engineering Research Council, Canada,
Israel Science Foundation, administered by the Israel Academy of Science and Humanities,
Minerva Gesellschaft,
Benoziyo Center for High Energy Physics,
Japanese Ministry of Education, Science and Culture (the Monbusho) and a grant under the Monbusho International Science Research Program,
Japanese Society for the Promotion of Science (JSPS),
German Israeli Bi-national Science Foundation (GIF),
Bundesministerium für Bildung und Forschung, Germany,
National Research Council of Canada,
Research Corporation, USA,
Hungarian Foundation for Scientific Research, OTKA T-029328, T023793 and OTKA F-023259.
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# 1 Introduction
## 1 Introduction
2D conformally–invariant field theories have become the subject of intense investigation in recent years, after the work of Belavin, Polyakov, and Zamolodchikov . One of the main reasons for this, is that 2D conformal field theories describe the critical behaviour of two dimensional statistical models \[2–5\]. Conformal field theory provides us with a simple and powerful means of calculating the critical exponents, as well as, the correlation functions of the theory at the critical point . Another application of conformal field theories is in string theories. Originally, string theory was formulated in flat 26–dimensional space–time for bosonic- and flat 10-dimensional space–time for supersymmetric-theories. It has been realized now that the central part of string theory is a 2D conformally–invariant field theory. It is seen that the tree–level string amplitudes may be expressed in terms of the correlation functions of the corresponding conformal field theory on the plane, whereas string loop amplitudes may be expressed in terms of the correlation functions of the same conformal field theory on higher–genus Riemann surfaces \[7–10\].
Supersymmtry is a $`Z_2`$ extension of the Poincaré algebra . But this can be enlarged to a superconformal algebra ( for example). If the dimension of the space–time is two, there are also fractional supersymmetric extensions of the Poincaré and conformal algebra \[14–17\]. Fractional supersymmetry is a $`Z_n`$ extension of the Poincaré algebra. In this paper, the genral form of the two–point functions of a theory with this symmetry is obtained.
According to Gurarie , conformal field theories with logarithmic correlation functions may be consistently defined. In some interesting physical theories like polymers , WZNW models \[20–23\], percolation , the Haldane-Rezayi quantum Hall state , and edge excitation in fractional quantum Hall effect , there appear logarithmic correlation functions. Logarithmic operators are also seen in 2D magnetohydrodynamic turbulence \[27–29\], 2D turbulence and some critical disordered models . Logarithmic conformal field theories for D–dimensional case ($`D>2`$) has also been studied . The general form of the correlators of the 2D conformal field theories and supersymmetric conformal field theories is investigated in and , respectively.
In this paper, the general case $`n=F+1`$ is considered. So the components of the superfield have grades 0, 1,…, and $`F`$. The complex plane is extended by introducing two independent paragrassmann variables $`\theta `$ and $`\overline{\theta }`$, satisfying $`\theta ^{F+1}=\overline{\theta }^{F+1}=0`$. One can develop an algebra, the fractional $`n=F+1`$ algebra, based on these variables and the derivatives with respect to them \[38–40\]. In , the fractional supersymmetry has been investigated by introducing a certain fractional superconformal action. We don’t consider any special action here. What we do, is to use only the structure of the fractional superconformal symmetry to obtain general restrictions on the form of the two–point functions. The scheme of the paper is the following. In section 2, infinitesimal superconformal transformations are defined. In section 3, the generators of these transformation and their algebra are investigated. In section 4, the two–point functions of such theories are obtained. In fact, sections 2–4 are a generalization of . In section 5, logarithmic superconformal field theories are investigated and the chiral- and full- two–point functions of primary and quasiprimary (logarithmic) fields are obtained. This is a generalization of .
## 2 Infinitesimal superconformal transformations
Consider a paragrassmann variable $`\theta `$, satisfying
$$\theta ^{F+1}=0,$$
(1)
where $`F2`$, is a positive integer. A function of a complex variable $`z`$, and this paragrassmann variable, will be of the form
$$f(z,\theta )=f_0(z)+\theta f_1(z)+\theta ^2f_2(z)+\mathrm{}+\theta ^Ff_F(z).$$
(2)
The covariant derivative is defined as
$$D:=_\theta +\frac{q^F\theta ^F}{(F)_{q^1}!}_z,$$
(3)
where
$$(F)_q!:=(F)_q(F1)_q\mathrm{}(1)_q,(m)_q=\frac{1q^m}{1q},$$
(4)
and
$$_\theta \theta =1+q\theta _\theta .$$
(5)
$`q`$ is an $`F+1`$th root of unity, with the property that there exisits no positive integer $`m`$ less than $`F+1`$ so that $`q^m=1`$.
An Infinitesimal transformation
$`z^{}`$ $`=`$ $`z+{\displaystyle \underset{k=0}{\overset{F}{}}}\theta ^k\omega _k(z),`$ (6)
$`\theta ^{}`$ $`=`$ $`\theta +{\displaystyle \underset{k=0}{\overset{F}{}}}\theta ^kϵ_k(z),`$ (7)
is called superconformal if
$$D=(D\theta ^{})D^{},$$
(8)
where
$$D^{}=_\theta ^{}+\frac{q^F\theta ^F}{(F)_{q^1}!}_z^{},$$
(9)
from these, it is found that an infinitesimal superconformal transformation is of the form
$`z^{}`$ $`=`$ $`z+\omega _0(z)+{\displaystyle \frac{q^F\theta ^F}{(F)_{q^1}!}}ϵ_0(z),`$ (10)
$`\theta ^{}`$ $`=`$ $`\theta +ϵ_0(z)+{\displaystyle \frac{1}{F+1}}\theta \omega _0^{}(z)+{\displaystyle \underset{k=2}{\overset{F}{}}}\theta ^kϵ_k(z),`$ (11)
where $`\omega _0^{}(z):=_z\omega _0(z)`$. It is also seen that the following commutation relations hold
$$ϵ_i\theta =q\theta ϵ_i,i1.$$
(12)
One can extend these naturally to functions of z and $`\overline{z}`$, and $`\theta `$ and $`\overline{\theta }`$ (full functions instead of chiral ones). It is sufficient to define a covariant derivative for the pair $`(\overline{z},\overline{\theta })`$, the analogue of (3), and extend the transformations (5), so that there are similar transformations for $`(\overline{z},\overline{\theta })`$ as well. Then, defining a superconformal transformation as one satisfying (6) and its analogue for $`(\overline{z},\overline{\theta })`$, one obtains, in addition to (8) and (9), similar expressions where $`(z,\theta ,\omega _k,ϵ_k)`$ are simply replaced by $`(\overline{z},\overline{\theta },\overline{\omega }_k,\overline{ϵ}_k)`$. So, the superconformal transformations consists of two distinct class of transformations, the holomorphic and the antiholomorphic, that do not talk to each other.
## 3 generators of superconformal field theory
The (chiral) superfield $`\varphi (z,\theta )`$, with the expansion,
$$\varphi (z,\theta )=\phi _0(z)+\theta \phi _1(z)+\theta ^2\phi _2(z)+\mathrm{}+\theta ^F\phi _F(z),$$
(13)
is a super–primaryfield of the weight $`\mathrm{\Delta }`$, if it transforms under a superconformal transformation as
$$\varphi (z,\theta )(D\theta ^{})^{(F+1)\mathrm{\Delta }}\varphi (z^{},\theta ^{}).$$
(14)
One can write this as
$$\varphi (z,\theta )[1+\stackrel{~}{T}(\omega _0)+\stackrel{~}{S}(ϵ_0)+\underset{k=2}{\overset{F}{}}\stackrel{~}{H}_k(ϵ_k)]\varphi (z^{},\theta ^{}),$$
(15)
to arrive at
$`\stackrel{~}{T}(\omega _0)`$ $`=`$ $`\omega _0_z+(\mathrm{\Delta }+{\displaystyle \frac{\mathrm{\Lambda }}{F+1}})\omega _0^{},`$ (16)
$`\stackrel{~}{S}(ϵ_0)`$ $`=`$ $`ϵ_0\left[\delta _\theta +{\displaystyle \frac{\theta ^F}{(F)_{q^1}!}}_z\right]+{\displaystyle \frac{F+1}{(F)_{q^1}}}\mathrm{\Delta }ϵ_0^{}\theta ^F,`$ (17)
$`\stackrel{~}{H}_k(ϵ_k)`$ $`=`$ $`\theta ^kϵ_k\delta _\theta .`$ (18)
Here $`\mathrm{\Lambda }`$ and $`\delta _\theta `$, are operators satisfying
$$[\mathrm{\Lambda },\theta ]=\theta ,[\mathrm{\Lambda },\delta _\theta ]=\delta _\theta ,$$
(19)
and
$$\delta _\theta \theta =q^F\theta \delta _\theta +1.$$
(20)
One can now define the classical generators
$`l_n`$ $`:=`$ $`\stackrel{~}{T}(z^{n+1}),`$ (21)
$`g_r`$ $`:=`$ $`\stackrel{~}{S}(z^{r+1/(F+1)}),`$ (22)
where, $`n`$ and $`r+1/(F+1)`$ are integers. We do not consider the generators corresponding to $`\stackrel{~}{H}`$, since there is no closed subalgebra, with a trivial central extension, containing these generators . The quantum generators of superconformal transformations are defined through
$`[L_n,\varphi (z,\theta )]`$ $`:=`$ $`l_n\varphi ,`$ (23)
$`[G_r,\varphi (z,\theta )]`$ $`:=`$ $`g_r\varphi .`$ (24)
One can chek that, apart from a possible central extension, these generator satisfy the following relations.
$`[L_n,L_m]`$ $`=`$ $`(nm)L_{n+m},`$ (25)
$`[L_n,G_r]`$ $`=`$ $`\left({\displaystyle \frac{n}{F+1}}r\right)G_{n+r},`$ (26)
and
$$\{G_{r_0}G_{r_1}\mathrm{}G_{r_F}\}_{per}=(F+1)!L_{\left(_{k=0}^Fr_k\right)}$$
(27)
where $`\{\mathrm{}\}_{per}`$, means, the sum of products of all possible permutations of generators $`G_{r_k}`$. This algebra has nontrivial central extensions, it is shown that there is only one subalgebra (containing $`G_r`$’s as well as $`L_n`$’s), the central extension for which is trivial. This algebra is the one generated by $`[L_1,L_0,G_{1/(F+1)}]`$ and their antiholomorphic counterparts $`[\overline{L}_1,\overline{L}_0,\overline{G}_{1/(F+1)}]`$. Having the effect of $`L_n`$’s and $`G_r`$’s on the superfield, it is not difficult to obtain their effect on the component fields. For $`L_n`$’s, the first equation of (18) leads directly to
$$[L_n,\phi _k]=z^{n+1}_z\phi _k+(n+1)z^n(\mathrm{\Delta }+\frac{k}{F+1})\phi _k.$$
(28)
This shows that the component field $`\phi _k`$, is simply a primary field with the weight $`\mathrm{\Delta }+\frac{k}{F+1}`$. One can write (22), also in terms of the operator–product expansion:
$$\mathrm{}[T(\omega )\phi _k(z)]\frac{_z\phi _k(z)}{\omega z}+\frac{(\mathrm{\Delta }+\frac{k}{F+1})\phi _k(z)}{(\omega z)^2},$$
(29)
where $`\mathrm{}`$ denotes the radial ordering and $`T(z)`$, is the holomorphic part of the energy–momentum tensor:
$$T(z)=\underset{n}{}\frac{L_n}{z^{n+2}}.$$
(30)
For $`G_r`$’s, a little more care is needed. One defines a $`\chi `$–commutator as
$$[A,B]_\chi :=AB\chi BA.$$
(31)
It is easy to see that
$$[A,BB^{}]_{\chi \chi ^{}}=[A,B]_\chi B^{}+\chi B[A,B^{}]_\chi ^{}.$$
(32)
Now, if we use
$$[G,\theta ]_q=0,$$
(33)
then the second equation of (18) leads to
$`[G_r,\phi _k]_{q^k}`$ $`=`$ $`z^{r+\frac{1}{F+1}}q^k(k+1)_{q^1}\phi _{k+1},0kF1,`$ (34)
$`[G_r,\phi _F]_{q^F}`$ $`=`$ $`{\displaystyle \frac{q^F}{(F)_{q^1}!}}[z^{r+\frac{1}{F+1}}_z\phi _0`$ (36)
$`+(F+1)(r+{\displaystyle \frac{1}{F+1}})\mathrm{\Delta }z^{r\frac{F}{F+1}}\phi _0].`$
This can also be written in terms of the operator–product expansion. To do this, however, one should first define a proper radial ordering for the supersymmetry generator and the component fields. Defining
$$\mathrm{}[S(w)\varphi _k(z)]:=\{\begin{array}{cc}S(w)\varphi _k(z),\hfill & |w|>|z|\hfill \\ q^k\varphi _k(z)S(w),\hfill & |w|<|z|\hfill \end{array}$$
(37)
where
$$S(z):=\underset{r}{}\frac{G_r}{z^{r+\frac{F+2}{F+1}}},$$
(38)
one arrives at
$`\mathrm{}[S(\omega )\phi _k(z)]`$ $``$ $`q^k(k+1)_{q^1}{\displaystyle \frac{\phi _{k+1}}{\omega z}},0kF1,`$ (39)
$`\mathrm{}[S(\omega )\phi _F(z)]`$ $``$ $`{\displaystyle \frac{q^F}{(F)_{q^1}!}}\left[{\displaystyle \frac{_z\phi _0}{\omega z}}+{\displaystyle \frac{(F+1)\mathrm{\Delta }\phi _0}{(\omega z)^2}}\right].`$ (40)
What we really use to restrict the correlation functions is that part of the algebra the central extension of which is trivial, that is, the algebra generated by $`[L_1,L_0,G_{1/(F+1)}]`$ and their antiholomorphic counterparts $`[\overline{L}_1,\overline{L}_0,\overline{G}_{1/(F+1)}]`$.
## 4 Two–point functions
The two–point functions should be invariant under the action of the subalgebra generated by $`L_1`$,$`L_0`$, and $`G_{1/(F+1)}`$. This means
$`0|[L_1,\varphi _k\varphi _k^{}^{}]|0`$ $`=`$ $`0,`$ (41)
$`0|[L_0,\varphi _k\varphi _k^{}^{}]|0`$ $`=`$ $`0,`$ (42)
$`0|[G_{1/(F+1)},\varphi _k\varphi _k^{}^{}]_{q^{kk^{}}}|0`$ $`=`$ $`0.`$ (43)
We have used the shorthand notation $`\varphi _k=\varphi _k(z)`$ and $`\varphi _k^{}^{}=\varphi _k^{}^{}(z^{})`$. $`\varphi `$ and $`\varphi ^{}`$ are primary superfields of weight $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$, respectively. Equations (32) and (33) imply
$$\varphi _k\varphi _k^{}^{}=\frac{A_{k,k^{}}}{(zz^{})^{\mathrm{\Delta }+\mathrm{\Delta }^{}+(k+k^{})/(F+1)}}.$$
(44)
This is simply due to the fact that $`\varphi _k`$ and $`\varphi _k^{}^{}`$ are primary fields of the weight $`\mathrm{\Delta }+k/(F+1)`$ and $`\mathrm{\Delta }^{}+k^{}/(F+1)`$, respectively. Note that it is not required that these weights be equal to each other, since we have not included $`L_1`$ in the subalgebra.
Relation (34) relates $`A_{k_1,k_1^{}}`$ with $`A_{k_2,k_2^{}}`$, if
$$k_1+k_1^{}(k_2+k_2^{})=0,\text{mod }F+1.$$
(45)
Therefore, there remains $`F+1`$ independent constants in the $`(F+1)^2`$ correlation functions. So there are correlation functions of grade 0, 1, 2, $`\mathrm{}`$, and $`F`$. We have the following relations between the constants $`A_{k,k^{}}`$’s:
$`A_{k+1,k^{}}`$ $`=`$ $`q^k^{}{\displaystyle \frac{1q^{(k^{}+1)}}{1q^{(k+1)}}}A_{k,k^{}+1},0k,k^{}F1,`$ (46)
$`A_{k+1,F}`$ $`=`$ $`{\displaystyle \frac{q(\mathrm{\Delta }+\mathrm{\Delta }^{}+\frac{F+k}{F+1})}{(F)_{q^1}!(k+1)_{q^1}}}A_{k,0},0kF1,`$ (47)
$`A_{F,k^{}+1}`$ $`=`$ $`{\displaystyle \frac{q^k^{}(\mathrm{\Delta }+\mathrm{\Delta }^{}+\frac{F+k^{}}{F+1})}{(F)_{q^1}!(k^{}+1)_{q^1}}}A_{0,k^{}},`$ (48)
$`A_{F,0}`$ $`=`$ $`q^FA_{0,F}.`$ (49)
Using these, one can write (35) as
$$\varphi _k\varphi _k^{}^{}=A_{k+k^{}}f_{k,k^{}}(zz^{}),$$
(50)
where
$$A_{k+k^{}}:=A_{k+k^{},0}=A_K.$$
(51)
So far, everything has been calculated for the chiral fields. But the genralization of this to full fields is not difficult. Following and using exactly the same reasoning, it is seen that
$$\phi _{k\overline{k}}(z,\overline{z})\phi _{k^{}\overline{k}^{}}^{}(z^{},\overline{z}^{})=A_{K\overline{K}}q^{k\overline{k}}f_{k,k^{}}(zz^{})\overline{f}_{\overline{k},\overline{k}^{}}(\overline{z}\overline{z}^{}).$$
(52)
Here $`\overline{f}`$ is the same as $`f`$ with $`\mathrm{\Delta }\overline{\mathrm{\Delta }}`$ and $`\mathrm{\Delta }^{}\overline{\mathrm{\Delta }}^{}`$, and
$$K=k+k^{}\text{mod }F+1,\overline{K}=\overline{k}+\overline{k}^{}\text{mod }F+1.$$
(53)
## 5 Logarithmic two–point functions
Suppose that the first component–field $`\phi _0(z)`$ of the chiral superprimary field $`\mathrm{\Phi }(z,\theta )`$, has a logarithmic counterpart $`\phi _0^{}(z)`$:
$$[L_n,\phi _0^{}(z)]=[z^{n+1}_z+(n+1)z^n\mathrm{\Delta }]\phi _0^{}(z)+(n+1)z^n\phi _0(z).$$
(54)
Following and , one can show that $`\phi _0^{}(z)`$ is the first component–field of a new superfield $`\varphi ^{}(z,\theta )`$, which is the formal derivative of the superfield $`\varphi (z,\theta )`$. One defines the fields $`f_r^{}(z)`$ through
$$[G_r,\phi _0^{}(z)]=:z^{r+1/(F+1)}f_r^{}(z).$$
(55)
Then, by a reasoning similar to that presented in and , that is acting on both sides by $`L_m`$ and using the genralized Jacobi identity, it is seen that all $`f_r`$’s are the same. Denoting this filed by $`\phi _1^{}`$, we have
$$[G_r,\phi _0^{}(z)]=:z^{r+1/(F+1)}\phi _1^{}(z).$$
(56)
One can continue and construct other component–fields. It is found that
$`[G_r,\phi _k^{}]_{q^k}`$ $`=`$ $`z^{r+\frac{1}{F+1}}q^k(k+1)_{q^1}\phi _{k+1}^{},0kF1,`$ (57)
$`[G_r,\phi _F^{}]_{q^F}`$ $`=`$ $`{\displaystyle \frac{q^F}{(F)_{q^1}!}}[z^{r+\frac{1}{F+1}}_z\phi _0^{}`$ (60)
$`+(F+1)\left(r+{\displaystyle \frac{1}{F+1}}\right)\mathrm{\Delta }z^{r\frac{F}{F+1}}\phi _0^{}`$
$`+(F+1)(r+{\displaystyle \frac{1}{F+1}})z^{r\frac{F}{F+1}}\phi _0],`$
and
$$[L_n,\phi _k^{}]=z^{n+1}_z\phi _k^{}+(n+1)z^n(\mathrm{\Delta }+\frac{k}{F+1})\phi _k^{}+(n+1)z^n\phi _k.$$
(61)
Combining the primed component–fields in the chiral superfield $`\varphi ^{}`$, one can write the action of $`L_n`$’s and $`G_r`$’s on it as
$`[L_n,\varphi ^{}]`$ $`=`$ $`\left[z^{n+1}_z+(n+1)z^n\left(\mathrm{\Delta }+{\displaystyle \frac{\mathrm{\Lambda }}{F+1}}\right)\right]\varphi ^{}+(n+1)z^n\varphi ,`$ (62)
$`[G_r,\varphi ^{}]`$ $`=`$ $`z^{r+\frac{1}{F+1}}\left[\delta _\theta +{\displaystyle \frac{\theta ^F}{(F)_{q^1}!}}_z\right]\varphi ^{}`$ (63)
$`+{\displaystyle \frac{F+1}{(F)_{q^1}!}}\left(r+{\displaystyle \frac{1}{F+1}}\right)z^{r\frac{F}{F+1}}\mathrm{\Delta }\theta ^F\varphi ^{}`$
$`+{\displaystyle \frac{F+1}{(F)_{q^1}!}}\left(r+{\displaystyle \frac{1}{F+1}}\right)\theta ^F\varphi .`$
It is easy to see that these are formal derivatives of (18) with respect to $`\mathrm{\Delta }`$, provided one defines
$$\varphi ^{}(z,\theta )=\frac{d\varphi }{d\mathrm{\Delta }}.$$
(64)
The two superfields $`\varphi `$ and $`\varphi ^{}`$, thus combine in a two–dimensional Jordanian block of quasi–primary fields. The generalization of the above results to an $`m`$–dimensional Jordanian block:
$`[L_n,\varphi ^{(i)}]`$ $`=`$ $`\left[z^{n+1}_z+(n+1)z^n\left(\mathrm{\Delta }+{\displaystyle \frac{\mathrm{\Lambda }}{F+1}}\right)\right]\varphi ^{(i)}`$ (66)
$`+(n+1)z^n\varphi ^{(i1)},1im1,`$
and
$`[G_r,\varphi ^{(i)}]`$ $`=`$ $`z^{r+\frac{1}{F+1}}\left[\delta _\theta +{\displaystyle \frac{\theta ^F}{(F)_{q^1}}}_z\right]\varphi ^{(i)}`$ (69)
$`+{\displaystyle \frac{F+1}{(F)_{q^1}!}}\left(r+{\displaystyle \frac{1}{F+1}}\right)z^{r\frac{F}{F+1}}\mathrm{\Delta }\theta ^F\varphi ^{(i)}`$
$`+{\displaystyle \frac{F+1}{(F)_{q^1}!}}\left(r+{\displaystyle \frac{1}{F+1}}\right)z^{r\frac{F}{F+1}}\theta ^F\varphi ^{(i1)},1im1.`$
One can regard these as formal derivatives of (18) with respect to $`\mathrm{\Delta }`$, provided it is defined
$$\varphi ^{(i)}=\frac{1}{i!}\frac{d^i\varphi ^{(0)}}{d\mathrm{\Delta }^i}.$$
(70)
It is then easy to see that
$$\phi _{}^{(i)}{}_{k}{}^{}\phi _{}^{(j)}{}_{k^{}}{}^{}=\frac{1}{i!}\frac{1}{j!}\frac{d^i}{d\mathrm{\Delta }^i}\frac{d^j}{d\mathrm{\Delta }^j}\phi _{}^{(0)}{}_{k}{}^{}\phi _{}^{(0)}{}_{k^{}}{}^{}.$$
(71)
The corralator in the left–hand side is obtained from (38), and in differentiating with respect to the weights one should regard the constants $`A_K`$ as functions of the weights.
One can similarly define Jordanian blocks of full–fields. This generalization is obvious:
$$\varphi ^{(ij)}(z,\overline{z},\theta ,\overline{\theta })=\frac{1}{i!}\frac{1}{j!}\frac{d^i}{d\mathrm{\Delta }^i}\frac{d^j}{d\mathrm{\Delta }^j}\varphi ^{(0)}(z,\overline{z},\theta ,\overline{\theta }).$$
(72)
The correlators of the component–fields are then
$$\phi _{k\overline{k}}^{(ij)}\phi _{k^{}\overline{k}^{}}^{(lm)}=\frac{1}{i!j!l!m!}\frac{d^i}{d\mathrm{\Delta }^i}\frac{d^j}{d\overline{\mathrm{\Delta }}^j}\frac{d^l}{d\mathrm{\Delta }^l}\frac{d^m}{d\overline{\mathrm{\Delta }}^m}\phi _{k\overline{k}}^{(00)}(z,\overline{z})\phi _{k^{}\overline{k}^{}}^{(00)}(z^{},\overline{z}^{}).$$
(73)
The correlator in the right–hand side is obtained from (40), and in differentiating with respect to the weights, $`A_{K\overline{K}}`$’s are regarded as functions of the weights.
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# Rotational Dynamics of Vortices in Confined Bose-Einstein Condensates
## I Introduction
Ever since Bose-Einstein condensation was first observed in a dilute atomic gas , considerable attention has been devoted to understanding its rotational properties. Numerous connections have already been made between superfluid phenomena in the weakly interacting atomic gases and superfluid phenomena observed in condensed matter systems (e.g. ). Recent demonstrations of the ability to create and observe singly quantized vortices , and vortex arrays , have opened up the possibility for investigations of vortex dynamics. Vortex dynamics in this context refers to the motion of topological defects within the superfluid.
A quantized vortex is represented by a singularity in a superfluid order parameter. Since the order parameter must be single-valued, the condition for quantized circulation derives from the fact that the superfluid phase $`\varphi `$ must undergo a $`2\pi l`$ change around any closed contour, where $`l`$ is an integer. Consequently, the velocity field in a superfluid is irrotational everywhere except at a vortex defect or singularity. The density at the defect is zero so that the current density vanishes even though the superfluid velocity $`𝒗`$ at that coordinate is infinite. The superfluid velocity is found from the gradient of the phase by $`𝒗=(\mathrm{}/m)\varphi `$, where $`m`$ is the mass . The vortex core is the region around the defect in which the density falls from its asymptotic value to zero. The spatial scale for the core is characterized by the “healing” length $`\zeta =(8\pi \rho a)^{1/2}`$, where $`\rho `$ is the local superfluid number density and $`a`$ is the $`s`$-wave scattering length characterizing the interactions.
It is anticipated that experimental observations on vortex dynamics in dilute atomic Bose-Einstein condensates should agree quantitatively with predictions of the mean-field theory. External parameters such as the trap properties and the density and composition of the cloud can typically be controlled with a high degree of precision. The positions of vortex cores can be directly observed by imaging techniques rather than indirectly inferred. These features, along with possibility for multiple condensate components to be simultaneously present, make the dilute gas experiments ideal candidates for studies of vortex dynamics.
In this paper, we derive analytic solutions of the equations of motion for vortices with unit angular momentum and make comparison with full numerical simulations. The paper is outlined as follows. We begin by examining in detail a model system composed of a straight line vortex in a uniform density single-component superfluid which is confined in an infinite cylindrical vessel. Analytic solutions are presented for this well-known system for the motion of the vortex defect and connections are made with numerical calculations of the free energy. We extend these results to consider the implications of a harmonic confining potential with the associated quadratic Thomas-Fermi density envelope.
Analysis of such model systems allows us to elucidate the systematic method for generalization to more complicated systems in which analytic solutions are not easily tractable. In particular we investigate numerically the two-component condensate (relevant to current experiments at JILA for example) in which the effects of “buoyancy” must be considered. Buoyancy in this context is used to refer descriptively to a net mean-field force on a constituent of a multi-component condensate due to the various interspecies and intraspecies interaction parameters.
## II Mean-Field Theory
For a single-component condensate, we derive the motion of the vortex defect by solving the evolution of the superfluid order parameter $`\mathrm{\Phi }(𝒓)`$ according to the Gross-Pitaevskii equation
$$i\mathrm{}\frac{d\mathrm{\Phi }(𝒓)}{dt}=\left(\frac{\mathrm{}^2}{2m}^2+V(𝒓)+g|\mathrm{\Phi }(𝒓)|^2\right)\mathrm{\Phi }(𝒓)$$
(1)
where $`V(𝒓)`$ is the external potential and $`g=4\pi \mathrm{}^2a/m`$. The initial condition is found by evaluating the lowest energy solution to the time-independent form of the Gross-Pitaevskii equation consistent with a given total number of atoms and a given position of the vortex defect. This requires minimizing the free energy $`E`$ where
$`E`$ $`=`$ $`{\displaystyle d^3𝒓\left(\frac{\mathrm{}^2}{2m}|\mathrm{\Phi }(𝒓)|^2+V(𝒓)|\mathrm{\Phi }(𝒓)|^2+\frac{g}{2}|\mathrm{\Phi }(𝒓)|^4\right)}`$ (2)
with the imposed constraints.
In the case of a two-component condensate, these equations are modified to account for the fact that the order parameter is a spinor $`(\mathrm{\Phi }_1(𝒓),\mathrm{\Phi }_2(𝒓))`$. The free energy is then given by
$`E`$ $`=`$ $`{\displaystyle }d^3𝒓({\displaystyle \frac{\mathrm{}^2}{2m}}(|\mathrm{\Phi }_1(𝒓)|^2+|\mathrm{\Phi }_2(𝒓)|^2)`$ (5)
$`+V_1(𝒓)|\mathrm{\Phi }_1(𝒓)|^2+V_2(𝒓)|\mathrm{\Phi }_2(𝒓)|^2`$
$`+{\displaystyle \frac{g_{11}}{2}}|\mathrm{\Phi }_1(𝒓)|^4+g_{12}|\mathrm{\Phi }_1(𝒓)\mathrm{\Phi }_2(𝒓)|^2+{\displaystyle \frac{g_{22}}{2}}|\mathrm{\Phi }_2(𝒓)|^4)`$
In the case of <sup>87</sup>Rb, which we will focus on here, the relevant matrix elements which characterize the interspecies and intraspecies interactions of the condensates in the applicable hyperfine states have similar values as indicated by the relationships $`g_{11}=0.97g_{12}`$ and $`g_{22}=1.03g_{12}`$.
For both the single-component case and the two-component case, we minimize the free energy using a steepest descents algorithm . This involves propagating the Gross-Pitaevskii equation in imaginary time (by making the simple substitution $`tit`$), and adjusting the order parameter at each step in the propagation to account for the imposed constraints. The algorithm is straightforward to implement and converges efficiently on the self-consistent lowest energy solution.
The constraints we impose depend on the symmetry of the state we wish to investigate. Since the imaginary time propagation does not preserve normalization, after each numerical step it is necessary to renormalize the order parameter to give the correct total number of atoms $`N`$. For a single-component condensate, the condition is $`N=d^3𝒓|\mathrm{\Phi }(𝒓)|^2`$. Placing the vortex defect at a given location is implemented by imposing a phase pattern about the chosen point. A single unit of circulation requires a $`2\pi `$ phase wrap in the order parameter. We therefore enforce that at each numerical step the order parameter has a complex phase argument of $`\mathrm{exp}(i\theta )`$ about the vortex defect, where $`\theta `$ is the usual counterclockwise angle measured in the plane perpendicular to the vortex line.
For a two-component condensate the method used is similar. Given a proportion $`p`$ of atoms in the first condensate component, the normalization constraints are $`Np=d^3𝒓|\mathrm{\Phi }_1(𝒓)|^2`$ and $`N(1p)=d^3𝒓|\mathrm{\Phi }_2(𝒓)|^2`$. We adopt the convention that the component 2 is the state which will contain the single vortex line, while component 1 will contain no vortices and will therefore tend to fill the vortex core of component 1 to plug the hole in the density. In order to approximate this situation using the method of steepest descents, we enforce that the phase of $`\mathrm{\Phi }_1(𝒓)`$ is spatially uniform at each numerical step. The phase of component 2 is fixed at each numerical step to give a $`2\pi `$ phase winding around the chosen position of the defect in the same manner described for a single-component condensate.
For simplicity, we take the system to be translationally invariant along the dimension parallel to the vortex line. This allows us to solve the Gross-Pitaevskii equation in two dimensions rather than in three dimensions. To allow for comparison with experiment, the number of atoms per unit length in the two-dimensional Gross-Pitaevskii equation should take the value which reproduces the same chemical potential as the equivalent total number of atoms in the real three-dimensional system.
## III Single-Component Condensate
We begin our conceptual treatment of vortex dynamics in model systems with simple cases involving the motion of a quantized vortex line in a cylinder.
### A Uniform Density Distribution
Consider a uniform density superfluid confined in an infinite cylindrical vessel of radius $`R`$. A vortex is placed in the superfluid with the vortex line displaced from the cylinder axis by $`𝒓_\mathrm{𝟎}`$. We define the circulation $`𝜿`$ in the usual way to have magnitude $`2\pi l\mathrm{}/m`$ and to be aligned parallel to the vortex line with direction determined according to the usual right-hand rule applied to the superfluid flow. The velocity field can be found using an image vortex argument . The effect of the boundary conditions at the cylinder walls is to require that the perpendicular component of the superfluid velocity is zero at the surface. This condition is satisfied by considering a formally equivalent situation of a uniform fluid of infinite extent with an additional image vortex of opposite circulation placed at $`𝒓_\mathrm{𝟏}=(R/r_0)^2𝒓_\mathrm{𝟎}`$. The velocity field for this situation is:
$`𝒗(𝒓)`$ $`=`$ $`{\displaystyle \frac{𝜿}{2\pi }}\times \left({\displaystyle \frac{𝒓𝒓_\mathrm{𝟎}}{|𝒓𝒓_\mathrm{𝟎}|^2}}{\displaystyle \frac{𝒓𝒓_\mathrm{𝟏}}{|𝒓𝒓_\mathrm{𝟏}|^2}}\right).`$ (6)
The motion of the vortex defect at $`𝒓_\mathrm{𝟎}`$ is found by computing the superfluid flow at that coordinate, which is due solely to the image vortex contribution:
$$𝒗(𝒓_\mathrm{𝟎})=\frac{𝜿\times 𝒓_\mathrm{𝟎}}{2\pi (R^2r_0^2)}$$
(7)
The angular frequency of precession of the vortex defect about the center axis of the cylinder is therefore given by $`𝝎=𝜿/2\pi (R^2r_0^2)`$. Note that the angular precession direction is always in the same sense as the circulation.
An alternative approach (but equivalent method) for finding the rate and direction of precession of the vortex defect is to compute the free energy due to the attractive interaction between the real vortex and the image vortex. Ignoring the kinetic energy within the core of the vortex line, the free energy per unit length is given by
$$E=\frac{\rho \kappa ^2m}{4\pi }\mathrm{log}\left[\frac{R^2r_0^2}{R\zeta }\right]$$
(8)
The defect velocity $`𝒗`$ is then found from the solution of
$$\rho m(𝜿\times 𝒗)=E$$
(9)
where $`E`$ is the gradient of the free energy with respect to the location of the defect, i.e. $`𝒓_\mathrm{𝟎}`$. The importance of this alternative approach is that, for a complex system, the free energy surface $`E`$ can be calculated numerically. Consequently, the implication is that one may find the gradient of the free energy with respect to displacement of the defect as a model to infer the behavior of the vortex dynamics even when a simple analytic expression such as Eq. (8) cannot easily be derived.
In order to illustrate this point, in Fig. 1 we show the results of a numerical solution of the free energy for a displaced vortex in a cylinder. The external potential $`V(𝒓)`$ for this case changes abruptly from zero inside the cylinder to infinity at the walls. The resulting superfluid density is approximately uniform except within a small distance of the surfaces as characterized by the healing length. The exact numerical results are compared with the approximate analytic free energy expression given in Eq. (8). According to the analytic expression, the free energy diverges at the edge of the cloud, since the boundary effects associated with the core size have not been taken into account. Apart from edge effects, the analytic expression for the free energy agrees very well with the total energy found from the numerical solution of the free energy.
At zero temperature there is no energy dissipation and the motion of the vortex defect is along an equipotential of the free energy, which is circular in this case. However, in the presence of dissipation, the defect propagates to regions of lower free energy. The vortex and image vortex attract each other in the uniform fluid since they have the opposite sign for the circulation. Consequently, for the situation considered in Fig. 1 the vortex will spiral out and eventually annhilate with the image vortex at the surface of the cylinder. At finite temperature, dissipation is generated by collisions between superfluid atoms and atoms from the non-condensed component of the cloud.
### B Thomas-Fermi Density Distribution
We now replace the superfluid of uniform density by the quadratic Thomas-Fermi density envelope. This is a good approximation to the density distribution which results from a harmonic confining potential and allows us to make contact with an experimentally more relevant situation for dilute atomic Bose-Einstein condensates. For this distribution, the density falls to zero at the Thomas-Fermi radius $`R`$ in a smooth and continuous manner so the importance of edge effects is reduced.
The free energy surface results from the velocity field associated with the minimum energy configuration at each value of the radial core displacement and this may be calculated analytically in a hydrodynamic approach . This approach is in the same spirit as the method used for the derivation of the free energy in Eq. (8) in that the kinetic energy associated with the core region, and the spatial dependence of the trapping potential and mean-field potential across the vortex core are neglected. Taking into account the Thomas-Fermi density distribution in this way gives the free energy expression
$`E`$ $`=`$ $`{\displaystyle \frac{\rho _0\kappa ^2m}{8\pi }}[{\displaystyle \frac{R^2r_0^2}{R^2}}\mathrm{ln}{\displaystyle \frac{R^2}{\zeta _0^2}}`$ (11)
$`+({\displaystyle \frac{R^2}{r_0^2}}+1{\displaystyle \frac{2r_0^2}{R^2}})\mathrm{ln}{\displaystyle \frac{R^2r_0^2}{R^2}}]`$
where $`\rho _0`$ is the number density of the gas at the center of the trap, and $`\zeta _0`$ is the corresponding healing length.
According to the implication that the gradient of the free energy surface is related to the precession frequency as given in Eq. (9), one may calculate the precession frequency within the approximations associated with Eq. (11) as
$`𝝎`$ $`=`$ $`{\displaystyle \frac{𝜿}{4\pi (R^2r_0^2)}}(2\mathrm{ln}{\displaystyle \frac{R}{\zeta _0}}`$ (13)
$`+({\displaystyle \frac{R^4}{r_0^4}}+2)\mathrm{ln}{\displaystyle \frac{R^2r_0^2}{R^2}}+{\displaystyle \frac{R^2}{r_0^2}}+2)`$
Note that the precession direction for this case is always in the same sense as the circulation, so that the qualitative behavior of the motion of the vortex defect is similar to that of the uniform fluid in a cylinder.
In Fig. 2 we compare the analytic expression for the free energy given in Eq. (11) with the exact numerical solution for the free energy. The agreement is remarkable and the role of edge effects is small on the free energy and consequently on the precession frequency. The sign of the free energy gradient indicates that in the presence of dissipation, the vortex is not energetically stable and will spiral outwards to the edge of the cloud where it will disappear.
## IV Two-Component Condensate
When multiple-components are simultaneously present in a trap, the interactions are characterized by a matrix of the interspecies and intraspecies collisions. Depending on the elements of the scattering matrix, various distinct kinds of behavior are possible for the density distribution of the components.
Condensate experiments on <sup>87</sup>Rb can typically trap two-components simultaneously which tend to phase separate and are approximately immiscible in a non-uniform confining potential. In addition, within the mean-field approximation, the lowest energy solution has the component with the smaller self-interaction scattering length in regions of highest density. If this component is displaced it will tend to float to the center of the trap, where in a harmonic potential the density is greatest. A pictorial analogy is often made between this physical mechanism and a buoyancy force in a fluid of varying density.
In terms of the implications for the vortex dynamics, the inclusion of this second component can be important even if the second component contains no vortex lines itself. The extra degree of freedom associated with the buoyancy behavior allows a more complex structure for the free energy surface.
In Fig. 3 we calculate the free energy for the case of $`10^6`$ total <sup>87</sup>Rb atoms in the condensate in an isotropic 7.8 Hz trap with about 40% in component 1 and 60% in component 2. The scattering parameter used is $`g_{12}=4\pi \mathrm{}^2(100a_0)/m`$ where $`a_0`$ is the Bohr radius. As mentioned previously the simulation is actually performed in two-dimensions with a number of particles per unit length set to give the same chemical potential as for the three dimensional system.
The characteristic of a local minimum in the free energy at the center of the trap has significant implications. The minimum implies a critical radius $`R_c`$, which is the radius corresponding to the maximum in the free energy curve. If the displacement of the vortex line is less than $`R_c`$, the vortex will dissipate energy by spiraling inwards to the center of the trap, where it will remain forever. There is no energetic path by which the vortex can propagate to a region of zero density and annihilate. Due to its non-trivial topological structure, annihilation with a vortex of opposite circulation is required to remove a vortex line and this can happen trivially only in a zero-density region. Consequently, this two-component condensate system studied here can support a persistent current which is metastable, provided the vortex line is generated near the trap center. In reality, inelastic processes which have not been included in our discussion and which result from spin-exchange collisions will limit the lifetime of the metastable state.
The complementary situation exists for a vortex line which is created with a displacement larger than $`R_c`$. In the presence of dissipation, such a vortex will spiral out of the system and annihilate at the surface. If the vortex line is displaced from the trap center by exactly $`R_c`$, the forces acting on the vortex defect balance precisely, and the velocity of the singularity in the superfluid flow is zero.
These phenomena require a sufficiently strong influence of the mean-field interactions. If the total number of atoms is reduced from $`10^6`$ to $`4\times 10^5`$ with the same fraction in the non-rotating component, there is no longer a barrier in the free energy, and the critical radius is zero. In this case no energetically stable persistent currents are possible.
The presence of a barrier in the free energy alters qualitatively the behavior of the vortex dynamics. As illustrated in Fig. 4, when the gradient of the free energy changes sign with respect to displacement of the vortex core, so does the direction of precession of the defect, as implied by Eq. (9). Consequently, in a two-component system, the non-rotating component can modify the rate and direction of precession of the core purely through its influence through the mean-field potential.
In Fig. 5 we illustrate density snapshots showing real-time images of the vortex dynamics for the case of $`10^6`$ total atoms. The vortex defect in the first column, which has a displacement less than the critical radius, is precessing clockwise, opposite to the chosen direction of the superfluid circulation. The others are precessing counterclockwise, which is the same sense as the intrinsic vortex fluid flow. For each of these columns, the rate of precession is in accordance with the gradient of the of the free energy curve calculated numerically for the same parameters. In Fig. 6 similar snapshots are shown for the case for $`4\times 10^5`$ total atoms. Here, the defect is precessing counterclockwise at all displacements, and the barrier in the free energy is absent.
## V Conclusion
We have found that for a single-component condensate the precession direction will always be in the same sense as the vortex fluid flow for both the uniform density in a cylinder and for the Thomas-Fermi density profile. Such a vortex is not energetically stable, and in the presence of dissipation, will spiral out to the edge of the cloud and annihilate. Including the effects of kinetic energy in the core region and edge effects at the surface of the cylinder in the free energy gave minor modifications to the vortex dynamics from that predicted by the hydrodynamic theories.
For a two-component condensate, numerical calculations demonstrated the possibility for the precession direction to be opposite to the sense of the vortex fluid flow. In this case, the vortex will dissipate energy by spiraling to the center of the trap. Such a situation occurs whenever there is a minimum in the free energy as a function of displacement of the vortex defect. This possibility allows a metastable persistent current which may remain indefinitely, even in the presence of thermal fluctuations. This is due to the topological nature of the vortex prohibiting annihilation and is a manifestation of superfluidity. When there are insufficient atoms, the absence of an energy barrier causes the vortex to spiral outward. This case is similar to the case of a vortex in a single-component condensate, with a modified precession frequency.
## VI Acknowledgements
We would like to thank A. Fetter, E. Cornell, B. Anderson, and P. Haljan for helpful discussions. This work was supported by the National Science Foundation.
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# Ultra High Energy Cosmic Rays from Compact Sources
## I INTRODUCTION
The interaction of protons with photons of the cosmic microwave background predicts a sharp drop in the cosmic ray flux above the GZK cutoff around $`510^{19}`$ eV . The available data shows no such drop. About 20 events above $`10^{20}`$ eV were observed by a number of experiments such as AGASA , Fly’s Eye , Haverah Park , Yakutsk and HiRes . Since above the GZK energy the attenuation length of particles is a few tens of megaparsecs if an ultra high energy cosmic ray (UHECR) is observed on earth it must be produced in our vicinity (except for UHECR scenarios based on weakly interacting particles, e.g. neutrinos).
Usually it is assumed that at these high energies the galactic and extragalactic magnetic fields do not affect the orbit of the cosmic rays, thus they should point back to their origin within a few degrees. In contrast to the low energy cosmic rays one can use UHECRs for point-source search astronomy. (For an extragalactic magnetic field of $`\mu `$G rather than the usually assumed nG there is no directional correlation with the source . Note, that due to the Local Supercluster the magnetic field is presumably not less than 10 nG which results in a Larmor radius of few tens of megaparsecs for protons above $`10^{20}`$ eV.) Though there are some peculiar clustered events, which we discuss in detail, the overall distribution of UHECR on the sky is practically isotropic. This observation is rather surprising since in principle only a few astrophysical sites (e.g. active galactic nuclei or the extended lobes of radio galaxies ) are capable of accelerating such particles, nevertheless none of the UHECR evets came from these directions. Hence it is generally believed that there is no conventional astrophysical explanation for the observed UHECR spectrum.
There are several ways to look for the source inhomogeneity from the energy spectrum and spatial directions of UHECRs. One possibility is to assume that the source density of UHECRs is proportional to the galaxy densities . Another approach is to analyze the clustering properties of the unknown sources by some correlation length.
Clearly, the arrival directions of the UHECRs measured by experiments show some peculiar clustering: some events are grouped within $`3^o`$, the typical angular resolution of an experiment. Above $`410^{19}`$ eV 92 cosmic ray events were detected, including 7 doublets and 2 triplets. Above $`10^{20}`$ eV one doublet out of 14 events were found . The chance probability of such a clustering from uniform distribution is rather small . (Taking the average bin $`3^o`$ the probability of generating one doublet out of 14 events is $`11\%`$.)
The clustered features of the events initiated an interesting statistical analysis assuming compact UHECR sources . The authors found a large number, $`400`$ for the number of sources<sup>*</sup><sup>*</sup>*approximately $`400`$ sources within the GZK sphere results in one doublet for 14 events. The order of magnitude of this result is in some sense similar to that of a “high-school” exercise: what is the minimal size of a class for which the probability of having clustered birthdays –at least two pupils with the same birthdays– is larger than 50%. In this case the number of “sources” is the number of possible birthdays $`400`$. In order to get the answer one should solve $`365!/[365^k(365k)!]<0.5`$, which gives as a minimal size k=23. inside a GZK sphere of 25 Mpc. They assumed that a.) the number of clustered events is much smaller than the total number of events (this is a reliable assumption at present statistics; however, for any number of sources the increase of statistics, which will happen in the near future, results in more clustered events than unclustered), b.) all sources have the same luminosity which gives a delta function for their distribution (this unphysical choice represents an important limit, it gives the smallest source density for a given number of clustered and unclustered events) c.) The GZK effect makes distant sources fainter; however, this feature depends on the injected energy spectrum and the attenuation lengths and elasticities of the propagating particles. In an exponential decay was used with an energy independent decay length of 25Mpc.
In our approach none of these assumptions are used. In addition we include spherical astronomy corrections and in particular give the upper and lower bounds for the source density at a given confidence level. As we show the most probable value for the source density is really large; however, the statistical significance of this result is rather weak. At present the small number of UHECR events allows a 95% confidence interval for the source density which spreads over four orders of magnitude. Since future experiments, particularly Pierre Auger , will have a much higher statistical significance on clustering (the expected number of events of $`10^{20}`$ eV and above is 60 per year ), we present our results on the density of sources also for larger number of UHECRs above $`10^{20}`$ eV.
In order to avoid the assumptions of a combined analytical and Monte-Carlo technique will be presented adopting the conventional picture of protons as the ultra high energy cosmic rays. Our analytical approach of Section II gives the event clustering probabilities for any space, luminosity and energy distribution of the sources by using a single additional function $`P(r,E;E_c)`$, the probability that a proton created at a distance $`r`$ with energy $`E`$ arrives at earth above the threshold energy $`E_c`$ . With our Monte-Carlo technique of Section III we determine the probability function $`P(r,E;E_c)`$ for a wide range of parameters. Our results for the present and future UHECR statistics are presented in Section IV We summarize in Section V.
## II ANALYTICAL APPROACH
The key quantity for finding the distribution functions for the source density is the probability of detecting $`k`$ events from one randomly placed source. The number of UHECRs emitted by a source of $`\lambda `$ luminosity during a period $`T`$ follows the Poisson distribution. However, not all emitted UHECRs will be detected. They might loose their energy during propagation or can simply go to the wrong direction.
For UHECRs the energy loss is dominated by the pion production in interaction with the cosmic microwave background radiation. In ref. the probability function $`P(r,E,E_c)`$ was presented for three specific threshold energies. This function gives the probability that a proton created at a given distance from earth (r) with some energy (E) is detected at earth above some energy threshold ($`E_c`$). The resulting probability distribution can be approximated over the energy range of interest by a function of the form
$$P(r,E,E_c)\mathrm{exp}[a(E_c)r^2\mathrm{exp}(b(E_c)/E)]$$
(1)
The appropriate values of $`a`$ and $`b`$ for $`E_c/(10^{20}\mathrm{eV})=`$1,3, and 6 are, respectively $`a/(10^4\mathrm{Mpc}^2)=`$1.4, 9.2 and 11, $`b/(10^{20}eV)=`$2.4, 12 and 28.
For the sources we use the second equatorial coordinate system: $`𝐱`$ is the position vector of the source characterized by ($`r,\delta ,\alpha `$) with $`\delta `$ and $`\alpha `$ beeing the declination and right ascension, respectively. The features of the Poisson distribution enforce us to take into account the fact that the sky is not isotropically observed. There is a circumpolar cone, in which the sources can always be seen, with half opening angle $`\delta ^{}`$ ($`\delta ^{}`$ is the declination of the detector, for the experiments we study $`\delta ^{}40^o50^o`$). There is also an invisible region with the same opening angle. Between them there is a region for which the time fraction of visibility, $`\gamma (\delta ,\delta ^{})`$ is a function of the declination of the source. It is straightforward to determine $`\gamma (\delta ,\delta ^{})`$ for any $`\delta `$ and $`\delta ^{}`$:
$$\gamma (\delta ,\delta ^{})=\{\begin{array}{ccc}0\hfill & \text{if }\hfill & \pi /2<\delta \delta ^{}\pi /2\hfill \\ 1\hfill & \mathrm{arccos}\hfill & (\mathrm{tan}\delta ^{}\mathrm{tan}\delta )/\pi \hfill \\ & \text{if }\hfill & \delta ^{}\pi /2<\delta \pi /2\delta ^{}\hfill \\ 1\hfill & \text{if }\hfill & \pi /2\delta ^{}<\delta \pi /2\hfill \end{array}$$
(2)
To determine the probability that a particle arriving from random direction at a random time is detected we have to multiply $`\gamma (\delta ,\delta ^{})`$ by the cosine of the zenith angle $`\theta `$. In the following we will use the time average of this function:
$$\eta (\delta ,\delta ^{})=\frac{1}{T}_0^T\gamma (\delta ,\delta ^{})\mathrm{cos}\theta (\delta ,\delta ^{},t)𝑑t$$
(3)
Since $`\delta ^{}`$ is constant, in the rest of the paper we do not indicate the dependence on it. Neglecting these spherical astronomy effects means more than a factor of two for the prediction of the source density.
The probability of detecting $`k`$ events from a source at distance $`r`$ with energy $`E`$ can be obtained by including $`P(r,E,E_c)A\eta (\delta )/(4\pi r^2)`$ in the Poisson distribution:
$`p_k(𝐱,E,j)={\displaystyle \frac{\mathrm{exp}\left[P(r,E,E_c)\eta (\delta )j/r^2\right]}{k!}}\times `$ (4)
$`\left[P(r,E,E_c)\eta (\delta )j/r^2\right]^k,`$ (5)
where we introduced $`j=\lambda TA/(4\pi )`$ and $`A\eta (\delta )/(4\pi r^2)`$ is the probability that an emitted UHECR points to a detector of area $`A`$. We denote the space, energy and luminosity distributions of the sources by $`\rho (𝐱)`$, $`c(E)`$ and $`h(j)`$, respectively. The probability of detecting $`k`$ events above the threshold $`E_c`$ from a single source randomly positioned within a sphere of radius $`R`$ is
$`P_k={\displaystyle _{S_R}}dV\rho (𝐱){\displaystyle _{E_c}^{\mathrm{}}}dEc(E){\displaystyle _0^{\mathrm{}}}djh(j)\times `$ (6)
$`{\displaystyle \frac{\mathrm{exp}\left[P(r,E,E_c)\eta (\delta )j/r^2\right]}{k!}}\left[P(r,E,E_c)\eta (\delta )j/r^2\right]^k.`$ (7)
Denote the total number of sources within the sphere of sufficiently large radius (e.g. several times the GZK radius) by $`N`$ and the number of sources that gave $`k`$ detected events by $`N_k`$. Clearly, $`N=_0^{\mathrm{}}N_i`$ and the total number of detected events is $`N_e=_0^{\mathrm{}}iN_i`$. The probability that for $`N`$ sources the number of different detected multiplets are $`N_k`$ is:
$$P(N,\{N_k\})=N!\underset{k=0}{\overset{\mathrm{}}{}}\frac{1}{N_k!}P_k^{N_k}.$$
(8)
The value of $`P(N,\{N_k\})`$ is the most important quantity in our analysis of UHECR clustering. For a given set of unclustered and clustered events ($`N_1`$ and $`N_2,N_3`$,…) inverting the $`P(N,\{N_k\})`$ distribution function gives the most probable value for the number of sources and also the confidence interval for it. If we want to determine the density of sources we can take the limit $`R\mathrm{},N\mathrm{}`$, while the density of sources $`S=N/(\frac{4}{3}R^3\pi )`$ is constant.
In order to illustrate the dominant length scale it is instructive to study the integrand $`f_k(r)`$ of the distance integration in eqn. (6)
$`P_k={\displaystyle _0^R}\left({\displaystyle \frac{dr}{R}}\right)f_k(r),`$
$`f_k(r)=Rr^2{\displaystyle }d\mathrm{\Omega }\rho (𝐱){\displaystyle _{E_c}^{\mathrm{}}}dEc(E){\displaystyle _0^{\mathrm{}}}djh(j)\times `$ (9)
$`{\displaystyle \frac{\mathrm{exp}\left[P(r,E,E_c)\eta (\delta )j/r^2\right]}{k!}}\left[P(r,E,E_c)\eta (\delta )j/r^2\right]^k.`$ (10)
Fig. 1 shows that $`f_1(r)`$, which leads to singlet events, is dominated by the distance scale of 10-15 Mpc, whereas $`f_2(r)`$, which gives doublet events, is dominated by the distance scale of 4-6 Mpc It is interesting that the dominant distance scale for singlet events is by an order of magnitude smaller than the attenuation length of the protons at these energies ($`l_a110`$ Mpc). This surprising result can be illustrated using a simple approximation. Assuming that the probability of detecting a particle coming from distance $`r`$ is proportional to $`\mathrm{exp}(r/l_a)/r^2`$, $`P_1`$ will be proportional to $`𝑑\mathrm{\Omega }𝑑rr^2\mathrm{exp}[j\mathrm{exp}(r/l_a)/r^2]\mathrm{exp}(r/l_a)/r^2`$. For the typical $`j`$ values the $`r`$ integrand has a maximum around 15 Mpc and not at $`l_a`$.. These typical distances partly justify our assumption of neglecting magnetic fields. The deflection of singlet events due to magnetic field does not change the number of multiplets, thus our conclusions. The typical distance for higher multiplets is quite small, therefore deflection can be practically neglected. Clearly, the fact that multiplets are coming from our “close” neighbourhood does not mean that the experiments reflect just the densities of these distances. The overwhelming number of events are singlets and they come from much larger distances. Note, that these $`f_1(r)`$ and $`f_2(r)`$ functions are obtained with our optimal $`j_{}`$ value (cf. Fig. 5 and explanation there and in the corresponding text). Using the largest possible $`j_{}`$ value allowed by the 95% confidence region the dominant distance scales for $`f_1(r)`$ and $`f_2(r)`$ functions turn out to be 30 Mpc and 20 Mpc, respectively.
Note, that $`P_k`$ and then $`P(N,\{N_k\})`$ are easily determined by a well behaving four-dimensional numerical integration (the $`\alpha `$ integral can be factorized) for any $`c(E)`$, $`h(j)`$ and $`\rho (r)`$ distribution functions . In order to illustrate the uncertainties and sensitivities of the results we used a few different choices for these distribution functions.
For $`c(E)`$ we studied three possibilities. The most straightforward choice is the extrapolation of the ‘conventional high energy component’ $`E^2`$. Another possibility is to use a stronger fall-off of the spectrum at energies just below the GZK cutoff, e.g. $`E^3`$. These choices span the range usually considered in the literature and we will study both of them. The third possibility is to assume that topological defects generate UHECRs through production of superheavy particles Note, that these particles are not superheavy DM particles , which are located most likely in the halo of our galaxy. These superheavy DM particles can also be considered as possible sources of UHECR with anisotrpies in the arrival direction .. According to these superheavy particles decay into quarks and gluons which initiate multi-hadron cascades through gluon bremstrahlung. These finally hadronize to yield jets. The energy spectrum, first calculated in for the Standard Model and in for the Minimal Supersymmetric Standard Model, in this case can be estimated by the function obtained from the HERWIG QCD generator:
$$c(x)=c_1\frac{\mathrm{exp}[c_2\sqrt{\mathrm{ln}(1/x)}](1x)^2}{x\sqrt{\mathrm{ln}(1/x)}},$$
(11)
where $`x=E/m_X`$ the ratio of the energy and the mass of the decaying particle. The best fit to the observed UHECR spectrum gives $`m_X10^{12}`$ GeV for the mass of the decaying particle. This corresponds to $`c_10.0086`$ and $`c_22.77`$. We will use these $`c_1`$ and $`c_2`$ values for our third choice of energy distribution, $`c(E)`$.
In ref. the authors have shown that for a fixed set of multiplets the minimal density of sources can be obtained by assuming a delta-function distribution for $`h(j)`$. We studied both this limiting case ($`h(j)=\delta (jj_{})`$) and a more realistic one with Schechter’s luminosity function :
$$h(j)dj=h(j/j_{})^{1.25}\mathrm{exp}(j/j_{})d(j/j_{}).$$
(12)
The space distribution of sources can be given based on some particular survey of the distribution of nearby galaxies or on a correlation length $`r_0`$ characterizing the clustering features of sources . For simplicity the present analysis deals with a homogeneous distribution of sources randomly scattered in the universe (Note, that due to the Local Supercluster the isotropic distribution is just an approximation.).
Fig. 2 shows the resulting $`P_k(j_{})`$ probability functions for the different choices of $`c(E)`$ and $`h(j)`$. The overall shapes of them are rather similar; nevertheless, relatively small differences lead to quite different predictions for the UHECR source density. The “shoulders” of the curves with Dirac-delta luminosity distributions got smoother for the Schechter’s distribution. The scales on the figures are chosen to cover the 98% confidence regions (see section IV for details).
Note, that – assuming that UHECRs point back to their sources – our clustering technique discussed above applies to practically any models of UHECR (e.g. neutrinos). One only needs a change in the $`P(r,E,E_c)`$ probability distribution function (e.g. neutrinos penetrate the microwave background uninhibited) and use the $`h(j)`$ and $`c(E)`$ distribution function of the specific model.
## III MONTE-CARLO STUDY OF THE PROPAGATION
Our Monte-Carlo model of UHECR studies the propagation of UHECR. The analysis of showed that both AGASA and Fly’s Eye data demonstrated a change of composition, a shift from heavy –iron– at $`10^{17}`$ eV to light –proton– at $`10^{19}`$ eV Thus, the chemical composition of UHECRs is most likely to be dominated by protons. In our analysis we use exclusively protons as UHECR particles. (for suggestions that air showers above the GZK cutoff are induced by neutrinos see .)
Using the pion production as the dominant effect of energy loss for protons at energies $`>10^{19}`$ eV ref. calculated $`P(r,E,E_c)`$, the probability that a proton created at a given distance (r) with some energy (E) is detected at earth above some energy threshold ($`E_c`$). For three threshold energies the authors of gave an approximate formula, which we used in the previous section.
In our Monte-Carlo approach we determined the propagation of UHECR on an event by event basis. Since the inelasticity of Bethe-Heitler pair production is rather small ($`10^3`$) we used a continuous energy loss approximation for this process. The inelasticity of pion-photoproduction is much higher ($`0.20.5`$) in the energy range of interest, thus there are only a few tens of such interactions during the propagation. Due to the Poisson statistics of the number of interactions and the spread of the inelasticity, we will see a spread in the energy spectrum even if the injected spectrum is mono-energetic.
In our simulation protons are propagated in small steps ($`10`$ kpc), and after each step the energy losses due to pair production, pion production and the adiabatic expansion are calculated. During the simulation we keep track of the current energy of the proton and its total displacement. This one avoids performing new simulations for different initial energies and distances. The propagation is completed when the energy of the proton goes below a given cutoff. For the proton interaction lengths and inelasticities we used the values of . The deflection due to magnetic field is not taken into account, because it is small for our typical distances illustrated in Fig. 1. This fact justifies our assumption that UHECRs point back to their sources (for a recent Monte-Carlo analysis on deflection see e.g. ).
Since it is rather practical to use the $`P(r,E,E_c)`$ probability distribution function we extended the results of by using our Monte-Carlo technique for UHECR propagation. In order to cover a much broader energy range than the parametrization of (1) we used the following type of function
$$P(r,E,E_c)=\mathrm{exp}\left[a(r/1\mathrm{Mpc})^b\right].$$
(13)
Fig. 3 demonstrates the reliability of this parametrization. The direct Monte-Carlo points and the fitted function (eqn. (13) with $`a=0.0019`$ and $`b=1.695`$) are plotted for $`E_c=10^{20}`$eV and $`E=210^{20}`$eV.
Fig. 4 shows the functions $`a(E/E_c)`$ and $`b(E/E_c)`$ for a range of three orders of magnitude and for five different threshold energies. Just using the functions of $`a(E/E_c)`$ and $`b(E/E_c)`$, thus a parametrization of $`P(r,E,E_c)`$ one can obtain the observed energy spectrum for any injection spectrum without additional Monte-Carlo simulation.
## IV RESULTS
In order to determine the confidence intervals for the source densities we used the frequentist method. We wish to set limits on S, the source density. Using our Monte-Carlo based $`P(r,E,E_c)`$ functions and our analytical technique we determined $`p(N_1,N_2,N_3,\mathrm{};S;j_{})`$, which gives the probability of observing $`N_1`$ singlet, $`N_2`$ doublet, $`N_3`$ triplet etc. events if the true value of the density is $`S`$ and the central value of luminosity is $`j_{}`$. The probability distribution is not symmetric and far from being Gaussian. For a given set of $`\{N_i,i=1,2,\mathrm{}\}`$ the above probability distribution as a function of $`S`$ and $`j_{}`$ determines the 68% and 95% confidence level regions in the $`Sj_{}`$ plane. Fig. 5 shows these regions for our “favorite” choice of model ($`c(E)E^3`$ and Schechter’s luminosity distribution) and for the present statistics (one doublet out of 14 UHECR events). The regions are deformed, thin ellipse-like objects in the $`\mathrm{log}(j_{})`$ versus $`\mathrm{log}(S)`$ plane. Since $`j_{}`$ is a completely unknown and independent physical quantity the source density can be anything between the upper and lower parts of the confidence level regions. For this model our final answer for the density is $`180_{165(174)}^{+2730(8817)}10^3`$ Mpc<sup>-3</sup>, where the first errors indicate the 68%, the second ones in the parenthesis the 95% confidence levels, respectively. The choice of –Dirac-delta like luminosity distribution– and, for instance, conventional $`E^2`$ energy distribution gives much smaller value: $`2.77_{2.53(2.70)}^{+96.1(916)}10^3`$ Mpc<sup>-3</sup>. For other choices of $`c(E)`$ and $`h(j)`$ see Table I. Our results for the Dirac-delta luminosity distribution are in agreement with the result of within the error bars. Neverthless, there is a very important message. The confidence level intervals are so large, that on the 95% confidence level two orders of magnitude smaller densities than suggested as a lower bound by are also possible.
As it can be seen there is a strong correlation between the luminosity and the source density. Physically it is easy to understand the picture. For a smaller source density the luminosities should be larger to give the same number of events. However it is not possible to produce the same multiplicity structure with arbitrary luminosities. Very small luminosities can not give multiplets at all, very large luminosities tend to give more than one doublet.
The same technique can be applied for any hypothetical experimental result. For fixed $`\{N_k\}`$ the above probability function determines the 68% confidence regions in $`S`$ and $`j_{}`$. Using these regions one can tell the 68% confidence interval for S. The most probable values of the source densities for fixed number of multiplets are plotted on Fig. 6 with the lower and upper bounds. The total number of events is shown on the horizontal axis, whereas the number of multiplets label the lines. Here again, our ”favorite” choice of distribution functions were used: $`c(E)E^3`$ and $`h(j)`$ of eqn. (12).
It is of particular interest to analyze in detail the present experimental situation having one doublet out of 14 events. Since there are some new unpublished events, too, we study the a hypothetical case of one or two doublets out of 24 events. The 68% and 95% confidence level results are summarized in Table I for our three energy and two luminosity distributions. It can be seen that Dirac-delta type luminosity distribution really gives smaller source densities than broad luminosity distribution, as it was proven by . Less pronounced is the effect on the energy distribution of the emitted UHECRs. The $`c(E)E^3`$ case gives somewhat larger values than the other two choices ($`c(E)E^2`$ or given by the decay of a superheavy particle). The confidence intervals are typically very large, on the 95% level they span 4 orders of magnitude. An interesting feature of the results is that ”doubling” the present statistics with the same clustering features (in the case studied by the table this means one new doublet out of 10 new events) reduces the confidence level intervals by an order of magnitude. The reduction is far less significant if we add singlet events only. Inspection of Fig. 6 leads to the coclusion that experiments in the near future with approximately 200 UHECR events can tell at least the order of magnitude of the source density.
## V SUMMARY
We presented a technique in order to statistically analyze the clustering features of UHECR. The technique can be applied for any model of UHECR assuming small deflection. The key role of the analysis is played by the $`P_k`$ functions defined by eqn. (6), which is the probability of detecting $`k`$ events above the threshold from a single source. Using a combinatorial expression of eqn. (8) the probability distribution for any set of multiplets can be given as a function of the source density.
We discussed several types of energy and luminosity distributions for the sources and gave the most probable source densities with their confidence intervals for present and future experiments.
The probability $`P(r,E,E_c)`$ that a proton created at a distance $`r`$ with energy $`E`$ arrives above the threshold $`E_c`$ is determined and parametrized for a wide range of threshold energies. This result can be used to obtain the observed energy spectrum of the UHECR for arbitrary injection spectrum.
In ref. the authors analyzed the statistical features of clustering of UHECR, which provided constraints on astrophysical models of UHECR when the number of clusters is small, by giving a bound from below. In our paper we have shown that there is some constraint, but it is far from being tight. At present statistics the 95% confidence level regions usually span 4 orders of magnitude. Two orders of magnitude smaller numbers than the prediction of (their eqn. (13) suggests for the density of sources $`610^3`$ Mpc<sup>-3</sup>) can also be obtained. Adding 10 new events with an additional doublet the confidence interval can be reduced to 3 orders of magnitude and the increase of the UHECR events to 200 can tell at least the order of magnitude of the source density.
## VI ACKNOWLEDGEMENTS
We thank K. Petrovay for clarifying some issues in spherical astronomy. This work was partially supported by Hungarian Science Foundation grants No. OTKA-T29803/T22929-FKP-0128/1997.
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# Examples, Counterexamples, and Enumeration Results for Foldings and Unfoldings between Polygons and Polytopes
## 1 Introduction
We explore the process of folding a simple polygon by gluing its perimeter shut to form a convex polyhedron, and its reverse, cutting a convex polyhedron open and flattening its surface to a simple polygon. We restrict attention to convex polyhedra (henceforth, polytopes), and to simple (i.e., nonself-intersecting, nonoverlapping) polygons (henceforth just polygons). The restriction to nonoverlapping polygons is natural, as this is important to the manufacturing applications \[O’R00\]. The restriction to convex polyhedra is made primarily to reduce the scope of the problem. See \[BDD<sup>+</sup>98\] and \[BDEK99\] for a start on unfolding nonconvex polyhedra.
Much recent work on unfolding revolves around an open problem that seems to have been first mentioned in print in \[She75\] but is probably much older: Can every polytope be cut along edges and unfolded flat to a (simple) polygon? Cutting along edges leads to edge unfoldings; we will not follow this restriction here. Thus our work is only indirectly related to this edge-unfolding question.
In some sense this report is a continuation of the investigation started in \[LO96\], which detailed an $`O(n^2)`$ algorithm for deciding when a polygon may be folded to a polytope, with the restriction that each edge of the polyon perimeter glues to another complete edge: edge-to-edge gluing. But here we do not following this restriction, permitting arbitrary perimeter gluings. Moreover, we do not consider algorithmic questions. Rather we concentrate on enumerating the number of foldings and unfoldings between polygons and polytopes. We pay special attention to convex polygons; following Shephard \[She75\], we call an unfolding of a polytope that produces a convex polygon a convex unfolding. Within the class of polytopes, we sometimes use the five regular polytopes as examples; within the class of convex polygons, we additionally focus on regular polygons.
The basic questions we ask are:
1. How many combinatorially different foldings of a polygon lead to a polytope?
2. How many geometrically different polytopes may be folded from one polygon?
3. How many combinatorially different cuttings of a polytope lead to polygon unfoldings?
4. How many geometrically different polygons may be unfolded from one polytope?
Our answers to these four questions are crudely summarized in Table 1, whose four rows correspond to the four questions above, and whose columns are for general, convex, and regular polygons. We will not explain the entries in the table here, but only remark that the increased constraints provided by convex and regular polygons reduces the number of possibilities.
A key tool in our work is a powerful theorem of Aleksandrov, which we describe and immediately apply in Section 2. We then define the two main combinatorial objects we study, cut trees and gluing trees, and make clear exactly how we count them. We then explore constraints on convex unfoldings in Section 4 before proceeding to the general enumeration bounds in Table 1 in Sections 5-8. A final section (9) concentrates on regular polygons
## 2 Aleksandrov’s Theorem
Aleksandrov proved a far-reaching generalization of Cauchy’s rigidity theorem in \[Ale58\] that gives simple conditions for any folding to a polytope. Let $`P`$ be a polygon and $`P`$ its boundary. A gluing maps $`P`$ to $`P`$ in a length-preserving manner, as follows. $`P`$ is partitioned by a finite number of distinct points into a collection of open intervals whose closure covers $`P`$. Each interval is mapped one-to-one (i.e., glued) to another interval of equal length. Corresponding endpoints of glued intervals are glued together (i.e., identified). Finally, gluing is considered transitive: if points $`a`$ and $`b`$ glue to point $`c`$, then $`a`$ glues to $`b`$.<sup>1</sup><sup>1</sup>1 What we call gluing is sometimes called pasting \[AZ67, p. 13\]. In the theory of complexes, it is sometimes called topological identification \[Hen79, p. 116\]. Aleksandrov proved that any gluing that satisfies these two conditions corresponds to a unique polytope:
1. No more than $`2\pi `$ total face angle is glued together at any point; and
2. The complex resulting from the gluing is homeomorphic to a sphere. (This condition is satisfied if, when $`P`$ is viewed as a topological circle, and the interval gluings as chords of the circle, then no pair of chords cross in the $`P`$-circle.)
Aleksandrov calls any complex (not necessarily a single polygon) that satisfies these properties a net \[Ale58\].<sup>2</sup><sup>2</sup>2 This may derive from the German translation, Netz. In fact, the Russian word Aleksandrov used is closer to “unfolding.” We call a gluing that satisfies these conditions an Aleksandrov gluing.
Although an Aleksandrov gluing of a polygon forms a unique polytope, it is an open problem to compute the three-dimensional structure of the polytope \[O’R00\]. Note that there is no specification of the fold (or “crease”) lines; and yet they are uniquely determined. Henceforth we will say a polygon folds to a polytope whenever it has an Aleksandrov gluing.
We should mention two features of Aleksandrov’s theorem. First, the polytope whose existence is guaranteed may be flat, that is, a doubly-covered convex polygon. We use the term “polytope” to include flat polyhedra. Second, condition (2) specifies a face angle $`2\pi `$. The case of equality with $`2\pi `$ leads to a point on the polytope at which there is no curvature, i.e., a nonvertex. We make explicit what counts as a vertex below.
##### Polygon/Polytope Notation.
We will use $`P`$ throughout the paper for a polygon, and $`Q`$ for a polytope. Their boundaries are $`P`$ and $`Q`$ respectively. The curvature $`\gamma (x)`$ of a point $`xQ`$ is $`2\pi `$ minus the sum of the face angles incident to $`x`$. This “angle deficit” corresponds to the notion of Gaussian curvature. We define vertices of polygons and polytopes to be essential in the sense that the boundary is not flat there: the interior angle at a polygon vertex is different from $`\pi `$, and the curvature at a polytope vertex is different from $`0`$. Because of these definitions, there is no direct correspondence between the vertices of a polytope $`Q`$ and the vertices of a polygon $`P`$ unfolding of $`Q`$: a vertex of $`Q`$ may or may not unfold to a vertex of $`P`$; and a vertex of $`P`$ may or may not fold to a vertex of $`Q`$ (see Section 3.3). At the risk of confusion, we will use the terms “vertex” and “edge” for both polygons and polytopes, but reserve “node” and “arc” for graphs. We will use $`n`$ for the number of vertices of $`P`$ or $`Q`$, letting the context determine which.
We will also freely employ two types of paths on the surface of a polytope: geodesics, which unfold (or “develop”) to straight lines, and shortest paths, geodesics which are in addition shortest paths between their endpoints. See, e.g., \[AAOS97\] for details and basic properties.
### 2.1 Perimeter Halving
As a straightforward application of Aleksandrov’s theorem, we prove that every convex polygon folds to a polytope. We will see in Section 4.2 that the converse does not hold.
For two points $`x,yP`$, define $`(x,y)`$ be the open interval of $`P`$ counterclockwise from $`x`$ to $`y`$, and let $`|x,y|`$ be its length. Define a perimeter-halving gluing as one which glues $`(x,y)`$ to $`(y,x)`$.
###### Lemma 2.1
Every convex polygon folds to a polytope via perimeter halving.
Proof: Let the perimeter of a convex polygon $`P`$ be $`L`$. Let $`xP`$ be an arbitrary point on the boundary of $`P`$, and let $`yP`$ be the midpoint of perimeter around $`P`$ measured from $`x`$, i.e., $`y`$ is the unique point satisfying $`|x,y|=|y,x|=L/2`$. See Fig. 1 for an example.
Now glue $`(x,y(`$ to $`(y,x)`$ in the natural way, mapping each point $`z`$ with $`|x,z|=d`$ to the point $`z^{}`$ the same distance from $`x`$ in the other direction: $`|z^{},x|=d`$. We claim this is an Aleksandrov gluing. It is a gluing by construction. Because $`P`$ is convex, each point along the gluing path has $`2\pi `$ angle incident to it: the gluing of two nonvertex points results in exactly $`2\pi `$, and if either point is a vertex, the total angle is strictly less than $`2\pi `$. The resulting surface is clearly homeomorphic to a sphere. By Aleksandrov’s theorem, this gluing corresponds to a unique polytope $`Q_x`$. $`\mathrm{}`$
In an Aleksandrov gluing of a polygon, a point in the interior of a polygon edge that glues only to itself, i.e., where a crease folds the edge in two, is called a fold point. A fold point corresponds to a leaf of the gluing tree, and becomes a vertex of the polytope with curvature $`\pi `$. Points $`x`$ and $`y`$ in the above proof are fold points. In Theorem 6.2 we will show that different choices of $`x`$ result in distinct polytopes $`Q_x`$, leading to the conclusion that every convex polygon folds to an infinite number of polytopes.
## 3 Cut Trees and Gluing Trees
The four main objects we study are polygons, polytopes, cut trees, and gluing trees. It will be useful in spots to distinguish between a geometric tree $`𝒯`$ composed of a union of line segments, and the more familiar combinatorial tree $`T`$ of nodes and arcs. A geometric cut tree $`𝒯_C`$ for a polytope $`Q`$ is a tree drawn on $`Q`$, with each arc a polygonal path, which leads to a polygon unfolding when the surface is cut along $`𝒯`$, i.e., flattening $`Q𝒯`$ to a plane. A geometric gluing tree $`𝒯_G`$ specifies how $`P`$ is glued to itself to fold to a polytope. There is clearly a close correspondence between $`𝒯_C`$ and $`𝒯_G`$, which are in some sense the same object, one viewed from the perspective of unfolding, one from the perspective of folding. It will nevertheless be useful to retain a distinction between them, and especially their combinatorial counterparts, which we define below after stating some basic properties.
### 3.1 Cut Trees
###### Lemma 3.1
If a polygon $`P`$ folds to a polytope $`Q`$, $`P`$ maps to a tree $`𝒯_CQ`$, the geometric cut tree, with the following properties:
1. $`𝒯_C`$ is a tree.
2. $`𝒯_C`$ spans the vertices of $`Q`$.
3. Every leaf of $`𝒯_C`$ is at a vertex of $`Q`$.
4. A point of $`𝒯_C`$ of degree $`d`$ (i.e., one with $`d`$ incident segments) corresponds to exactly $`d`$ points of $`P`$. Thus a leaf corresponds to a unique point of $`P`$.
5. Each arc of $`T_C`$ is a polygonal path on $`Q`$.
Proof:
1. If $`𝒯_C`$ contained a cycle, then it would unfold to disconnected pieces, contradicting the assumption that $`Q`$ is folded from a single polygon $`P`$. Thus $`𝒯_C`$ is a forest. But because $`𝒯_C`$ is constructed by gluing the connected path $`P`$ to itself, it must be connected. So $`𝒯_C`$ is a tree.
2. If a vertex $`v`$ of $`Q`$ is not touched by $`𝒯_C`$, then, because $`Q`$ is not flat at $`v`$, $`P`$ is not planar, a contradiction to the assumption that $`P`$ is a polygon.
3. Suppose a leaf $`x`$ of $`𝒯_C`$ is interior to a face or edge of $`Q`$. Then it is surrounded by $`2\pi `$ face angle on $`Q`$, and so unfolds to a point $`x`$ of $`P`$ similarly surrounded. But by assumption, $`x`$ is on the boundary of a simple polygon $`P`$, a contradiction.
4. Gluing exactly two distinct points of $`x,yP`$ together implies that neighborhoods of $`x`$ and $`y`$ are glued, which leads to the interior of an arc of the cut tree, i.e., a degree-$`2`$ point of $`𝒯_C`$. Note that either or both of these points might be vertices of $`P`$. In general, if $`p𝒯_C`$ has $`d`$ incident cut segments, $`p`$ unfolds to $`d`$ distinct points of $`P`$.
5. If an arc of $`𝒯_C`$ is not a polygonal path, then neither side unfolds to a polygonal path, contradicting the assumption that $`P`$ is a polygon.
$`\mathrm{}`$
When counting cut trees, we will rely on their combinatorial structure. There are several natural definitions of this structure, which are useful in different circumstances. We first discuss some of the options.
1. Make every segment of $`𝒯_C`$ an arc of $`T_C`$. Although this is very natural, it means there are an infinite number of different cut trees for any polytope, for the path between any two polytope vertices could be an arbitrarily complicated polygonal path, leading to different combinatorial trees.
2. Make every point where a path of $`𝒯_C`$ crosses an edge of the polytope a node of $`T_C`$. This again leads to trivially infinite numbers of cut trees when a path of $`𝒯_C`$ zigzags back and forth over an edge of $`Q`$.
3. Exclude this possibility by forcing the paths between polytope vertices to be geodesics, and again make polytope edge crossings nodes of $`T_C`$. This excludes many interesting cut trees—all those where a polygon vertex is glued to a point with angle sum $`2\pi `$.
4. Make every maximal path of $`𝒯_C`$ consisting only of degree-$`2`$ points a single arc of $`T_C`$. This has the undesirable effect of having polytope vertices in the interior of such a path disappear from $`T_C`$.
Threading between these possibilities, we define the combinatorial cut tree $`T_C`$ corresponding to a geometric cut tree $`𝒯_C`$ as the labeled graph with a node (not necessarily labeled) for each point of $`𝒯_C`$ with degree not equal to $`2`$, and a labeled node for each point of $`𝒯_C`$ that corresponds to a vertex of $`Q`$ (labeled by the vertex label); arcs are determined by the polygonal paths of $`𝒯_C`$ connecting these nodes. An example is shown in Fig. 2. Note that not every node of the tree is labeled, but every polytope vertex label is used at some node. All degree-$`2`$ nodes are labeled.
Although this definition avoids some of the listed pitfalls, it does have the undesirable consequence of counting different geodesics on $`Q`$ between two polytope vertices as the same arc of $`T_C`$. Thus the two unfoldings shown in Fig. 5 (below) have the same combinatorial cut tree under our definition, even though the geodesic in (c) spirals twice around compared to once around in (a).
### 3.2 Gluing Trees
Let a convex polygon $`P`$ have vertices $`v_1,\mathrm{},v_n`$, labeled counterclockwise, and edge $`e_i`$, $`i=1,\mathrm{},n`$ the open segment of $`P`$ after $`v_i`$. There is less need to discuss the geometric gluing tree, so we concentrate on the combinatorial gluing tree $`T_G`$. $`T_G`$ is a tree representing the identification of $`P`$ with itself. Any point of $`P`$ that is identified with more or less than one other distinct point of $`P`$ becomes a node of $`T_G`$, as well as any point to which a vertex is glued. (Note that this means there may be nodes of degree $`2`$.) So every vertex of $`P`$ maps to a node of $`T_G`$; each node is labeled with the set of all the elements (vertices or edges) that are glued together there. A leaf that is a fold point is labeled by the edge label only. Every nonleaf node has at least one vertex label, and at most one edge label. A simple example is shown in Fig. 3.<sup>3</sup><sup>3</sup>3 Gluing trees can be drawn by folding up the polygon toward the viewer (as in this figure), or folding the polygon away. We employ both conventions but always note which is followed. Here the central node of $`T_G`$ is assigned the label $`\{v_1,v_3,e_3\}`$.
A more complicated example is shown in Fig. 4.<sup>4</sup><sup>4</sup>4 We found this example by an enumeration algorithm that will not be discussed in this report. The polygon shown folds (amazingly!) to a tetrahedron by creasing as illustrated in (a). All four tetrahedron vertices are fold points. The corresponding gluing tree is shown in (b) of the figure. The two interior nodes of $`T_G`$ have labels $`\{v_1,v_6,e_1\}`$ and $`\{v_2,v_5,e_5\}`$.
Later (Lemma 5.3) will show that the gluing tree is determined by a relatively sparse set of gluing instructions.
### 3.3 Comparison of Cut and Gluing Trees
###### Lemma 3.2
Let $`T_C`$ be a combinatorial cut tree for polytope $`Q`$ that unfolds to a polygon $`P`$, and let $`T_G`$ be the combinatorial gluing tree that folds $`P`$ to $`Q`$. If all degree-$`2`$ nodes are removed by contraction, $`T_C`$ and $`T_G`$ are isomorphic as unlabeled graphs.
Proof: Let $`(a,b,c)`$ be three consecutive nodes on a path in a tree $`T`$, with $`b`$ of degree $`2`$. Removing $`b`$ by contraction deletes $`b`$ and replaces it with the arc $`(a,c)`$. Applying this to both $`T_C`$ and $`T_G`$ produces two trees $`T_C^{}`$ and $`T_G^{}`$ without degree-$`2`$ nodes. As the trees were defined to include nodes for each point whose degree differs from $`2`$, it must be that $`T_C^{}`$ and $`T_G^{}`$ have isomorphic structures. Of course they are labeled differently, but without the labels, they are isomorphic graphs. $`\mathrm{}`$
Note that vertices in $`Q`$ and vertices in $`P`$ do not necessarily map to one another: A vertex of $`Q`$ can map to an interior point of $`P`$, and a vertex of $`P`$ can map to a point interior to a face or edge of $`Q`$. This affects the labeling of the two trees, but they have essentially the same structure.
## 4 Cut Trees for Convex Unfoldings
Before embarking on general enumeration results, we specialize the discussion to convex unfoldings, and derive some constraints on the possible cut trees that lead to convex unfoldings.
### 4.1 Stronger Characterization
We now sharpen the characterization of cut trees (and via Lemma 3.2, of gluing trees) under the restriction that the unfolding must be a convex polygon. We first strengthen Lemma 3.1(5), which only required arcs to be polygonal paths:
###### Lemma 4.1
Every arc of a cut tree $`T_C`$ that leads to a convex unfolding must be a geodesic on $`Q`$ (paths that unfold to straight segments), but arcs might not be shortest paths on $`Q`$.
Proof: Suppose an arc $`a`$ of $`T`$ is not a geodesic. Then it does not unfold to a straight line. Suppose a point $`xa`$ is a point in the relative interior of $`a`$ at which the unfolding is locally not straight. Then only one of the two points of $`P`$ that correspond to $`x`$ can have an interior angle $`\pi `$ in $`P`$, showing that $`P`$ has at least one reflex angle. This establishes that arcs of $`T_C`$ must be geodesics. We now show that this claim cannot be strengthened to shortest paths by an explicit example.
Let $`Q`$ be a doubly-covered rectangle with vertices $`v_i`$, $`i=1,2,3,4`$, as shown in Fig. 5(a). Let $`x`$ be the midpoint of edge $`v_1v_4`$. Let $`T_C`$ be the path $`(v_1,v_2,x,v_3,v_4)`$, where the subpath $`(v_2,x,v_3)`$ is half on the upper rectangular face, and half on the bottom face. Clearly this subpath is not a shortest path, although it is a geodesic. The corresponding convex unfolding is shown in Fig. 5(b).
This example can be modified to a nondegenerate “sliver” tetrahedron by perturbing one vertex to lie slightly out of the plane of the other three. $`\mathrm{}`$
Fig. 5(c-d) shows that we cannot even bound the length of a geodesic arc of $`T_C`$.
One immediate corollary of Lemma 4.1 is that cuts need not follow polytope edges (which are all shortest paths), i.e., not every convex unfolding is an edge unfolding.
### 4.2 Necessary Conditions: Sharp Vertices
We define a vertex of a polytope to be sharp if it has curvature $`\pi `$, and round if its curvature is $`<\pi `$. The following theorem gives a simple necessary condition for a polytope to have a convex unfolding. We employ this fact implied by the Gauss-Bonnet theorem:
###### Fact 4.1
The sum of the curvatures of all the vertices of a polytope is exactly $`4\pi `$.
###### Theorem 4.2
If a polytope $`Q`$ has a convex unfolding via a cut tree $`T_C`$, then each leaf of $`T_C`$ is at a sharp vertex. Moreover, $`Q`$ must have at least two sharp vertices.
Proof: Let $`P`$ be a convex polygon to which $`Q`$ unfolds via cut tree $`T_C`$. By Lemma 3.1(3), the leaves of $`T_C`$ are at vertices of $`Q`$. Let $`x`$ be a leaf of $`T_C`$, at a vertex $`v`$ with curvature $`\gamma (v)=\gamma `$. Point $`xQ`$ corresponds to a unique point $`yP`$ by Lemma 3.1(4). The internal angle at $`y`$ in $`P`$ is $`2\pi \gamma `$. Because $`P`$ is convex, we must have
$$2\pi \gamma \pi $$
and so $`\gamma (v)\pi `$. Thus $`v`$ is sharp. Because $`T_C`$ must have at least two distinct leaves, the lemma follows. $`\mathrm{}`$
###### Corollary 4.3
Of the five Platonic solids, only the regular tetrahedron has a convex unfolding.
Proof: The curvatures at the vertices of the solids are:
$`2\pi 3(\pi /3)=`$ $`\pi `$
$`2\pi 3(\pi /2)=`$ $`\pi /2`$ $`<\pi `$
$`2\pi 4(\pi /3)=`$ $`2\pi /3`$ $`<\pi `$
$`2\pi 3(3\pi /5)=`$ $`\pi /5`$ $`<\pi `$
$`2\pi 5(\pi /3)=`$ $`\pi /3`$ $`<\pi `$
Only the tetrahedron has sharp vertices. $`\mathrm{}`$
We next show that two natural extensions of the previous results fail.
###### Lemma 4.4
There is a tetrahedron with no convex unfolding.
Proof: Let $`Q_1`$ be a tetrahedron whose vertices $`v_1,v_2,v_3`$ form an equilateral triangle base in the $`xy`$-plane, with apex $`v_4`$ centered at a great height $`z`$ above. See Fig. 7. Let $`\gamma _i`$ be the curvature of vertex $`v_i`$. If the face angle of each triangle incident to $`v_4`$ is $`ϵ`$, then $`\gamma _4=2\pi 3ϵ`$, and $`\gamma _i`$ for $`i=1,2,3`$ is
$$2\pi [\pi /3+2(\pi ϵ/2)]=2\pi /3+ϵ$$
Choosing $`z`$ large makes $`ϵ`$ small, and then $`Q_1`$ has just one sharp vertex. Theorem 4.2 then establishes the claim. $`\mathrm{}`$
###### Lemma 4.5
There is a polytope with two sharp vertices but with no convex unfolding.
Proof: Our proof of this lemma is less straightforward, although the example is simple. Let $`Q_2`$ be the polytope formed by joining two copies of $`Q_1`$ from Lemma 4.4 at their bases, as shown in Fig. 7. $`Q_2`$ is a $`5`$-vertex polytope, with vertices $`v_1,\mathrm{},v_4`$ as in $`Q_1`$, and $`v_5`$ the reflection of $`v_4`$ in the central triangle $`C=\mathrm{}v_1v_2v_3`$. Again let $`ϵ`$ be the face angle incident to $`v_4`$ (and symmetrically $`v_5`$), and choose $`ϵ`$ small so that only $`v_4`$ and $`v_5`$ are sharp vertices.
By Lemma 3.1(3), if $`Q_2`$ has a convex unfolding, the cut tree must be a path with its two leaves at the two sharp vertices. By Lemma 3.1(5), the path must be composed of geodesics. We now analyze the geodesics starting at $`v_5`$ and show that there can be no piecewise simple geodesic path that passes through all the vertices of $`Q_2`$.
We group the geodesics starting at $`v_5`$ into three classes:
1. The three geodesics that pass through a midpoint of an edge of triangle $`C`$. Each of these passes through $`v_4`$ before encountering any of the other vertices, and so cannot serve as the cut path.
2. The three geodesics that pass through a vertex of $`C`$. Because these vertices have low curvature ($`2ϵ`$), the geodesic must emerge nearly headed toward $`v_4`$: it cannot turn to hit another vertex of $`C`$ without creating a reflex angle in the unfolding. If the geodesic goes directly to $`v_4`$, then again this cannot serve as the cut path. So it must head towards $`v_4`$ but miss it. We group this type of geodesic with the third class.
3. Geodesics that pass though an interior point of an edge of $`C`$, but not the midpoint. These geodesics all head toward $`v_4`$ but miss it.
We now argue that all the geodesics in the third class (the only remaining candidates) self-intersect after looping around $`v_4`$. This will then establish the lemma.
An unfolding of a typical geodesic is shown in Fig. 8. By choosing $`ϵ`$ small, we can arrange that every such geodesic crosses several unfoldings of the three faces incident to $`v_4`$ before returning back down to triangle $`C`$. As can be seen from the copy of face $`\mathrm{}v_1v_2v_4`$ to the side, the path crosses each face several times slanting one way, and then returns slanting the other way. In the vicinity of the closest approach to $`v_4`$, the path must self-cross. We now establish this more formally.
Consider the unfolding of the three faces incident to $`v_4`$ (now viewed as a unit) that includes the point $`p`$ of closest approach between the geodesic and $`v_4`$; see Fig. 9. Let the geodesic cross the edge $`v_1v_4`$ at points $`a`$, $`x`$, $`y`$, and $`b`$ in that order, with $`xy`$ including $`p`$. Then $`|v_4b|>|v_4x|`$ and $`|v_4a|>|v_4y|`$, because the distance from $`v_4`$ monotonically increases on either side of $`p`$. Thus the images $`a^{}`$ and $`b^{}`$ of $`a`$ and $`b`$ must fall below $`y`$ and $`x`$ respectively in the figure. Thus the geodesic must cross somewhere in the stretch immediately before and after the closest approach.
All we need for this argument to hold in general is for the geodesic to cross three complete unfoldings of the three faces incident to $`v_4`$ before returning to the lower half of $`Q_2`$. But this is easily arranged by choosing $`ϵ`$ small.
We have shown that no geodesic starting from $`v_5`$ may serve as a cut path for a convex unfolding. Therefore $`Q_2`$ has no convex unfolding. $`\mathrm{}`$
### 4.3 Necessary Conditions: Combinatorial Structure
We now study the combinatorial structure of cut trees that lead to convex unfoldings. The following theorem is due to Shephard \[She75\], although under different assumptions and with a different proof.<sup>5</sup><sup>5</sup>5 Shephard concludes that cut trees cannot have four leaves, an incorrect claim under our assumptions.
###### Theorem 4.6
If a polytope $`Q`$ of $`n4`$ vertices has a convex unfolding, then the corresponding cut tree $`T_C`$ has two or three leaves: it is either a path, or a ‘Y’ (a single degree-$`3`$ node). If $`n=4`$, then additionally it may have four leaves, and have the combinatorial structure of ‘+’ (a single degree-$`4`$ node), or two degree-$`3`$ nodes connected by an edge, which we will call a ‘I’.
Proof: Let the cut tree $`T_C`$ unfold $`Q`$ to a convex polygon. By Theorem 4.2, each leaf of $`T_C`$ must be at a sharp vertex $`v`$, and so have curvature $`\gamma (v)\pi `$. If $`T_C`$ has more than four leaves $`v`$ (and therefore $`n>4`$, i.e., we are in the $`n4`$ case of the theorem claim), $`_v\gamma (v)>4\pi `$, which violates the Gauss-Bonnet theorem. Therefore $`T_C`$ has no more than four leaves. If $`T_C`$ has just two or three leaves, then the only possible combinatorial structures for $`T_C`$ are the two claimed in the theorem: a path, and a ‘Y’. (Note that it is possible that $`n=3`$, when $`Q`$ is a doubly-covered triangle.)
So assume that $`T_C`$ has exactly four leaves. Because each leaf vertex is sharp, $`_v\gamma (v)4\pi `$; on the other hand, we know the sum over all vertices is equal $`4\pi `$. Therefore we know that each leaf has curvature exactly $`\pi `$ and that the leaves of $`T_C`$ are at the only vertices of $`Q`$. Thus $`n=4`$ and $`Q`$ is a tetrahedron. The only additional possible combinatorial structures for a tree with four leaves are the two claimed in the theorem: a ‘+’ and a ‘I’. Note that in both these cases, the internal node(s) of $`T_C`$ are not at vertices of $`Q`$. $`\mathrm{}`$
A simple example of the ‘I’ possibility is shown in Fig. 10. If the rectangle is modified to become a square, the ‘I’ becomes a ‘+’.
## 5 Counting Foldings: Gluing Trees
In this section we move beyond Lemma 2.1, which shows that every convex polygon folds to a polytope, and explore how many different ways there are to fold a given polygon, as measured by the number of combinatorially distinct Aleksandrov gluing trees. In Section 6 we count instead the number of distinct polytopes that might be produced from a given polygon. In both cases, we will also examine the restriction to convex polygons, which not surprisingly yields sharper results.
### 5.1 Unfoldable Polygons
We start with a natural and easily proved claim:
###### Lemma 5.1
Some polygons cannot be folded to any polytope.
Proof: Consider the polygon $`P`$ shown in Fig. 11.
$`P`$ has three consecutive reflex vertices $`(a,b,c)`$, with the exterior angle $`\beta `$ at $`b`$ small. All other vertices are convex, with interior angles strictly larger than $`\beta `$.
Either the gluing “zips” at $`b`$, leaving $`b`$ a leaf of $`T_G`$, or some other point(s) of $`P`$ glue to $`b`$. The first possibility forces $`a`$ to glue to $`c`$, exceeding $`2\pi `$ there; so this gluing is not Aleksandrov. The second possibility cannot occur with $`P`$, because no point of $`P`$ has small enough internal angle to fit at $`b`$. Thus there is no Aleksandrov gluing of $`P`$. $`\mathrm{}`$
It is natural to wonder what the chances are that a random polygon could fold to a polytope. This is difficult to answer without a precise definition of “random,” but we feel any reasonable definition would lead to the same answer:
###### Conjecture 5.1
The probability that a random polygon of $`n`$ vertices can fold to a polytope approaches $`0`$ as $`n\mathrm{}`$.
Proof: (Sketch.) Assume that random polygons on $`n`$ vertices satisfy two properties:
1. The distribution of the polygon angles approaches the uniform distribution on the interval $`(0,2\pi )`$ as $`n\mathrm{}`$. In particular, the number of reflex and convex vertices approaches balance.
2. The distribution of polygon edge lengths approaches some continuous density distribution.
For large $`n`$, we expect $`P`$ to have $`r=n/2`$ reflex vertices. Each of these reflex vertices $`a`$ faces one of two fates in the gluing tree: either it becomes a leaf by “zipping” at $`a`$; or at least one convex vertex $`b`$ (of sufficiently small angle) is glued to $`a`$. The number of reflex vertices that can be zipped is limited by Fact 4.1: if $`a`$ has angle $`\alpha `$, zipping there adds $`2\pi \alpha `$ to the curvature; but the total curvature is limited to $`4\pi `$. Suppose we zip the largest $`k`$ angles out of the $`r`$ reflex vertices (the largest angles increment the curvature the least). Then one can compute that, under the uniform angle distribution assumption, these $`k`$ angles have an expected curvature sum of
$$\frac{1}{2}\frac{\pi }{r}k^2.$$
(1)
(For example, for $`r=100`$, the largest $`k=10`$ have an expected curvature sum of $`\pi /2`$.) Limiting this to $`4\pi `$ implies that the expected maximum number of reflex vertices that can be zipped without exceeding $`4\pi `$ curvature is
$$k2\sqrt{2r}=2\sqrt{n}.$$
(2)
(For example, for $`r=1000`$ reflex vertices, the largest $`k=89`$ lead to a curvature of $`4\pi `$.) Thus, at most a small portion of the reflex vertices can be zipped; the remainder (expected number: $`n/22\sqrt{n}`$) must be glued to convex vertices. We now show that this gluing is not in general possible.
Let $`a`$ be a reflex vertex with angle $`\alpha `$, and $`b`$ a convex vertex whose angle $`\beta `$ satisfies $`\beta 2\pi \alpha `$, so that $`b`$ can glue to $`a`$. It could be that this gluing forces one or more reflex vertices adjacent to $`a`$ or $`b`$ to glue to edges incident to $`a`$ or $`b`$, in which case the gluing is not possible (i.e., it is not an Aleksandrov gluing). If the adjacent vertices are convex, and/or the edge lengths are such that the gluing is Aleksandrov, then, in general, two new reflex vertices are created, as is illustrated in Fig. 12.
To be more precise, let $`A=|aa_1|`$ be the length of the edge incident to $`a`$ which is glued to the length $`B=|bb_1|`$ of an edge incident to $`b`$. If $`A>B`$ and $`b_1`$ is reflex, the gluing is not Aleksandrov; but if $`b_1`$ is convex, a new reflex vertex is created at $`b_1`$. Symmetrically, if $`B>A`$ and $`a_1`$ is reflex, the gluing is not possible; but if $`a_1`$ is convex, a new reflex vertex is created at $`a_1`$. The only circumstance in which the gluing is Aleksandrov and a new reflex vertex is not created is when $`A=B`$ and both $`a_1`$ and $`b_1`$ are convex with an angle sum of no more than $`\pi `$.
Under the assumption that the edge lengths approach some continuous distribution, the probability that two lengths match exactly approaches $`0`$. Thus we conclude that gluing convex vertices to reflex vertices does not remove reflex vertices, but rather creates new ones in shorter polygonal chains, one new reflex vertex in each of the two chains produced by the gluing. Note that gluing several convex vertices to one reflex vertex does not change matters: we can view the first convex vertex as simply leaving a reflex remainder, and argue as above.
Thus, any gluing of a random polygon for large $`n`$ will lead to shorter and and shorter chains “pinched” between reflex-convex gluings, each of which will contain at least one reflex vertex (actually, two reflex vertices for those pinched on both sides). Eventually these chains reach the point where either there are no convex vertices that fit into the reflex vertex gap, or there are no convex vertices at all. In either case, the chain cannot be glued: the reflex vertex would have to glue to a point interior to an edge, violating the Aleksandrov condition that no point have more than $`2\pi `$ glued angle. $`\mathrm{}`$
The proof above hinges on the unlikeliness of matching edge lengths. It is therefore natural to wonder if the same result holds for polygons all of whose edge lengths are the same. Again we believe it does:
###### Conjecture 5.2
The probability that a random polygon of $`n`$ vertices, all of whose edges have unit length, can fold to a polytope, approaches $`0`$ as $`n\mathrm{}`$.
Proof: (Sketch.) Assume a model of random polygons such that the angles are probabilistically independent and uniformly distributed in $`(0,2\pi )`$ as $`n\mathrm{}`$. The restriction to unit edge lengths means that all gluings are vertex to vertex (no vertex is ever glued to the interior of an edge). The gluing is Aleksandrov iff the angles glued together sum to at most $`2\pi `$ everywhere.
Consider gluing two vertices to one another. Because their angles are independent, the chance that the gluing is legal is $`1/2`$ (the sum of their distributions is uniform between $`0`$ and $`4\pi `$). Gluing $`k`$ pairs then has a $`1/2^k`$ chance of being Aleksandrov.
As in the above proof sketch, the gluing tree cannot have too many leaves. Zipping just $`2\sqrt{n}`$ reflex vertices uses up all $`4\pi `$ of curvature. So the number of leaves is only about $`2\sqrt{n}`$. As we will see in Theorem 5.11 below, specifying a “source” for each leaf pins down the whole tree structure. So by selecting $`4\sqrt{n}`$ vertices for the leaves and their sources, the gluing tree is determined.
Therefore we should compare the number of different gluing trees,
$$\left(\begin{array}{c}n\\ 4\sqrt{n}\end{array}\right)$$
(3)
to the probability that each one is Aleksandrov,
$$\frac{1}{2^{n4\sqrt{n}}}$$
(4)
Note here we conservatively only concern ourselves with degree-$`2`$ vertex-to-vertex gluings; junctions of degree $`d>2`$ have a lower probability of summing to no more than $`2\pi `$. We also ignore the change to the angle distribution caused by the removal of the leaf vertices.
Using Stirling’s approximation shows that the $`\mathrm{log}`$ of Eq. (3) grows as $`2\sqrt{n}\mathrm{log}n`$; but the $`\mathrm{log}`$ of Eq. (4) grows as $`n`$. So their ratio approaches $`0`$ as $`n\mathrm{}`$. $`\mathrm{}`$
We leave these results on random polygons as conjectures, as it would require a more precise definition of what constitutes a random polygon, and more careful probabilistic analyses, to establish them formally.
### 5.2 Lower Bound: Exponential Number of Gluing Trees
In contrast to the likely paucity of foldable polygons, some polygons generate many foldings.
###### Theorem 5.2
For any even $`n`$, there is a polygon $`P`$ of $`n`$ vertices that has $`2^{\mathrm{\Omega }(n)}`$ combinatorially distinct Aleksandrov gluings.
Proof: The polygon $`P`$ is illustrated in Fig. 13(a).
It is a centrally symmetric star, with $`m`$ vertices, $`m`$ even, with a small convex angle $`\alpha 0`$, alternating with $`m`$ vertices with large reflex angle $`\beta <2\pi `$. All edges have the same (say, unit) length. We call this an $`m`$-star. We first specify the constraints on $`\alpha `$ and $`\beta `$.
$`P`$ has $`n^{}=2m`$ vertices (ignoring $`x`$ and $`y`$, to be described shortly). So $`m(\alpha +\beta )=(n^{}2)\pi `$, which implies that
$$\alpha +\beta =(1\frac{1}{m})2\pi .$$
(5)
We choose $`\alpha `$ small enough so that $`m`$ copies of $`\alpha `$ can join with one of $`\beta `$ and still be less than $`2\pi `$:
$$m\alpha +\beta <2\pi .$$
(6)
Substituting this relationship into Eq. (5) and solving for $`\alpha `$ yields:
$$\alpha <\frac{2\pi }{m(m1)}.$$
(7)
Now we add two vertices $`x`$ and $`y`$ at the midpoints of edges, symmetrically placed so that $`y`$ is half the perimeter around $`P`$ from $`x`$. Let $`n=n^{}+2`$ be the total number of vertices of $`P`$.
The “base” gluing tree is illustrated in Fig. 13(b). $`x`$ and $`y`$ are fold vertices of the gluing. Otherwise, each $`\alpha `$ is matched with a $`\beta `$. Because all edge lengths are the same, and because $`\alpha +\beta <2\pi `$ by Eq. (5), this path is an Aleksandrov gluing. We label it $`T_{00\mathrm{}0,00\mathrm{}0}`$, where $`m/2`$ zeros $`00\mathrm{}0`$ represent the top chain, and another $`m/2`$ zeros represent the bottom chain.
The other gluing trees are obtained via “contractions” of the base tree. A contraction makes any particular $`\beta `$-vertex not adjacent to $`x`$ or $`y`$ a leaf of the tree by gluing its two adjacent $`\alpha `$-vertices together. Label a $`\beta `$-vertex $`0`$ or $`1`$ depending on whether it is uncontracted or contracted respectively. Then a series of contractions can be identified with a binary string. For example, Fig. 13(c) displays the tree $`T_{010100\mathrm{},00110\mathrm{}0}`$. Note that $`k`$ adjacent contractions result in $`2k`$ $`\alpha `$-vertices glued together.
We now claim that if the number of contractions in the top chain is the same as the number in the bottom chain (call such a series of contractions balanced), the resulting tree represents an Aleksandrov gluing. Fix the position of $`x`$ to the left, and contract leftwards, as in Fig. 13(c). Then it is evident that the alternating “parity” pattern of $`\alpha `$’s and $`\beta `$’s is not changed by contractions. Ignoring the arcs attached to the central path, each contraction replaces $`\alpha 2\alpha `$, and shortens the path by $`2`$ units. Because the contraction shortens by an even number of units, it does not affect the parity pattern. If the top and bottom chains are contracted the same number of times (twice each in (c) of the figure), then their lengths are the same.
Thus after a balanced series of contractions, we have a number of $`\beta `$-leaves, and gluings of $`2k`$ $`\alpha `$-vertices to one $`\beta `$-vertex. The $`\beta `$-leaves are legal gluings because $`\beta <2\pi `$. Because there are $`m/21`$ contractible $`\beta `$-vertices in each chain, the longest series of adjacent contractions is $`m/21`$. So $`km/21`$, and $`2k<m`$. Eq. (6) then shows that each gluing produces less than $`2\pi `$ angle, and so is Aleksandrov.
Finally, we bound the number of gluings. There are $`2^{m/21}`$ binary numbers of $`m/21`$ bits. Thus there are this many ways to contract the top chain. The bottom chain must be contracted with the same number of $`1`$’s for a balanced series. Rather than count this explicitly, we simply note that $`P`$ has at least $`2^{m/21}`$ Aleksandrov gluings, and because $`P`$ has $`n=n^{}+2=2m+2`$ vertices, $`\mathrm{\Omega }(2^{m/21})=\mathrm{\Omega }(2^{(n6)/4})=2^{\mathrm{\Omega }(n)}`$. $`\mathrm{}`$
Fig. 14 shows six gluings of a $`4`$-star. The first two in the top row correspond to the perimeter-halving construction used in the proof. By Aleksandrov’s theorem, each corresponds to a unique polytope, but as mentioned in Section 2, we do not know how to compute the 3D structure of these polytopes. Nevertheless, our hand-exploration suggest that all fold to noncongruent polytopes, each with the combinatorial structure of the regular octahedron. Two of our conjectured crease patterns are shown in Fig. 15.
### 5.3 Upper Bound: Few Leaves
Our goal is now to provide upper bounds on the number of gluings, both for arbitary polygons and for convex polygons. Both will rely on upper bounds for gluing trees with a small number of leaves. Let a gluing tree $`T_G`$ have $`\lambda `$ leaves. In this section, we prove results for $`\lambda =2`$ and $`\lambda =3`$. We then use these to obtain a general upper bound in Section 5.4, and a bound for convex polygons in Section 5.6. In between, we summarize the structural properties of gluing trees in Section 5.5.
It will sometimes be easier to work with “gluing instructions” rather than with gluing trees. Toward that end, we define the combinatorial type of a gluing. Again let polgyon $`P`$ have vertices $`v_i`$ and edges $`e_i`$, labeled counterclockwise, The combinatorial type $`\mathrm{\Gamma }_G`$ of a gluing $`G`$ specifies to which vertex or edge of $`P`$ each vertex of $`P`$ glues via a set of ordered pairs: $`\mathrm{\Gamma }_G=\{(v_i,z_j)\}`$, where $`z_j`$ is either $`v_j`$ or $`e_j`$, the first element $`j`$ to which $`v_i`$ glues counterclockwise around $`P`$. If $`v_i`$ is a leaf of the cut tree, then the pair $`(v_i,v_i)`$ is included; otherwise $`v_i`$ must glue to an element different from itself. For example, the combinatorial type of the gluing illustrated earlier in Fig. 3a is
$$\{(v_1,v_3),(v_2,v_2),(v_3,e_3)\}$$
We now prove that the combinatorial type of a gluing determines the gluing tree.
###### Lemma 5.3
The combinatorial type $`\mathrm{\Gamma }_G`$ of a gluing $`G`$ determines the gluing tree $`T_G`$.
Proof: A node of degree $`2`$ of $`T_G`$ is directly labeled in $`\mathrm{\Gamma }_G`$ as either $`(v_i,v_j)`$ or $`(v_i,e_j)`$. It is only nodes of degree $`2`$ for which $`T_G`$ contains information not immediately supplied by $`\mathrm{\Gamma }_G`$. Nodes of degree $`1`$ (leaves) of $`T_G`$ correspond to two possible types of gluings: either $`(v_i,v_i)`$, which are directly labeled in $`\mathrm{\Gamma }_G`$, or fold vertices, a vertex produced by folding at a point $`x`$ in the interior of an edge $`e_j`$. (Cf. Fig. 1 for an example of fold vertices.) Fold vertices can be identified in $`\mathrm{\Gamma }_G`$ as gluings of $`v_i`$ to either $`e_i`$ or $`e_{i1}`$: gluing to an incident edge necessarily implies a fold vertex on that edge. Or $`v_i`$ can be glued to the next vertex, folding the edge in half. In Fig. 3, the pair $`(v_3,e_3)`$ identifies fold vertex $`x`$ as labeled with $`e_3`$; that $`v_1`$ also glues to incident edge $`e_3`$ is known after the degree $`3`$ node’s labels are determined.
Nodes of degree $`d>2`$ in $`T_G`$ have $`d`$ labels. Because every such node can involve at most one edge (because two edges glued to a point already gives an angle of $`2\pi `$ there, and the other elements glued to the same point would cause the angle sum to exceed this), the labels can be gathered by following the gluings counterclockwise:
$$(v_{i_1},v_{i_2}),(v_{i_2},v_{i_3}),\mathrm{},(v_{i_{d2}},v_{i_{d1}}),(v_{i_{d1}},e_j).$$
In Fig. 3, the node at point $`z`$ has labels $`\{v_1,v_3,e_3\}`$, which can be identified from the pairs $`(v_1,v_3),(v_3,e_3)`$ of $`\mathrm{\Gamma }_G`$. $`\mathrm{}`$
This lemma permits us to count gluing trees by counting combinatorial types of gluings.
###### Lemma 5.4
A polygon $`P`$ of $`n`$ vertices has $`\mathrm{\Theta }(n^2)`$ different gluing trees of two leaves, i.e., paths.
Proof: View $`P`$ as rolling continuously between the two leaves $`x`$ and $`y`$, like a conveyor belt or tank tread. Each specific position corresponds to a perimeter-halving gluing $`G`$ (Fig. 1). The combinatorial type $`\mathrm{\Gamma }_G`$ changes each time a vertex $`v_i`$ either passes another vertex $`v_j`$, or becomes the leaf $`x`$ or $`y`$. Each such event corresponds to two distinct types: the type at the event, and the type just beyond it: e.g., $`(v_i,v_j)`$ and $`(v_i,e_j)`$. So counting events undercounts by half. If we count the possible pairs $`(v_i,v_j)`$ for all $`ij`$, we will double count each type: the event $`(v_i,v_j)`$ leads to the same type as $`(v_j,v_i)`$. The undercount by half and overcount by double cancel; thus $`n(n1)`$ is the number of types without a vertex at a leaf. Adding in the $`n`$ possible $`(v_i,v_i)`$ events, each of which leads to two types, yields an upper bound of $`n(n1)+2n=O(n^2)`$ on the number of combinatorial types.
A lower bound of $`\mathrm{\Omega }(n^2)`$ is achieved by the example illustrated in Fig. 16(a). Here $`n/2`$ vertices of $`P`$ are closely spaced within a length $`L`$ of $`P`$, and $`n/2`$ vertices are spread out by more than $`L`$ between each adjacent pair. Then each of the latter vertices (on the lower belt in the figure) can be placed between each pair of the former vertices (on the upper belt), yielding $`n^2/4`$ distinct types. This example can be realized geometrically by making the internal angle at each vertex nearly $`\pi `$, i.e., by a convex polyon that approximates a circle.
Lemma 5.3 shows that the bound just obtained of $`\mathrm{\Theta }(n^2)`$ on the number of combinatorial types applies as well to the number of gluing trees. $`\mathrm{}`$
###### Lemma 5.5
A polygon $`P`$ of $`n`$ vertices folds to at most $`O(n^4)`$ different gluing trees of three leaves, i.e., ‘Y’s.
Proof: Observe that the degree-$`3`$ node of the ‘Y’ is either comprised by the gluing two vertices and an edge together (call this type-vve), or three vertices (type-vvv). It is not possible to glue two or more edges together without violating the $`2\pi `$ angle restriction of an Aleksandrov gluing.
There are $`O(n^3)`$ possible type-vvv nodes for the ‘Y’. Once this type of node is specified, the entire gluing tree is determined, so this bounds the number of ‘Y’s with type-vvv nodes. Now consider type-vve nodes. There are $`O(n^2)`$ possible vv-gluings, which determine one branch of the ‘Y’. The remainder of the ‘Y’ can be viewed as a path between its leaves; essentially this view corresponds to a conveyor belt with an appendage. Applying Lemma 5.4 yields a bound of $`O(n^4)`$.
$`\mathrm{}`$
We leave open the question of whether this bound is tight. We will improve it for convex polygons in Section 5.6.
#### 5.3.1 Four Fold-Point Gluing Trees
We now embark on a study of a special case that will play two roles: in the proof of our main combinatorial upper bound, Theorem 5.11, and in counting noncongruent polytopes in Section 6. Define a four fold-point gluing tree to be a gluing tree with (at least) four leaves, each fold points, i.e., creases in the interior of polygon edges leading to polytope vertices of curvature $`\pi `$. We have already encountered one such tree in Fig. 4(b). We start with this straightforward lemma.
###### Lemma 5.6
A four fold-point gluing tree must have exactly four leaves, and so have combinatorial structure ‘+’ or ‘I’.
Proof: Because each fold point leads to a vertex of the resulting polytope $`Q`$ which has curvature $`\pi `$, Fact 4.1 implies that all the curvature of the polytope is at the four fold vertices. Thus all vertices of $`P`$ must glue to points that have total angle $`2\pi `$, so that the curvature there is zero.
A leaf of a gluing tree cannot have zero curvature. This is because a leaf is either a fold point (curvature $`\pi `$) or a “zipped” polygon vertex $`v`$. The only way to achieve zero curvature at a zipped vertex is to have an internal polygon angle at $`v`$ of $`2\pi `$. But this violates simplicity of $`P`$: all internal angles are strictly less than $`2\pi `$.
Therefore, a four-fold gluing tree must have exactly four leaves. So there are only two possible combinatorial structures: ‘+’ and ‘I’ (as in Lemma 3.1). $`\mathrm{}`$
Before counting the number of gluing trees, we detail one example that will be the basis for the remainder of our analysis. Start with an $`L\times W`$ rectangle $`P`$, and fold it as follows. Glue the two opposite edges of length $`W`$ together to form a cylinder. Now glue the bottom rim of the cylinder to itself by creasing at two diametrically opposed points $`x_1`$ and $`y_1`$. Similarly glue the top rim to itself by creasing at two points $`x_2`$ and $`y_2`$. The gluing tree is of structure ‘I’: see Fig. 17.
It is easy to see this is an Aleksandrov gluing. Note both internal nodes of the gluing tree glue two $`\pi /2`$ rectangle corners to the interior of an $`L`$-edge; so the angle sum there is $`2\pi `$. The gluing is Aleksandrov even if the crease points on the top and bottom are not located at corresponding points on their rims. In particular, identify the points $`x_i`$ with their distance from the rectangle corner to the left. If $`x_1=x_2`$, then the crease points correspond, and the gluing produces a flat, $`L/2\times W`$ rectangle. If $`x_1x_2`$, the gluing is still Aleksandrov, but the “twist” in the gluing results in a nondegenerate tetrahedron, with all vertices of curvature $`\pi `$. Let $`x=|x_2x_1|`$ characterize amount of the twist, with $`x=0`$ representing no twist.
Because the $`1`$-skeleton of a tetrahedron is combinatorially $`K_4`$, each vertex is adjacent to all the others via polytope edges. This makes it trivial to decide the structure of the polytope $`Q_x`$ created by this rectangle gluing with twist $`x`$. The six distances between pairs of vertices are easily computed from the gluing, and each represents an edge length. These six lengths uniquely determine the 3D shape of the tetrahedron. It is not difficult to compute 3D vertex coordinates from the six lengths, and we have written code for this computation. An example is shown in Fig. 18. Here a $`2\times 2`$ rectangle is folded with a variety of different twists $`x`$. For both $`x=0`$ and $`x=1`$, the result is a flat $`1\times 2`$ rectangle, with a smooth interpolation between for $`0<x<1`$.
We have proven this lemma:
###### Lemma 5.7
Any rectangle may fold via a ‘I’ gluing tree to a uncountably infinite number of noncongruent tetrahedra.
Proof: Two tetrahedra with different edges lengths are not congruent. The edge lengths of $`Q_x`$ for twist $`x`$ are $`L/2`$ (twice), $`u(x)=\sqrt{x^2+W^2}`$ (twice), and $`v(x)=\sqrt{(1x)^2+W^2}`$ (twice). For $`ab`$, $`u(a)u(b)`$; and for $`x<L/4`$, $`u(x)v(x)`$. Thus the number of noncongruent tetrahedra is at least the number of distinct $`x[0,L/4)`$, which is nondenumerable. $`\mathrm{}`$
We return now to the task of upper-bounding the number of four fold-point gluing trees possible for a polygon of $`n`$ vertices. Although we do not at this point have tight bounds, they suffice for our purposes in the next section.
Define a conveyor belt (or just belt) in a gluing tree to be a path between two leaf fold points. Let a belt have fold points $`x`$ and $`y`$, with $`x`$ an interior point of edge $`e`$. A belt can roll if there is a nonzero-length interval $`Ie`$ such that for every $`xI`$, the belt folded at $`x`$ is an Aleksandrov gluing. A belt could instead have a finite number of distinct gluings, perhaps just one. We first show that rolling belts must be vertex-free in four fold-point trees.
###### Lemma 5.8
A rolling belt in a four fold-point tree $`T`$ cannot contain any vertices except those at the attachment points to other branches of $`T`$.
Proof: Suppose to the contrary that a rolling belt contains at least one vertex $`v`$ in its interior, i.e., not at an attachment point. Because under our definition, all vertices of $`P`$ are essential, the internal angle at $`v`$ is different from $`\pi `$. Let $`xI`$ be a particular fold point that determines the gluing of the belt. In this position, $`v`$ must match up with another vertex $`v^{}`$ with supplementary angle. Rolling the belt in a neighborhood of $`x`$ breaks the match, leaving both $`v`$ and $`v^{}`$ glued to points internal to an edge. At these points, the curvature is greater than zero, violating the fact that all curvature at a four fold-point gluing are concentrated at the leaves. $`\mathrm{}`$
Note that the angles at the attachment points must be $`\pi `$.
###### Lemma 5.9
A belt in a four-fold gluing tree $`T`$ has at most $`O(n)`$ combinatorially distinct gluings.
Proof: Let $`B`$ be a belt with attachment points $`a`$ and $`b`$. Note that because each attachment point is an internal node of $`T`$, the limited structural possibilities established in Lemma 5.6 allow only one or two attachment points. Consider two cases:
1. $`B`$ can roll. Then by Lemma 5.8, $`B`$ contains no internal vertices. Thus its only vertices are at $`a`$ and $`b`$. Rolling can produce just two combinatorially distinct positions of the belt.
2. $`B`$ cannot roll. Then $`B`$ can assume a finite number of possible positions. Define a kink in $`B`$ to be either a vertex, or an attachment point at which the angle is different from $`\pi `$, composed of two glued vertices. The kinks must match up in pairs. Matching one pair forces the remaining matches. Thus This can be seen by distributing the kinks around a topological circle representing $`B`$. Once one chord $`v_1v_i`$ is drawn in this circle, all other chords are forced by the pairwise matching requirement. Because there are only $`m1<n`$ choices for the mate for $`v_1`$, $`B`$ has only $`O(n)`$ legal gluings.
$`\mathrm{}`$
###### Lemma 5.10
The number of four fold-point gluing trees for a polygon of $`n`$ vertices is $`\mathrm{\Omega }(n^2)`$ and $`O(n^4)`$.
Proof: The lower bound is established by a variation on the foldings of a rectangle to tetrahedra (Fig. 18). The idea is to make each of the two conveyor belts in a ‘I’ structure (Fig. 17) realize $`\mathrm{\Omega }(n)`$ gluings independently. This can be accomplished by alternating supplementary angles along the belt at equal intervals. This is illustrated in Fig. 19 with angles $`\pi /2`$ and $`3\pi /2`$. The figure illustrates one possible folding; the fold points are midpoints of edges. The tetrahedra produced are the same as that obtained by folding a rectangle: the “teeth” mesh seamlessly.
For the upper bound, Lemma 5.6 restricts the structures to ‘+’ and ‘I’.
1. +’. Here we rely on the crude $`O(n^4)`$ bound determined by the four vertices, or three vertices and one edge, glued together to form the interior node of $`T_G`$. This fixes the combinatorial type of the gluing, which by Lemma 5.3 determines $`T_G`$.
2. I’. Let $`a`$ and $`b`$ be the upper and lower nodes of the ‘I’. There are two cases to consider for the upper node:
1. $`a`$ is of type ‘vv’: vertices $`v_i`$ and $`v_j`$ glue to form the belt attachment point. Then the path from $`a`$ to $`b`$ is determined by the requirement that the curvature must be zero at each point: the two sides “zip” closed from $`v_i`$/$`v_j`$ until the first point at which the curvature is nonzero, which then must be the lower node $`b`$.
2. $`a`$ is of type ‘ve’: vertex $`v_i`$ glues to the interior of edge $`e_j`$ to form the attachment point. The “zipping” down to $`b`$ is again determined, but this takes more argument. Let $`v_{i+1}`$ and $`v_j`$ be the two vertices closest to $`a`$ on the path to $`b`$. Both of their angles must differ from $`\pi `$ (because all vertices are essential). They must glue to one another with an angle sum of $`2\pi `$ (because the curvature must be zero). We want to show that $`v_i`$ cannot “slide” along $`e_j`$ to another position and still result in an Aleksandrov gluing. Sliding $`v_i`$ up $`e_j`$ places $`v_{i+1}`$ in the interior of $`e_j`$, and sliding $`v_i`$ down places $`v_j`$ in the interior of $`e_i`$, in both cases producing a point of nonzero curvature. Therefore no sliding is possible. Because any respositioning of $`v_i`$ on $`e_j`$ can be viewed as such a sliding, no repositioning is possible.
In both cases there are at most $`O(n^2)`$ choices for the constituents glued at $`a`$. Together with the $`O(n^2)`$ bound on the two belts from Lemma 5.9, this establishes the claimed $`O(n^4)`$ bound. $`\mathrm{}`$
It seems likely that this lemma could be strengthened:
###### Conjecture 5.3
The number of four fold-point gluing trees for a polygon of $`n`$ vertices is $`\mathrm{\Theta }(n^2)`$.
### 5.4 Upper Bound: General Case
We finally are positioned to establish an upper bound on the number of gluing trees, as a function of the number leaves.
###### Theorem 5.11
The number of gluing trees with $`\lambda `$ leaves for a polygon $`P`$ with $`n`$ vertices is $`O(n^{2\lambda 2})`$.
Proof: Let $`g(n,\lambda )`$ be the number of gluing trees for $`P`$ that have $`\lambda `$ leaves. The proof is by induction on $`\lambda `$. We know from Lemma 5.6 that at most four leaves can be fold-points. We assume for the general step of the induction that $`\lambda >4`$, and so there is at least one non-fold-point leaf. The base cases for $`\lambda 4`$ will be considered later.
The bound will use one consequence of the angles or curvature of a gluing (described in this paragraph), and one consequence of the matching edge lengths of a gluing (described in the next paragraph). Because a point interior to an edge of $`P`$ has angle $`\pi `$, a node of degree $`d`$ of a gluing tree ($`d=1,2,\mathrm{}`$) glues together $`d`$ vertices of $`P`$ or $`d1`$ vertices and one edge of $`P`$. Apart from this, we will use nothing else about the angles of the polygon, and in fact, our argument will hold more generally for a closed chain of $`n`$ vertices, with specified edge lengths.
Given a tree $`T_G`$ that is not a path, and a leaf $`l`$, define the source of $`l`$ as the first node of degree more than 2 along the (unique) path from $`l`$ into $`T`$. The path in $`T_G`$ from $`l`$ to its source is called the branch of $`l`$. For a tree $`T_G`$ and a non-fold-point leaf corresponding to polygon vertex $`l`$, let $`s(l)`$ be a vertex of $`P`$ closest to $`l`$ glued at the source of the leaf. Note that there must be such a vertex, since we cannot glue together two points interior to polygon edges at the source of the leaf. For example, in Fig. 3, $`s(v_2)`$ can be $`v_3`$ or $`v_1`$. Note—this is the single consequence of matching edge lengths referred to above—that the pair $`(l,s(l))`$ determines the portion of $`P`$’s boundary that is glued together to form the branch of $`l`$. We can simplify $`T`$ by cutting off $`l`$’s branch, resulting in a tree with $`\lambda 1`$ leaves. The corresponding simplification of $`P`$ is to excise the portion of its chain of length $`2d(l,s(l))`$ centered at $`l`$, resulting in a closed chain on at most $`n1`$ vertices. Since there are $`n`$ choices for $`l`$ and at most $`n`$ choices for $`s(l)`$ we obtain $`g(n,\lambda )n^2g(n1,\lambda 1)`$. For the general case there are at most 3 fold leaves, hence: $`g(n,\lambda )n^{2(\lambda 3)}g(n(\lambda 3),3)`$.
Lemmas 5.4 and 5.5 established the base cases $`g(n,2)=O(n^2)`$ and $`g(n,3)=O(n^4)`$. Substituting, this yields
$`g(n,\lambda )`$ $`n^{2(\lambda 3)}O([n(\lambda 3)]^4)`$ (8)
$`=`$ $`n^{2(\lambda 3)}O(n^4)`$ (9)
$`=`$ $`O(n^{2\lambda 2})`$ (10)
It remains to handle the case of $`\lambda =4`$ leaves. We will separate into the cases when at least one of these leaves is not a fold-point leaf, where arguments as above yield $`O(n^6)`$, and the case when all 4 vertices are fold leaves. In this case, Lemma 5.10 establishes a bound of $`O(n^4)`$, smaller than that claimed by the lemma. $`\mathrm{}`$
Of course because $`\lambda `$ could be $`\mathrm{\Omega }(n)`$, there is no contradiction between this upper bound and the exponential lower bound in Theorem 5.2. We specialize the upper bound to convex polygons in Section 5.6, but first we summarize the structural characteristics of gluing trees we have uncovered.
### 5.5 Gluing Tree Characterization
Our previous results imply that gluing trees are fundamentally discrete structures, with one or two rolling conveyor belts, and two such belts only in very special circumstances.
###### Theorem 5.12
Gluing trees satisfy these properties:
1. At any gluing tree point of degree $`d2`$, at most one point of $`P`$ in the interior of an edge may be glued, i.e., at most one nonvertex may be glued there.
2. At most four leaves of the gluing tree can be fold points, i.e., points in the interior of an edge of $`P`$. The case of four fold-point leaves is only possible when the tree has exactly four leaves, with the combinatorial structure ‘+’ or ‘I’.
3. A gluing tree can have at most two rolling belts.
4. A gluing tree with two rolling conveyor belts must have the structure ‘I’, and result from folding a polygon that can be viewed as a quadrilateral with two of its opposite edges replaced by complimentary polygonal paths.
Proof:
1. That $`d2`$ points of a gluing tree have at most one edge-interior points glued is immediate from the definition of an Aleksandrov gluing, and our insistence that all vertices are essential.
2. The structure of four fold-point trees was established in Lemma 5.6.
3. The definition of “rolling belts” (p. 5.7) implies four fold points, so the constraints from the previous item apply.
4. Two rolling belts cannot be accommodated by the ‘+’ structure, which is determined by the four vertices glued at the central node. So the tree structure must be ‘I’. Lemma 5.8 established that the belts are vertex-free, corresponding to two opposite edges of the quadrilateral. The central path of the ‘I’ must be formed by gluing vertices together whose angle sum is $`2\pi `$, and they are in this sense complimentary polygonal paths.
$`\mathrm{}`$
Thus a generic gluing tree has one rolling belt, with trees hanging off it, and one of those trees having a fold-point leaf. See Fig. 20.
### 5.6 Upper Bound: Convex Polygons
For convex polygons we can prove a polynomial upper bound. We first handle the special case of $`\lambda =4`$.
###### Lemma 5.13
A convex polygon $`P`$ may fold to gluing trees of four leaves only if it is a quadrilateral, a pentagon, or a hexagon; $`P`$ may fold to $`O(1)`$ such gluing trees.
Proof: As in the proof of Theorem 4.6, the two conditions $`\gamma (v)\pi `$ and $`_v\gamma (v)=4\pi `$ for the four leaves $`v`$ of the tree imply that $`\gamma (v)=\pi `$ for each. This implies that the internal angle at $`v`$ in $`P`$ is $`\pi `$, which, by our assumption that all vertices are essential, implies that all four are fold vertices.
Because all available curvature is consumed by the four leaves, the internal nodes of the gluing tree must be flat. If the $`T`$ has shape ‘+’, four vertices whose angles sum to $`2\pi `$ join there. Recalling that the turn angle at each vertex is $`\tau _i=\pi \alpha _i`$ and that the total turn angle is $`2\pi `$, this angle sum implies that $`_i\tau _i=4\pi 2\pi =2\pi `$, for the four vertices at the ‘+’, and so the turn angle is completely consumed by these four vertices. Thus $`P`$ must be a quadrilateral, and there is just one way to form the gluing tree.
If $`T`$ is a ‘I’ shape, then each of the two internal nodes of the ‘I’ are formed either by gluing together three vertices, or two vertices and an edge. For a three-vertex node, the turn angle sum is $`3\pi 2\pi =\pi `$; for a two-vertex and edge node, the turn angle sum is $`2\pi \pi =\pi `$. So both nodes together consume of all the $`2\pi `$ turn angle. Therefore $`P`$ has at most six vertices. The hexagon permits the most groupings of vertices, six; and so there are at most six gluing trees. $`\mathrm{}`$
See Fig. 21 for an irregular hexagon that folds with a ‘I’ gluing tree.
###### Theorem 5.14
A convex polygon $`P`$ of $`n`$ vertices folds to at most $`O(n^3)`$ different gluing trees.
Proof: Theorem 4.6 limits the combinatorial possibilities to trees with four or fewer leaves. We have settled each case for $`\lambda 4`$ earlier:
Lemma 5.13: $`O(1)`$.
Lemma 5.5: $`O(n^4)`$.
Lemma 5.4: $`O(n^2)`$.
We now improve the $`\lambda =3`$ case to $`O(n^3)`$ for convex polygons. We can tighten the $`O(n^4)`$ bound with the following two observations:
1. The two internal angles at the two vertices glued at a type-vve node must sum to no more than $`\pi `$.
2. A convex polygon cannot have too many vertices with small angles.
To quantify the second observation, define the turn angle $`\tau _i`$ at a vertex $`v_i`$ with internal angle $`\alpha _i`$ to be $`\tau _i=\pi \alpha _i`$. For any polyon, we must have $`_i=2\pi `$. For a convex polygon, $`\tau _i>0`$. Now suppose $`v_i`$ and $`v_j`$ glue at a type-vve node. Then $`\alpha _i+\alpha _j\pi `$, and so $`\tau _i+\tau _j\pi `$. Thus two distinct vv-gluings, involving four different vertices, already consume the available $`2\pi `$ turn angle of the polygon. Call two pairs of vertices disjoint if all four vertices are different. The turn angle bound implies that any polygon can have at most two disjoint pairs of vertices glued to a type-vve node. We now show that this implies a $`O(n)`$ bound on the number of type-vve nodes.
Construct a bipartite graph $`H`$ with $`n`$ nodes for the first vertex and $`n`$ nodes for the second vertex of vv-gluings. Let $`(v_i,v_j)`$ and $`(v_k,v_l)`$ be disjoint pairs in vv-gluings, as depicted in Fig. 16(b). Then because there cannot be another pair disjoint from either of these, every other pair must be incident to one of the four vertices $`v_i,v_j,v_k,v_l`$. This limits $`H`$ to at most $`4n`$ edges (and even this bound is loose, for this permits as many as four disjoint pairs, as is evident in the figure).
Thus there are at most $`O(n)`$ type-vve nodes. Repeating the argument that a vv-gluing determines one leg of the ‘Y’ and Lemma 5.4 bounds the remaining path to $`O(n^2)`$ possibilities, leads to the claimed $`O(n^3)`$ bound. $`\mathrm{}`$
We leave open the question of whether this bound is tight.
It is straightforward to list all possible gluing trees for a given convex polygon with an $`O(n^3\mathrm{log}n)`$ time algorithm. We have implemented such an algorithm, with, however, less than maximally efficient data structures.
## 6 Counting Foldings: Noncongruent Polytopes
We have so far been counting the number of different ways to fold up a given polygon, but have not addressed the question of whether all these foldings produce distinct polytopes. There are several notions of what constitutes distinctness. One natural definition relies on the combinatorial structure of the polytope, as explored by Shephard \[She75\]. We will have little to say on this topic here. Instead, we will focus on counting noncongruent polytopes.
We have already established in Lemma 5.7 that any rectangle can fold to an uncountably infinite number of noncongruent tetrahedra. We extend this result in this section to the “obvious” fact that any convex polygon folds (via perimeter-halving) to an uncountably infinite number of noncongruent polytopes. Despite the naturalness of this claim, our inability to determine the 3D structure of the polytope guaranteed by an Aleksandrov gluing makes our proof less than satisfactory. In the absence of any 3D information, we concentrate instead on the pattern of geodesics between vertices, for of course two congruent polytopes have the exact same set of geodesics.
###### Lemma 6.1
A polytope $`Q`$ resulting from a perimeter-halving fold of polygon $`P`$ has a countable number of geodesics between any pair of vertices.
Proof: Let $`x`$ and $`y`$ be the fold vertices produced by the perimeter-halving (as in Fig. 1). We will assign each geodesic a unique integer, which establishes that there are only a countable number of them. The integers are based on a “layout” of the surface of $`Q`$ in the plane. Fix $`P`$ in the plane, and designate it as level-$`0`$ of the layout. Around $`P`$ layout $`2n`$ copies of $`P`$ (where $`P`$ has $`n`$ vertices) corresponding to the perimeter gluing. These are level-$`1`$ $`P`$ copies of the layout. This level is illustrated in Fig. 22. For example, because edge $`e_4=v_4v_5`$ of $`P`$ is glued into the edge $`e_1=v_1v_2`$ by the perimeter halving, a level-$`1`$ copy of $`P`$ is placed exterior to $`e_4`$ arranged so that the glued portions of $`e_4`$ and $`e_1`$ match. There are $`2n`$ level-$`1`$ copies of $`P`$ because the $`n`$ vertices around $`P`$ are interspersed by a reversed sequence of the same $`n`$ vertices.
Continuing the construction, level-$`i`$ of the layout is formed by surrounding each level-$`(i1)`$ copy of $`P`$ with $`2n`$ additional copies. Give these copies a “sequence number” $`j=1,\mathrm{},2n`$. Now every copy of $`P`$ at level-$`i`$ in the layout may be assigned a unique integer by listing the sequence numbers for each level $`0,\mathrm{},i`$ and interpreting it as a base-$`2n`$ number.
It is clear from the layout construction that any geodesic on $`Q`$ “unrolls” to a straightline in the layout. Because we can number the copies of $`P`$, we can number the geodesics between any given pair of vertices. Therefore the number of geodesics is denumerable. $`\mathrm{}`$
Although this proof is specialized to polytopes formed from perimeter halving, it would not be difficult to extend it to all polytopes formed by gluings including a “rolling” fold-point.
###### Theorem 6.2
Any convex polygon $`P`$ folds, via perimeter halving, to a uncountably infinite number of noncongruent polytopes.
Proof: Select $`x`$, a fold point for perimeter halving, interior to an edge $`e_i=v_iv_{i+1}`$ of $`P`$. The segment $`xv_iP`$ in level-$`0`$ of the layout used in the previous lemma corresponds to a geodesic on $`Q_x`$. Now let $`x`$ vary within some neighborhood $`Ne_i`$; let $`x^{}x`$ be a point in $`N`$. The segment $`x^{}v_i`$ corresponds to a geodesic on $`Q_x^{}`$ of a different length. We use this fact to establish our claim.
Let $`𝒬=\{Q_x^{}:x^{}N\}`$ be the set of all the polytopes produced as $`x`$ varies over the neighborhood. Assume, for the purposes of contradiction, that the number of distinct, noncongruent polytopes in $`𝒬`$ is denumerable: $`Q_1,Q_2,\mathrm{}`$. By Lemma 6.1, each has a countable number of geodesics: a pair of numbers suffice to uniquely identify them. Thus the total number of distinct lengths of geodesics represented by all these polytopes is denumerable. But this contradicts the nondenumerable number of lengths of segments $`|x^{}v_i|`$ for $`x^{}N`$. Therefore the number of noncongruent polytopes in $`𝒬`$ is nondenumerable. $`\mathrm{}`$
Although this theorem establishes the result even for regular polygons, there is much more to say about the structure of the polytopes that can be folded from regular polygons. We explore this in Section 9.
## 7 Counting Unfoldings: Cut Trees
In this section we explore unfolding from the point of view of cut trees. The general situation is that we are given one polytope $`Q`$ of $`n`$ vertices, and we would like to know how many different ways it can be cut and unfolded to a polygon. We start with some straightforward observations before proving enumeration bounds.
First, every polytope admits at least the $`n`$ cut trees provided by the star unfolding \[AO92\], one with each vertex as source. So in particular, every polytope unfolds to at least one polygon. (As we mentioned in the Introduction, the corresponding question for edge-unfoldings remains open.)
Second, because we permit arbitary polygonal paths between the nodes of a cut tree (Section 3), there is no upper bound on the number of polygon vertices in potential unfoldings of a given polytope. This might lead one to wonder if *any* polygon (of the appropriate area) can be unfolded from a given polytope. The answer is no, as is easily established by this lemma.
###### Lemma 7.1
Every polygon $`P`$ cut from $`Q`$ must have at least two vertices whose interior angles are of the form $`2\pi \gamma _i`$ for some $`i=1,\mathrm{},n`$, where $`\gamma _i`$ are the curvatures of the vertices of $`Q`$.
Proof: Let the $`n`$ vertices of $`Q`$ have curvatures $`\gamma _i`$, $`i=1,\mathrm{},n`$. The cut tree $`T_C`$ must have at least two leaves, and by Lemma 3.1 these leaves must be vertices of $`Q`$. Say they coincide with the vertices of curvatures $`\gamma _1`$ and $`\gamma _2`$. Then any polygon $`P`$ that unfolds from $`T_C`$ must have two vertices with interior angles $`2\pi \gamma _1`$ and $`2\pi \gamma _2`$. $`\mathrm{}`$
So let $`P`$ be a polygon with no interior angle equal to $`2\pi \gamma _i`$ for $`i=1,\mathrm{},n`$. Then $`P`$ cannot be cut from $`Q`$.
### 7.1 Lower Bound: Exponential Number of Unfoldings
In this section we provide an exponential lower bound.
###### Theorem 7.2
There is a polytope $`Q`$ of $`n`$ vertices that may be cut open with exponentially many ($`2^{\mathrm{\Omega }(n)}`$) combinatorially distinct cut trees, which unfold to exponentially many geometrically distinct simple polygons.
Proof: $`Q`$ is a truncated cone, as illustrated in Fig. 23:
the hull of two regular $`n`$-gons of different radii, lying in parallel planes and similarly oriented. We call this the volcano example. We require that $`n`$ be even; in the figure, $`n=16`$. Label the vertices on the top face $`a_0,\mathrm{},a_{n1}`$ and $`b_0,\mathrm{},b_{n1}`$ correspondingly on the bottom face. The “base” cut tree, which we notate as $`T_{0000000}`$, unfolds $`Q`$ as shown in Fig. 25. $`T_{0000000}`$ consists of a path on the top face $`(a_0,a_1,\mathrm{},a_{n1})`$ supplemented by arcs $`(a_i,b_i)`$ for all $`i=0,\mathrm{},n1`$. The polygon $`P`$ produced consists of the base face, $`n`$ attached trapezoids $`(b_i,b_{i+1},a_{i+1},a_i)`$, with the top face attached to $`a_{n1}a_0`$.
Define a cut tree $`T_{m_{(n1)/2}\mathrm{}m_2m_1m_0}`$, where $`m_i`$ are the digits of a binary number of $`n/21`$ bits, as an alteration of the base tree $`T_{0\mathrm{}0}`$ as follows. If $`m_i=1`$, then the arc $`(a_{2i+1},b_{2i+1})`$ is deleted and replaced by $`(a_{2i},b_{2i+1})`$. If $`m_i=0`$, then the arc $`(a_{2i+1},b_{2i+1})`$ is used as in $`T_{0\mathrm{}0}`$. Thus the cut tree $`T_{1001101}`$ shown in Fig. 25 replaces $`(a_1,b_1)`$ with $`(a_0,b_1)`$ because $`m_0=1`$, $`(a_5,b_5)`$ with $`(a_4,b_5)`$ because $`m_2=1`$, and so on.
There are $`2^{n/21}=2^{\mathrm{\Omega }(n)}`$ cut trees.
It is clear by construction that all these cut trees lead to simple polygon unfoldings. It only remains to argue that each leads to a distinct polygon, not congruent to any other. This is not strictly true for $`Q`$ as defined, for any bit pattern leads to a $`P`$ that is congruent (by reflection) to the polygon obtained from the reverse of the bit pattern. However, it is a simple matter to introduce some asymmetry, by, for example, lengthening edge $`a_{n1}a_0`$ slightly. Then all cut trees lead to distinct polygons. $`\mathrm{}`$
A simpler example is a drum, the convex hull of two regular polygons in parallel planes. Because some of the unfoldings used in the above proof overlap, there is a bit more argument needed to establish the exponential lower bound.
Even restricting the cut tree to a path permits an exponential number of unfoldings:
###### Theorem 7.3
There is a polytope $`Q`$ of $`n`$ vertices that may be cut open with exponentially many ($`2^{\mathrm{\Omega }(n)}`$) combinatorially distinct cut trees, all of which are paths, which unfold to exponentially many geometrically distinct simple polygons.
Proof: $`Q`$ is formed by pasting two halves of a regular $`2n`$-gon together to form a semicircle approximation with some small thickness $`w>0`$. Label the vertices on the front face $`a_0,\mathrm{},a_n`$ and $`b_0,\mathrm{},b_n`$ correspondingly on the back face, as illustrated in Fig. 26(a). Let $`\alpha =2\pi /n`$ be the turn angle at each vertex $`a_i`$ (and at $`b_i`$), i.e., the angle $`\alpha =\pi \mathrm{}(a_{i1},a_i,a_{i+1})`$. (In the figure, $`\alpha =\pi /3211^{}`$.) We specify a series of cut trees $`T_m`$, where $`m`$ is an $`n`$-digit base-$`3`$ integer $`m_n\mathrm{}m_2m_1`$, with the following interpretation. $`T_{0\mathrm{}00}`$ is the “base” cut tree on which all others are variations:
$$T_{0\mathrm{}00}=(a_0,a_1,\mathrm{},a_n,b_n,b_{n1},\mathrm{},b_1,b_0)$$
(11)
Note that $`T_{0\mathrm{}00}`$ is a path, as are all the $`T_m`$. We call the half of the path on the front face the $`a`$-path, and that on the back the $`b`$-path. The unfolding $`P`$ determined by $`T_{0\mathrm{}00}`$ is a regular $`2n`$-gon, fattened by a strip $`(a_0,b_0,b_n,a_n)`$ of width $`w`$ down its middle, with a “tail” of $`n`$ rectangles attached to edge $`a_0b_0`$. If $`|a_ia_{i+1}|=h`$, then each rectangle is $`w\times h`$.
In cut tree $`T_{m_n\mathrm{}m_2m_1}`$ the index $`m_i`$ is $`1`$ if the $`a`$-path deviates to touch $`b_i`$ on the back face via the path $`(\mathrm{},a_{i1},b_i,a_i,a_{i+1},\mathrm{})`$, and the index $`m_i`$ is $`2`$ if the $`b`$-path similarly deviates to include $`a_i`$ on the front face via the path $`(\mathrm{},b_{i1},a_i,b_i,b_{i+1},\mathrm{})`$. In both cases, the opposite path skips the vertex deviated to: if $`m_i=1`$, the $`b`$-path skips $`b_i`$ by shortcutting on the back face, and if $`m_i=2`$, the $`a`$-path skips $`a_i`$ by shortcutting on the front face. Fig. 26(a) illustrates $`T_{m_n\mathrm{}0022020100}`$, with $`a`$-path
$$(a_0,a_1,a_2,b_3,a_3,a_4,a_6,a_9,a_{10},\mathrm{},a_n)$$
(12)
and $`b`$-path
$$(b_0,b_1,b_2,b_4,a_5,b_5,b_6,a_7,b_7,a_8,b_8,b_9,b_{10},\mathrm{},b_n)$$
(13)
Note that when $`m_i0`$, the rectangle bounded between $`a_{i1}b_{i1}`$ and $`a_ib_i`$ is crossed by an $`ab`$-diagonal. We insist that $`m_1=0`$, so that the cut tree starts with an uncrossed rectangle $`(a_0,b_0,b_1,a_1)`$. Finally, the edge $`a_nb_n`$ is included in $`T_m`$, so that it is a path from $`a_0`$ to $`a_n`$ to $`b_n`$ and returning to $`b_0`$. The digits $`m_n\mathrm{}m_2`$ are each free to be any one of $`\{0,1,2\}`$. Thus there are an exponential number of combinatorially distinct $`T_m`$: $`3^{n1}`$. We return below to the issue of how many of these lead to geometrically distinct unfoldings.
It should be clear by construction that $`T_m`$ spans the vertices. To show that it is a tree, we need to argue that it is non-self-intersecting. This is again clear by construction, for each nonzero $`m_i`$ uses a diagonal in the rectangle prior to $`a_ib_i`$, and because $`m_i`$ has only one value, no such rectangle has both diagonals used. Together with the shortcutting that prevents the $`a`$\- and $`b`$-paths from touching the same vertex, it follows that $`T_m`$ is indeed a tree; so it is a legitimate cut tree. Thus it unfolds to a single piece. It only remains to show that this unfolding is a simple polygon, i.e., it avoids overlap.
This is obvious for $`T_{0\mathrm{}00}`$, as mentioned previously. For general $`T_m`$, consider the layout of the unfolding $`P`$ that places $`a_0b_0`$ horizontally, as in Fig. 26(b). Let $`L`$ be the horizontal line through $`a_1b_1`$; this segment is necessarily horizontal because we stipulated that $`m_1=0`$. We will argue that $`L`$ strictly separates the tail of $`P`$ (the portion attached above $`a_1b_1`$) from its body (the portion attached below $`a_1b_1`$).
First, it is clear that the body unfolds without overlap. For it is simply truncations (due to path shortcuttings) of halves of a regular polygon glued to either side of the rectangle $`(a_0,b_0,b_n,a_n)`$, with attached triangle “spikes” for each nonzero $`m_i`$. None of these spikes can overlap, even when adjacent, for their length-$`w`$ edge juts out orthogonal to their length-$`h`$ edge glued to the body (see the body image of $`b_7`$ in Fig. 26(b)).
The tail consists of $`h\times w`$ rectangles, or half-rectangles, glued end-to-end, with turns to the right by $`\alpha `$ for every $`m_i=1`$ digit, and turns to the left by $`\alpha `$ for every $`m_i=2`$ digit. Thus, $`T_{\mathrm{}0022020100}`$ in (b) of the figure turns right once and left three times. Because there are at most $`n1`$ nonzero digits, the tail can turn at most $`n1`$ times. Because $`\alpha `$ is the turn angle of a regular $`2n`$-gon, it takes $`n`$ turns of $`\alpha `$ to turn a full $`\pi `$. Thus the tail turns strictly less than $`\pi `$, and so cannot return to line $`L`$. Thus the tail remains strictly above $`L`$. Choosing $`w<h`$ guarantees that no body spike protrudes vertically as much as $`h`$ above $`a_0b_0`$; so the body remains strictly below $`L`$.
It remains to argue that the tail does not self-intersect. But this follows from the same turn argument above. By construction, there are no local overlaps between two adjacent tail rectangles or half-rectangles. Thus the only overlap conceivable would result from the tail curling back to overlap itself. Choosing $`wh`$ makes the tail essentially a series of segments of length $`h`$, with attached pieces of the regular polygon clipped by shortcutting. For the tail segments to curl back and overlap would require a total turn by at least $`\pi `$, contradicting the bound on the sum of $`\alpha `$’s.
Finally, we turn to the question of how many of the $`3^{n1}`$ combinatorially distinct $`T_m`$ lead to geometrically distinct (noncongruent) $`P`$. Let $`x`$ and $`y`$ be two base-$`3`$ numbers, and let $`S(x)`$ be the base-$`3`$ number obtained by changing each $`1`$-digit in $`x`$ to a $`2`$, and each $`2`$-digit in $`x`$ to a $`1`$. (For example, $`S(1021)=2012`$.) Then if $`S(x)=y`$, $`T_x`$ and $`T_y`$ lead to congruent $`P`$, in that $`P_y`$ is the reflection of $`P_x`$ about a vertical line (in the layout used above).
Although we could easily ensure noncongruency for all $`T_m`$ by altering $`Q`$ to be less symmetric, we opt here for a counting argument. Let $`x`$ be a base-$`3`$ number. Define $`B(x)`$ to be the binary number obtained by changing each $`2`$-digit in $`x`$ to a $`1`$. (For example, $`B(2012)=1011`$.) Now it should be clear that for any two base-$`3`$ numbers $`x`$ and $`y`$, if $`B(x)B(y)`$, then $`P_x`$ is noncongruent to $`P_y`$. For then the pattern of spikes on the body are different in $`P_x`$ and $`P_y`$. Thus, among the $`3^{n1}`$ combinatorially distinct $`P`$, there are at least $`2^{n1}`$ geometrically distinct $`P`$. $`\mathrm{}`$
### 7.2 Lower Bound: Convex Unfoldings
It seems possible that the exponential lower bound holds even in the case of convex unfoldings, via an example similar to that used in Fig. 26.
###### Conjecture 7.1
There is a polytope with an exponential number of convex unfoldings.
This represents the only ‘?’ in Table 1.
### 7.3 Upper Bound
###### Theorem 7.4
The maximum number of edge-unfolding cut trees of a polytope of $`n`$ vertices is $`2^{O(n)}`$, and the maximum number of arbitary cut trees $`2^{O(n^2)}`$.
Proof: For edge unfoldings, the bound depends on the number of spanning trees of a polytope graph. We may obtain a bound here as follows.<sup>6</sup><sup>6</sup>6 We thank B. McKay \[personal communication, Jan. 2000\] for guidance here. First, triangulating a planar graph only increases the number of spanning trees, so we may restrict attention to triangulated planar graphs. Second, it is well known that the number of spanning trees of a connected planar graph is the same as the number of spanning trees of its dual. So we focus just on $`3`$-regular (cubic) planar graphs. Finally, a result of McKay \[McK83\] proves an upper bound of $`O((16/3)^n/n)`$ on the number of spanning trees for cubic graphs. This bound is $`2^{O(n)}`$.
For arbitrary cut trees, the underlying graph might conceivably have a quadratic number of edges, which leads to the bound $`2^{O(n^2)}`$. (Note that our definition of cut tree in Section 3.1 would not count different polygonal paths between two vertices as distinct arcs of $`T_C`$.) $`\mathrm{}`$
## 8 Counting Unfoldings: Noncongruent Polygons
We have already seen in Theorem 7.2 that one polytope can have an exponential number of noncongruent polygon unfoldings. In fact the possibilities range from $`0`$ to $`\mathrm{}`$, even for convex unfoldings, as this simple counterpart of Theorem 6.2 shows:
###### Theorem 8.1
Although some polytopes unfold to a nondenumerable number of noncongruent convex polygons, others have only a finite number of convex unfoldings.
Proof: For the former claim, consider a doubly-covered equilateral triangle. Choose any point $`x`$ interior to the top face, as shown in Fig. 27(a). This leads to a ‘Y’ cut tree that unfolds to a convex polygon (b) for every choice of $`x`$. All these polygons have different angles, and so are noncongruent.
The second claim of the theorem is trivially satisfied by polytopes with zero convex unfoldings. To establish it for a polytope that has at least one convex unfolding is more difficult, and we only sketch a construction. Consider the doubly-covered trapezoid shown in Fig. 28. It has just two sharp vertices, $`v_1`$ and $`v_4`$, and so, by Theorem 4.6, the cut tree must be a path connecting those vertices. The path $`(v_1,v_2,v_3,v_4)`$ unfolds to a convex polygon. Now consider a geodesic that starts with the segment $`v_1v_3`$ as illustrated. As in the proof of Lemma 4.5, this geodesic will either hit $`v_4`$ directly, in which case it is not a valid cut tree because $`v_2`$ is not spanned, or it spirals around the trapezoid and self-crosses. We will not prove this claim. $`\mathrm{}`$
## 9 Folding Regular Polygons
In this section we study folding *regular* polygons of $`n`$ vertices. Because all polygon vertices have the same interior angle $`\theta =(n2)\pi /n`$, only a limited variety of different polytope vertex curvatures may be created. We find, not surprisingly, that this leads to a limited set of possibilities: in general, only one “class” of nonflat polytopes can be produced. This is established in Lemma 9.2.
Let $`\alpha _k`$, $`k1`$, be the curvature at a polytope vertex formed by gluing $`k`$ $`P`$-angles of $`\theta `$ together, and $`\beta _k`$, $`k0`$, be the curvature at a vertex formed by gluing $`k`$ angles to a point interior to an edge of $`P`$. The next lemma details the possible $`\alpha _k`$ and $`\beta _k`$ values achievable.
Throughout this section we will find that the situation is more uniform for $`n>6`$ than it is for small $`n`$.
###### Lemma 9.1
For $`n>6`$, only four vertex curvatures can be obtained by folding a regular $`n`$-gon $`P`$; for $`n6`$, additional curvature values are possible. More precisely, for all $`n`$, these four curvature values are always achievable:
* $`\alpha _1=\pi (1+2/n)`$.
* $`\alpha _2=\pi (4/n)`$.
* $`\beta _0=\pi `$.
* $`\beta _1=\pi (2/n)`$.
The additional values possible for $`n6`$ are detailed in Tables 2 and 3.
Proof:
1. $`\alpha _1=2\pi \theta =2\pi (n2)\pi /n=\pi (1+2/n)`$. This vertex is a leaf of the gluing/cut tree. We call this a zipped vertex, for $`P`$ is “zipped shut” at the vertex.
2. $`\alpha _2=2\pi 2\theta =2\pi 2(n2)\pi /n=\pi (4/n)`$. This vertex is a degree-$`1`$ node of the gluing tree.
3. $`\beta _0=\pi `$. This is a fold vertex, when nothing is glued to an edge of $`P`$, and therefore a leaf of the gluing tree.
4. $`\beta _1=2\pi [\pi +\theta ]=2\pi [\pi +(n2)\pi /n]=\pi (2/n)`$. This is a degree-$`1`$ node of the gluing tree.
The additional possibilities for $`n6`$ are as follows. $`\alpha _3`$ is possible for all $`n6`$; $`\alpha _4`$ is possible only for $`n=4`$; and no other $`\alpha _k`$ is possible. See Table 2.
For $`n=3`$, $`\beta _2,\beta _3`$, and for $`n=4`$, $`\beta _2`$, are all possible. See Table 3.
Explicit computation shows that all higher values of $`k`$ lead to nonconvex vertices, whose total face angle exceeds $`2\pi `$ and so which have negative curvature. $`\mathrm{}`$
Let $`a_i`$ and $`b_i`$ be the number of polytope vertices of curvature $`\alpha _i`$ and $`\beta _i`$ respectively, formed by folding a regular $`n`$-gon $`P`$. Of course $`a_i`$ and $`b_i`$ are nonnegative integers, but there are additional significant restrictions imposed by the requirement that the total curvature be $`4\pi `$:
$$\underset{i=1}{}a_i\alpha _i+\underset{i=0}{}b_i\beta _i=4\pi .$$
(14)
We now explore the implications of this constraint, separately for $`n>6`$ and for $`n6`$. Note that our notation implies that
$$\underset{i=1}{}a_ii+\underset{i=0}{}b_ii=n,$$
(15)
because the subscripts on $`\alpha `$ and $`\beta `$ indicate the number of vertices involved in the gluing.
Now we prove that perimeter-halving is the only possible kind of folding for $`n>6`$.
###### Lemma 9.2
For all $`n3`$, regular $`n`$-gons fold via perimeter-halving, using path gluing trees, to two classes of polytopes:
1. A continuum of “pita” polytopes of $`n+2`$ vertices.
2. One or two flat, “half-$`n`$-gons”:
1. $`n`$ even: Two flat polytopes, of $`\frac{n}{2}+2`$ and $`\frac{n}{2}+1`$ vertices.
2. $`n`$ odd: One flat polytope, of $`\frac{n+1}{2}+1`$ vertices.
For $`n>6`$, these are the only foldings possible of a regular $`n`$-gon.
Proof: A perimeter-halving fold produces a path gluing tree. This has two leaves and all other nodes internal. From Lemma 9.1, the only two curvatures can be leaves: $`\{\alpha _1,\beta _0\}`$; and only two can be degree-$`1`$ nodes: $`\{\alpha _2,\beta _1\}`$. Moreover, these are the only curvatures possible for $`n>6`$. Thus Eq. (14) reduces to
$$a_1\alpha _1+a_2\alpha _2+b_0\beta _0+b_1\beta _1=4\pi .$$
(16)
Substituting the curvature values from Lemma 9.1 and solving for $`n`$ yields
$$n=\frac{2(a_1+2a_2+b_1)}{4(a_1+b_0)}$$
(17)
Because only $`\alpha _1`$ and $`\beta _0`$ are leaf vertex curvatures, we must have $`a_1+b_02`$. The requirement that the denominator of Eq. (17) be positive yields $`a_1+b_0<4`$. Therefore we know that $`a_1+b_0\{2,3\}`$. We now show that the case $`a_1+b_0=3`$ is not possible when $`n>6`$.
As both $`a_1`$ and $`b_0`$ count leaves, a tree formed with $`a_1+b_0=3`$ must have at least three leaves. By Theorem 4.6, because $`n4`$, it cannot have more than three leaves. So it has exactly three leaves, and has the combinatorial structure of a ‘Y’. The interior node must be formed by gluing three distinct points of $`P`$ together (by Lemma 3.1(4)). This corresponds to curvatures $`\alpha _k`$, $`k3`$, or $`\beta _k`$, $`k2`$. But Lemma 9.1 shows that none of these are possible for $`n>6`$. (Note, for later reference, that for $`n6`$, these possibilities will need consideration.)
Therefore we must have $`a_1+b_0=2`$. Therefore the gluing tree must be a path for $`n>6`$, and the folding must be a perimeter-halving folding. We now explore the three possible solutions to $`a_1+b_0=2`$.
The two leaves are both folds at interior points of edges of $`P`$, a perimeter-halving folding similar to that previously illustrated in Fig. 1. If neither fold point $`x`$ and $`y`$ is the midpoint of its edge, then no pair of vertices glue together, so $`a_2=0`$ and therefore $`b_1=n`$. This produces a continuum of polytopes $`Q_x`$ of $`n+2`$ vertices. We call these pita polytopes (Fig. 32), and will study them in Section 9.1 below.
Suppose one fold point $`x`$ is at an edge midpoint. If $`n`$ is even, then $`y`$ is also at a midpoint, and $`P`$’s vertices are glued in pairs. Therefore $`a_2=n/2`$ and $`b_1=0`$. The polytope is a flat half-$`n`$-gon of $`n/2+2`$ vertices. See Fig. 29(a). If $`n`$ is odd, then $`y`$ must be at a vertex. This means that $`a_10`$, and this case does not apply.
Both leaves are at vertices, and so $`n`$ must be even. All other vertices are glued in pairs, so $`a_2=(n2)/2`$ and $`b_1=0`$. The folding produces a flat half-$`n`$-gon of $`n/2+1`$ vertices. See Fig. 29(b).
One vertex is zipped to a leaf; half the perimeter around is a fold vertex. This implies that $`n`$ is odd. All other vertices are glued in pairs, so $`a_2=(n1)/2`$ and $`b_1=0`$. The folding produces a flat half-$`n`$-gon of $`(n1)/2+2`$ vertices. See Fig. 29(c).
The details derived above are gathered into Table 4, and the flat foldings illustrated in Fig. 29. $`\mathrm{}`$
###### Lemma 9.3
For $`n6`$, regular polygons fold to additional polytopes (beyond those listed in Lemma 9.2) as detailed in Table 5.
Proof: Lemma 3.1 limits the possible nonpath cut trees to ‘Y’, ‘+’, and ‘I’. We first argue that ‘I’ is only possible for $`n=6`$. The two interior nodes of the tree must have curvatures in $`\{\alpha _3,\beta _2\}`$. For $`n=3`$, there are not enough vertices to make these nodes. For $`n=4`$, there are enough vertices to make two $`\beta _2`$ nodes, but this then forces the ‘+’ structure, i.e., the interior edge of the ‘I’ has length zero. For $`n=5`$ and $`n=6`$, $`\beta _2`$ is not possible. For $`n=5`$, there are not enough vertices to make two $`\alpha _3`$ vertices. And finally, for $`n=6`$, there are enough vertices, and the folding produces a flat rectangle.
Thus only ‘Y’ and ‘+’ are possible. The ‘+’ can only be realized in two ways: by gluing four vertices together, which is only possible for $`n=4`$ (see $`\alpha _4`$ column in Table 2), and by gluing three vertices to an edge, which is only possible for $`n=3`$ (see $`\beta _3`$ column in Table 3).
There are a number of ways to realize ‘Y’-trees. The constraint that the curvature add to $`4\pi `$, Eq. (14), together with the discrete set of possible curvatures in Tables 2 and 3, lead to the possibilities listed in Table 5. (The second line of the table was previously illustrated in Fig. 3.) $`\mathrm{}`$
If we treat the $`n`$ vertices of a regular $`n`$-gon as assigned the same label (as seems appropriate), Lemmas 9.2 and 9.3 together show that there are only $`O(1)`$ ways to fold up a regular polygon, justifying the entry in Table 1. If we label the vertices with distinct labels, then there are $`O(n)`$ foldings.
### 9.1 Pita Polytopes
We define a pita polytope as one obtained by a perimeter-halving folding of a regular polygon at a point on an edge that is not a midpoint, as per the first line of Table 4. Let the regular $`n`$-gon $`P`$ have unit edge length, and let the fold points $`x`$ be distance $`a`$ from $`v_0`$ along edge $`v_0v_1`$. Let $`b=12a`$. Call the point along $`P`$ to which $`v_i`$ glues $`v_i^{}`$. See Fig. 30 for an example with $`n=12`$. We will use this example throughout the section.
As mentioned in Section 1, we have no method for computing the 3D structure of the unique polytope determined by a particular Aleksandrov gluing. Moreover, we do not even have a general method for computing the creases, i.e., the edges of the polytope. We will therefore largely conjecture the structure of the pita polytopes in this section, although we will establish a subset of the creases. We will only explore the situation for even $`n`$. Let $`\alpha =2\pi /n`$, the turn angle at each vertex of the polygon.
We view each pita polytope as composed of four parts:<sup>7</sup><sup>7</sup>7 The relationship to the example in Fig. 26 should be evident.
1. A central parallelogram with short side $`a`$: $`(x,v_0^{},y,v_{n/2}^{})`$.
2. A top, nearly half-$`n`$-gon: $`(v_0^{},v_{n1}^{},v_{n2}^{},\mathrm{},v_{n/2+1}^{})`$.
3. A bottom, nearly half-$`n`$-gon, congruent by reflection to the top: $`(x,v_1^{},v_2^{},\mathrm{},v_{n/2}^{})`$.
4. A “mouth,” a strip of triangular teeth; see Fig. 31. $`n2`$ of the triangles in the strip are congruent; call their generic shape $`T_1`$. $`T_1`$ has sides of length $`b`$, $`2a`$, and $`1`$, with an angle $`\alpha `$ between $`b`$ and $`2a`$. The two extreme triangles of the mouth are smaller, of shape $`T_2`$: lengths $`b`$ and $`a`$ surrounding an angle $`\alpha `$.
We conjecture that the central parallelogram’s edges are creases, as is its central perimeter-splitting diagonal $`xy`$. Call the top and bottom nearly half-$`n`$-gons pita polygons. We have no conjectures about how the pita polygons are triangulated (except that they are triangulated the same). Finally, we prove below in Lemma 9.7 that the mouth is creased at the edges displayed in Fig. 31.
The final 3D shape looks something like Fig. 32. As $`n\mathrm{}`$, the polytope approaches a doubly-covered flat semicircle.
We now establish the structure of the mouth of pita polytopes. We start with this obvious claim:
###### Lemma 9.4
Pita polytopes are not flat.
Proof: A flat polytope is a pasting of two congruent polygons, oriented and aligned the same. The vertices of the polygons are the only spots on the polytope surface with curvature. We know the location of all these $`n+2`$ vertices: $`x`$, $`y`$, and $`v_0,\mathrm{},v_{n1}`$. Thus the two polygons must be $`(y,x,v_1,\mathrm{},v_{n/2})`$ $`(x,y,v_{n/2+1},\mathrm{},v_0)`$. However, because $`x`$ and $`y`$ are not at the midpoints of their edges (by definition of a pita polytope), these two polygons are not congruent. $`\mathrm{}`$
Our tools will be two facts about edges of triangulated polytopes, neither of which we will prove:
###### Fact 9.1
Every edge of a polytope is a shortest path between its endpoint vertices.
Call two polytope edges incident to the same polytope vertex $`v`$ adjacent if they are consecutive in a circular sorting around $`v`$.
###### Fact 9.2
The smaller surface angle between two adjacent edges incident to a polytope vertex is less than $`\pi `$. In other words, within every open semicircle of face angle at a polytope vertex $`v`$, there is at least one edge incident to $`v`$.
We use Fact 9.1 to eliminate certain geodesics as candidates for polytope edges. The following lemma gathers together some basic distance relationships to be used later to show that some geodesics are not shortest paths:
###### Lemma 9.5
The following distance relationships hold for the length of chords between points of a pita polygon:
1. $`|v_iv_j^{}|=|v_i^{}v_j|`$.
2. $`|v_i^{}v_j^{}|<|v_iv_j|`$ for all $`|ij|>1`$, i.e., for all $`ji`$ and $`ji\pm 1`$.
3. $`|v_i^{}x|<|v_ix|`$ for all $`i0`$.
4. $`|v_i^{}y|<|v_iy|`$ for all $`in/2`$.
Proof:
1. The polygons cut off by the chords $`(v_i,v_j^{})`$ and $`(v_i^{},v_j)`$ are congruent. For example, in Fig. 30, the chord $`(v_4,v_0^{})`$ cuts off a polygon of edge lengths $`(b,1,1,1)`$, and the chord $`(v_4^{},v_0)`$ cuts off a polygon of lengths $`(b,1,1,1)`$, both of whose outer interior angles are all $`\alpha `$.
2. Distances between the $`v_i^{}`$ vertices are in general less than distances between the corresponding unprimed vertices, because the primed vertices form a regular figure inscribed in the $`n`$-gon. A particular instance is illustrated in Fig. 33. For $`j=i+1`$, the distances are equal.
3. Here the reason is similar: the primed vertices are inscribed in the $`n`$-gon determined by the unprimed vertices. For example, $`|v_1x|=a+b`$, but $`|v_1^{}x|`$ is the length of the hypotenuse of a $`T_2`$ triangle, with sides $`a`$ and $`b`$, which is shorter by the triangle inequality. We will not detail the computations necessary to establish this claim for all $`i`$. The only exception to the inequality is for $`v_0`$, when $`|v_0x|=|v_0^{}x|=a`$.
4. Symmetric with previous case.
$`\mathrm{}`$
To eliminate the equal-length geodesics in Lemma 9.5(1), we will need the following:
###### Lemma 9.6
An edge $`e=vu`$ of a nonflat polytope $`Q`$ is a uniquely shortest path, i.e., there is not another geodesic of the same length from $`v`$ to $`u`$.
Proof: Suppose $`e=vu`$ is an edge of $`Q`$. Let $`g`$ be another geodesic between $`v`$ and $`u`$ of the same length as $`e`$. Then because $`e`$ is a straight segment in 3D, and because any nonstraight path is strictly longer, it must be that $`g`$ is also a straight segment in 3D. Thus it must be coincident with $`e`$. If $`e`$ and $`g`$ are nevertheless distinct, then they must be on opposite sides of a flat surface. But then $`Q`$ must be flat (“pinched”) at $`e=g`$, which by convexity implies that $`Q`$ is entirely flat. This contradicts the assumption of the lemma. $`\mathrm{}`$
We now have assembled enough information to pin down the structure of the mouth:
###### Lemma 9.7
The mouth of a pita polytope is triangulated as in Fig. 31: the edges
$$(x,v_1,\mathrm{},v_{n/2},y,v_{n/2+1},\mathrm{},v_{n1},v_0)$$
surrounding the mouth, and the diagonals $`(v_i,v_{ni})`$ and $`(v_{ni+1},v_i)`$ that delimit its “teeth” (cf. Fig. 31), are all polytope edges.
Proof: Let $`v_i`$, $`i\{1,\mathrm{},n/21\}`$ be a vertex of the mouth. (It may help to think of $`v_4`$ in Fig. 30 as a typical $`v_i`$ in this proof.) By Fact 9.2, there must be a polytope edge $`e`$ incident to $`v_i`$ on the top face in the half plane bounded by the line through $`v_{i1}v_i`$. By Lemma 9.5(2), the other endpoint of $`e`$ cannot be any $`v_j`$, $`|ij|>1`$, for all those are longer than $`|v_i^{}v_j^{}|`$, the length of an alternate geodesic. So they are not shortest paths, and are ruled out by Fact 9.1. By Lemma 9.5(3-4), the other endpoint of $`e`$ cannot be $`x`$ or $`y`$, for we have restricted $`i`$ so that $`i0`$ and $`in/2`$. This leaves $`v_j^{}`$ as a possible endpoint of $`e`$. But by Lemma 9.5(1), $`v_iv_j^{}`$ is not uniquely shortest, which by Lemma 9.6 then implies that $`Q`$ must be flat, which we know is false by Lemma 9.4.
We have excluded all candidates for the endpoint of $`e`$ except for $`j=i\pm 1`$. Because we are examining the semicircle bounded by $`v_{i1},v_i`$, this leaves $`v_{i+1}`$ as the only possible endpoint. Thus $`v_iv_{i+1}`$ is an edge of the polytope.
Repeating this argument for the bottom face, $`i\{n/2+1,\mathrm{},n1\}`$, establishes the outer boundary of the mouth, excluding the edges incident to $`x`$ and $`y`$. Those can be argued similarly. The teeth diagonals are now easy to see. We illustrate with $`v_4`$ in Fig. 30. We have just proved that no edge is incident to $`v_4`$ across the top face. But that top face must be triangulated somehow. The only way to triangulate it without using a diagonal incident to $`v_4`$ is to include the diagonal $`v_9^{}v_8^{}`$. This means that $`v_4v_8`$ and $`v_4v_9`$ are edges of the polytope. $`\mathrm{}`$
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# 1 Introduction
## 1 Introduction
The correspondence between a string theory on a Kähler manifold and an $`𝒩=2`$ Landau-Ginzburg theory is interesting and is very largely investigated . But, the most results are limited to the cases of compact Calabi-Yau manifolds. Recently, it is conjectured that in the case of a noncompact Calabi-Yau manifold, the associated CFT consists of an $`𝒩=2`$ Liouville theory and a Landau-Ginzburg theory . They claimed that when the Calabi-Yau $`n`$-fold $`X`$ is written as a hypersurface $`F(z_1,\mathrm{},z_{n+1})=0`$ in $`^{n+1}`$ by a quasi-homogeneous polynomial $`F`$, then the string theory on $`X`$ is equivalent to the CFT on
$`_\varphi \times S^1\times LG(W=F).`$
Here $`_\varphi `$ is a real line parametrized by $`\varphi `$ with linear dilaton background and $`LG(W=F)`$ is the IR theory of the Landau-Ginzburg model with the superpotential $`F`$. In this case, the background charge $`Q`$ of $`_\varphi `$ is determined by a condition of the total central charge. From the condition that $`Q`$ is a non-zero real number, we find that the base manifold $`X/^\times `$ should be curved positively.
The boson with linear dilaton background is strongly coupled in the region $`\varphi \mathrm{}`$, so we should introduce the Liouville potential or consider an $`SL(2)/U(1)`$ Kazama-Suzuki model to avoid the strong coupling singularity . But we do not care about this point in this paper.
In , they also claim that the string theory on this singular noncompact Calabi-Yau manifold $`X`$ is holographic dual to the “little string theory”.
In the case of a compact Calabi-Yau manifold, the string theory is “solved” in the description of Gepner model in a special point of the moduli space. We want to describe also the string theory on the noncompact Calabi-Yau manifold $`X`$ by the Gepner-like solvable model. If we can do it successfully, it will be possible to analyze more deeply a noncompact Calabi-Yau manifold and the little string theory.
In , they treat the string theories with ADE simple singularities. They construct the modular invariant partition functions, and show the consistency of these string theories.
In this paper, we consider more general cases, in which the Landau-Ginzburg part is described by a direct product of a number of minimal models. A typical example of ours is the Calabi-Yau $`n`$-fold $`X`$ described in the form
$`z_1^{N_1}+z_2^{N_2}+\mathrm{}+z_{n+1}^{N_{n+1}}=0,\text{ in }^{n+1}.`$
We construct the modular invariant partition functions and show the string theory actually exists consistently in these cases. We also calculate the elliptic genus, and find that it is factorised into two parts — a rather trivial one and a rather non-trivial one. We analyze the non-trivial one in detail, and find that it has the information on the cohomology of the positively curved base manifold $`X/^\times `$ except the elements generated by cup products of a Kähler form.
The organization of this paper is as follows. In the next section, we explain the setup and review the correspondence between a noncompact Calabi-Yau manifold and an $`𝒩=2`$ Liouville theory $`\times `$ Landau-Ginzburg theory. In section 3, we construct the modular invariant partition function. In section 4, we calculate the elliptic genus and compare it with the geometric property of the associated Calabi-Yau manifold $`X`$. In the last section, we summarize the results and discuss the problems and prospects. In Appendix A. we collect some useful equations of theta functions and characters that we use in this paper.
## 2 The string theory on a noncompact singular Calabi-Yau manifold
We consider the string compactification to a noncompact, singular Calabi-Yau $`n`$-fold $`X`$. The total target space is expressed by a direct product of a $`d`$ dimensional flat spacetime and the manifold $`X`$
$`^{d1,1}\times X.`$ (2.1)
Here, $`n`$ is related to $`d`$ by the constraint on the total dimension $`2n+d=10`$.
For simplicity, we concentrate the case that the noncompact singular Calabi-Yau manifold $`X`$ is realized as the hypersurface in $`^{n+1}`$ determined by the algebraic equation with a quasi homogeneous polynomial $`F`$
$`F(z_1,\mathrm{},z_{n+1})=0.`$
By the term “quasi-homogeneous”, we mean that the polynomial $`F`$ satisfies
$`F(\lambda ^{r_1}z_1,\mathrm{},\lambda ^{r_{n+1}}z_{n+1})=\lambda F(z_1,\mathrm{},z_{n+1}),`$ (2.2)
for some exponents $`\{r_j\}`$ and for an arbitrary $`\lambda ^\times `$.
This manifold $`X`$ is singular at $`(z_1,\mathrm{},z_{n+1})=(0,\mathrm{},0)`$. If we consider the manifold $`(X(0,\mathrm{},0))/^\times `$, where the action of $`^\times `$ is $`(z_1,\mathrm{},z_{n+1})(\lambda ^{r_1}z_1,\mathrm{},\lambda ^{r_{n+1}}z_{n+1})`$ with the exponents $`\{r_j\}`$ of (2.2), then we get a compact manifold. We denote this compact manifold simply as $`X/^\times `$ and call it “the base manifold of $`X`$”.
It is conjectured in that the string theory on the space (2.1) is equivalent to the theory including flat spacetime $`^{d1,1}`$, a line with the linear dilation background $`_\varphi `$, $`S^1`$, and the Landau-Ginzburg theory with a superpotential $`W=F`$;
$`^{d1,1}\times _\varphi \times S^1\times LG(W=F).`$
The part $`(_\varphi \times S^1)`$ has a world sheet $`𝒩=2`$ superconformal symmetry. Let $`\varphi `$ be the parameter of $`_\varphi `$ , $`Y`$ be the parameter of $`S^1`$, and $`\psi ^+,\psi ^{}`$ be the fermionic part of $`(_\varphi \times S^1)`$. The $`𝒩=2`$ superconformal currents are written in terms of the above fields
$`T={\displaystyle \frac{1}{2}}(Y)^2{\displaystyle \frac{1}{2}}(\varphi )^2{\displaystyle \frac{Q}{2}}^2\varphi {\displaystyle \frac{1}{2}}(\psi ^+\psi ^{}\psi ^+\psi ^{}),`$
$`G^\pm ={\displaystyle \frac{1}{\sqrt{2}}}\psi ^\pm (iY\pm \varphi ){\displaystyle \frac{Q}{\sqrt{2}}}\psi ^\pm ,`$
$`J=\psi ^+\psi ^{}QiY.`$ (2.3)
The associated central charge of this algebra is $`\widehat{c}(=c/3)=1+Q^2`$.
In this paper, we consider the case in which the Landau-Ginzburg theory with superpotential $`W=F`$ can be described by a direct product of $`𝒩=2`$ minimal models. Let $`M_{G,N}`$ be the minimal model corresponding to simply laced Lie algebra $`G=A,D,E`$ with dual Coxeter number $`N`$. We consider the theory in the following;
$`^{d1,1}\times _\varphi \times S^1\times M_{G_1,N_1}\times \mathrm{}\times M_{G_R,N_R},`$
where $`R`$ is the number of the minimal models. The cases with $`R=1`$ are treated in and $`R=0`$ in
A typical example is the case that all $`G_j`$ are $`A`$ type. In this example, the quasi-homogeneous polynomial is written as
$`F(z)=z_1^{N_1}+\mathrm{}+z_R^{N_R}+z_{R+1}^2+\mathrm{}+z_{n+1}^2.`$
The background charge of $`_\varphi `$ is determined by the criticality condition. To cancel the conformal anomaly, the total central charge is to be $`0`$. The central charge of the ghost sector is $`15`$, so the total central charge of the matter sector is to be $`15`$. The central charge of the flat spacetime is $`3/2`$ for each pair of a boson and a fermion, and that of the $`_\varphi \times S^1`$ is $`3+3Q^2`$ as mentioned above, and that of a minimal model $`M_{G,N}`$ is $`\frac{3(N2)}{N}`$. Therefore, the criticality condition leads to the equation
$`{\displaystyle \frac{3d}{2}}+3+3Q^2+{\displaystyle \underset{j=1}{\overset{R}{}}}{\displaystyle \frac{3(N_j2)}{N_j}}=15.`$
From this criticality condition, we obtain the value of $`Q^2`$ as
$`Q^2=4{\displaystyle \frac{d}{2}}{\displaystyle \underset{j}{}}{\displaystyle \frac{(N_j2)}{N_j}}.`$ (2.4)
By the condition $`Q^2>0`$ for a real number $`Q`$, the right-hand side should be positive;
$`4{\displaystyle \frac{d}{2}}{\displaystyle \underset{j}{}}{\displaystyle \frac{(N_j2)}{N_j}}>0.`$ (2.5)
It is equivalent to a condition that the singularity is in finite distance in the moduli space of deformation of singular Calabi-Yau manifold $`X`$ . In view of the base manifold $`X/^\times `$, the finite distance condition is equivalent to that $`X/^\times `$ is positively curved.
## 3 Modular invariant partition function
Now, let us construct the modular invariant partition function. We take the light-cone gauge, then the associated CFT to consider is
$`^{d2}\times _\varphi \times S^1\times M_{G_1,N_1}\times \mathrm{}\times M_{G_R,N_R}.`$
The toroidal partition function can be separated into 2 parts: the one $`Z_{GSO}`$ concerning to the GSO projection and the other $`Z_0`$ not concerning to it. We construct the total partition function $`Z`$ as
$`Z={\displaystyle \frac{d^2\tau }{\tau _2^2}Z_0(\tau ,\overline{\tau })Z_{GSO}(\tau ,\overline{\tau })},`$
where $`\tau =\tau _1+i\tau _2`$ is the moduli parameter of the torus, and $`d^2\tau /\tau _2^2`$ is the modular invariant measure.
First, we study the rather easy part $`Z_0`$, then we investigate the rather complicated part $`Z_{GSO}`$.
### 3.1 GSO independent part of the partition function
In this subsection, we discuss the $`Z_0`$ : the GSO independent part. It is completely the same as that in .
The $`Z_0`$ includes the contribution from the flat spacetime bosonic coordinates $`X^I,(I=2,\mathrm{},d1)`$ and the linear dilation $`\varphi `$.
The partition function of each flat spacetime boson is represented by the Dedekind eta function $`\eta (\tau )`$ as
$`{\displaystyle \frac{1}{\sqrt{\tau _2}|\eta (\tau )|^2}}.`$
The partition function of $`\varphi `$ is defined as $`Z_L=\mathrm{Tr}q^{L_0c_L/24}\overline{q}^{\overline{L}_0c_L/24},(q=\mathrm{exp}(2\pi i\tau ),c_L=1+3Q^2)`$ in the canonical formalism. Here the trace is taken over delta function normalizable primary fields
$`\mathrm{exp}(ip\varphi ),p={\displaystyle \frac{iQ}{2}}+\mathrm{},\mathrm{},`$ (3.1)
and their excitations by oscillators. Then we obtain $`Z_L`$ as
$`Z_L`$ $`=`$ $`{\displaystyle \frac{1}{|_{n=1}(1q^n)|^2}}{\displaystyle 𝑑p\mathrm{exp}\left[4\pi \tau _2\left(\frac{1}{2}p^2+\frac{i}{2}pQ\frac{1+3Q^2}{24}\right)\right]}`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\tau _2}|\eta (\tau )|^2}},`$
where the region of the integral of $`p`$ is as (3.1). As a result, the partition function of $`\varphi `$ is the same as that of an ordinary boson. So we obtain $`Z_0`$ as the partition function of effectively $`(d1)`$ free bosons;
$`Z_0=\left({\displaystyle \frac{1}{\sqrt{\tau _2}|\eta (\tau )|^2}}\right)^{d1}.`$
The primary fields of (3.1) correspond to “Principal continuous series” in terms of the representation of $`SL(2)`$. To include the other sectors is an interesting problem and is postponed as a future work.
### 3.2 GSO dependent part of the partition function : $`d=2,6`$ cases
Now, let us proceed to the GSO dependent part $`Z_{GSO}`$. In this subsection, we treat $`d=2,6`$ cases.
This part includes $`(d2)`$ flat spacetime fermions $`\psi ^I,(I=2,\mathrm{},d1)`$ , two free fermions $`\psi ^\varphi ,\psi ^Y`$ associated to $`_\varphi \times S^1`$, minimal models $`M_{G_1,N_1},\mathrm{},M_{G_R,N_R}`$ , and an $`S^1`$ boson $`Y`$. We combine the $`d`$ free fermions $`\psi ^I,(I=2,\mathrm{},d1),\psi ^\varphi `$ and $`\psi ^Y`$ and construct the affine Lie algebra $`\widehat{SO(d)}_1`$. The Verma module of $`\widehat{SO(d)}_1`$ is characterized by an integer $`s_0=0,1,2,3`$ , which labels the representations of $`SO(d)`$, that is scalar, spinor, vector, and cospinor, respectively.
Let us turn to Verma modules of $`𝒩=2`$ minimal models. The Verma module of an $`𝒩=2`$ minimal model is specified with three indices $`(\mathrm{},m,s)`$, which satisfy the following conditions
$`\mathrm{}=0,\mathrm{},N2,`$
$`m=0,\mathrm{},2N1,`$
$`s=0,1,2,3,`$
$`\mathrm{}+m+s0mod2.`$ (3.2)
We denote $`\chi _m^{\mathrm{},s}(\tau ,z)`$ as the character of the Verma module labeled by the set $`(\mathrm{},m,s)`$. Some properties of this character is collected in Appendix A.
The Verma module of the whole GSO dependent parts is specified by the index $`s_0`$ of $`\widehat{SO(d)}_1`$ representation, the indices $`(\mathrm{}_j,m_j,s_j),(j=1,\mathrm{},R)`$ of the minimal models, and the $`S^1`$ momentum $`p`$. We combine these indices except $`p`$ into two vectors $`\lambda ,\mu `$.
$`\lambda :=(\mathrm{}_1,\mathrm{},\mathrm{}_R),`$
$`\mu :=(s_0;s_1,\mathrm{},s_R;m_1,\mathrm{},m_R).`$
We shall introduce the inner product between $`\mu `$ and $`\mu ^{}`$ as in
$`\mu \mu ^{}:={\displaystyle \frac{d}{2}}{\displaystyle \frac{s_0s_0^{}}{4}}{\displaystyle \underset{j=1}{\overset{R}{}}}{\displaystyle \frac{s_js_j^{}}{4}}+{\displaystyle \underset{j=1}{\overset{R}{}}}{\displaystyle \frac{m_jm_j^{}}{2N_j}}.`$
Also it is convenient to introduce special vectors $`\beta _0,\beta _j(j=1,\mathrm{},R)`$
$`\beta _0:=(1;1,\mathrm{},1;1,\mathrm{},1),`$
$`\beta _j:=(2;0,\mathrm{},0,\underset{\underset{S_j}{}}{2},0,\mathrm{}0;0,\mathrm{},0).`$
Here the $`\beta _0`$ is the vector with all components $`1`$, and $`\beta _j`$ is the vector with $`s_0`$ and $`s_j`$ components $`2`$ and the others zero.
With these notations, the criticality condition (2.4) can be written in a rather simple form as
$`Q^2=4(1+\beta _0\beta _0).`$ (3.3)
When we define an integer $`K:=\mathrm{lcm}(2,N_j)`$, $`KQ^2`$ is shown to be an even integer because of the equations (2.4) and (3.3). Therefore, it is convenient to define an integer $`J`$ by the equation
$`J:=2K(1+\beta _0\beta _0)(=KQ^2/2).`$ (3.4)
In terms of $`J`$, the finite distance condition (2.5) can be expressed as $`J>0`$.
Now, let us consider the character of the Verma module $`(\lambda ,\mu ,p)`$,
$`\chi _\mu ^\lambda (\tau ){\displaystyle \frac{q^{\frac{1}{2}p^2}}{\eta (\tau )}},`$
where $`\chi _\mu ^\lambda (\tau )`$ is the product of characters of the minimal models and the $`\widehat{SO(d)}_1`$ character $`\chi _{s_0}(\tau )`$ of the $`s_0`$ representation
$`\chi _\mu ^\lambda (\tau ):=\chi _{s_0}(\tau )\chi _{m_1}^{\mathrm{}_1,s_1}(\tau )\mathrm{}\chi _{m_R}^{\mathrm{}_R,s_R}(\tau ).`$
In this character, $`\chi _\mu ^\lambda (\tau )`$ has good modular properties, but $`q^{\frac{1}{2}p^2}/\eta (\tau )`$ has bad ones. So, we will sum up the characters with respect to certain values of $`p`$ and make the modular properties good .
Let us consider the GSO projection. By the GSO projection, we pick up the states with odd integral $`U(1)`$ charges of the $`𝒩=2`$ superconformal symmetry. The $`U(1)`$ charge of the states in the above Verma module is expressed as
$`2\beta _0\mu +pQ={\displaystyle \frac{d}{2}}{\displaystyle \frac{s_0}{2}}{\displaystyle \underset{j}{}}{\displaystyle \frac{s_j}{2}}+{\displaystyle \underset{j}{}}{\displaystyle \frac{m_j}{N_j}}+pQ.`$
From the condition that this $`U(1)`$ charge must be an odd integer $`(2u+1)`$ with $`u`$, the $`S^1`$ momentum $`p`$ is written as
$`p(u)={\displaystyle \frac{1}{Q}}\left(2u+12\beta _0\mu \right).`$
If we sum up the characters for all $`u`$, we obtain the theta function with a fractional level, which does not have good modular properties. So we perform the following trick.
Let us write $`u=Jv+w`$ with integers $`v,w`$ and sum up the characters for $`v`$. Then the sum leads to the following theta function
$`{\displaystyle \underset{v}{}}q^{\frac{1}{2}p(u=Jv+w)^2}=\mathrm{\Theta }_{2K\beta _0\mu +K(2w+1),KJ}(\tau ).`$ (3.5)
Note that $`2K\beta _0\mu +K(2w+1)`$ is an integer, and the above theta function has good modular properties.
Now, including oscillator modes and other sectors, we can define the building blocks $`f_{\underset{\stackrel{~}{}}{\mu }}^\lambda (\tau )`$ by
$`f_{\underset{\stackrel{~}{}}{\mu }}^\lambda (\tau ):=\chi _\mu ^\lambda (\tau )\mathrm{\Theta }_{M,KJ}(\tau )/\eta (\tau ),`$
$`\underset{\stackrel{~}{}}{\mu }:=(\mu ,M),M_{2KJ}.`$
We should use only the building blocks $`f_{\underset{\stackrel{~}{}}{\mu }}^\lambda `$ with the conditions
$`M=2K\beta _0\mu +K(2w+1)\text{ for }{}_{}{}^{}w,`$ (3.6)
$`s_0s_1\mathrm{}s_Rmod2.`$ (3.7)
The condition (3.6) comes from the formula of (3.5), and the condition (3.7) implies that the boundary condition of the fermionic currents are the same in all the sub-theories, i.e. they must be all in the NS sector, or all in the R sector.
The modular invariant partition function can be systematically obtained by “the beta method”.
The inner product between of two vectors $`\underset{\stackrel{~}{}}{\mu },\underset{\stackrel{~}{}}{\mu }^{}`$ is defined as
$`\underset{\stackrel{~}{}}{\mu }\underset{\stackrel{~}{}}{\mu }^{}:=\mu \mu ^{}{\displaystyle \frac{MM^{}}{2KJ}}.`$
We also extend the vectors $`\beta _0,\beta _j`$ to $`\underset{\stackrel{~}{}}{\beta _0},\underset{\stackrel{~}{}}{\beta _j}`$ as
$`\underset{\stackrel{~}{}}{\beta _0}:=(\beta _0,J),`$
$`\underset{\stackrel{~}{}}{\beta _j}:=(\beta _j,0),`$
and evaluate the inner products of these $`\underset{\stackrel{~}{}}{\beta _0},\underset{\stackrel{~}{}}{\beta _j}`$ vectors
$`\underset{\stackrel{~}{}}{\beta _0}\underset{\stackrel{~}{}}{\beta _0}=\beta _0\beta _0{\displaystyle \frac{J^2}{2KJ}}=1,`$
$`\underset{\stackrel{~}{}}{\beta _j}\underset{\stackrel{~}{}}{\beta _j}=\beta _j\beta _j={\displaystyle \frac{d}{2}}1,`$
$`\underset{\stackrel{~}{}}{\beta _j}\underset{\stackrel{~}{}}{\beta _0}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{d}{2}}1\right).`$ (3.8)
Note that $`\underset{\stackrel{~}{}}{\beta _0}\underset{\stackrel{~}{}}{\beta _0}`$ is an odd integer, $`\underset{\stackrel{~}{}}{\beta _j}\underset{\stackrel{~}{}}{\beta _j}`$ are even integers (Recall that we consider the cases $`d=2,6`$), and $`\underset{\stackrel{~}{}}{\beta _j}\underset{\stackrel{~}{}}{\beta _0}`$ are integers. Using these special vectors, the conditions (3.6) and (3.7) are written in a simple form
$`2\underset{\stackrel{~}{}}{\beta _0}\underset{\stackrel{~}{}}{\mu }2+1,`$
$`\underset{\stackrel{~}{}}{\beta _j}\underset{\stackrel{~}{}}{\mu }.`$ (3.9)
We call this condition “the beta condition”.
Using these notations, and the modular transformation laws of theta functions and $`𝒩=2`$ characters written in the Appendix A, we can calculate the modular transformation laws of $`f_{\underset{\stackrel{~}{}}{\mu }}^\lambda `$ as
$`f_{\underset{\stackrel{~}{}}{\mu }}^\lambda (\tau +1)=𝐞\left[{\displaystyle \underset{j}{}}{\displaystyle \frac{\mathrm{}_j(\mathrm{}_j+1)}{4N_j}}{\displaystyle \frac{1}{2}}\underset{\stackrel{~}{}}{\mu }\underset{\stackrel{~}{}}{\mu }{\displaystyle \frac{1}{24}}\left({\displaystyle \underset{j}{}}{\displaystyle \frac{N_j2}{N_j}}+{\displaystyle \frac{d}{2}}+1\right)\right]f_{\underset{\stackrel{~}{}}{\mu }}^\lambda (\tau ),`$
$`f_{\underset{\stackrel{~}{}}{\mu }}^\lambda (1/\tau )={\displaystyle \underset{\lambda ^{},\underset{\stackrel{~}{}}{\mu }^{}}{\overset{\mathrm{even}}{}}}A_{\lambda \lambda ^{}}\left({\displaystyle \underset{j}{}}{\displaystyle \frac{1}{\sqrt{8N_j}}}\right){\displaystyle \frac{1}{\sqrt{8KJ}}}𝐞\left[\underset{\stackrel{~}{}}{\mu }\underset{\stackrel{~}{}}{\mu }^{}\right]f_{\underset{\stackrel{~}{}}{\mu }^{}}^\lambda ^{}(\tau ),`$
where the sums of the $`\lambda ^{},\mu ^{}`$ are taken only for the range (3.2) and for $`M=0,\mathrm{},2KJ1`$. Especially we must impose the condition $`\mathrm{}_j+m_j+s_j0mod2`$ for each minimal model. $`A_{\lambda \lambda ^{}}`$ is the products of the $`\widehat{SU(2)}_{N_j2}`$ S matrices $`A_{\mathrm{}_j\mathrm{}_j^{}}^{(N_j)}`$;
$`A_{\lambda \lambda ^{}}={\displaystyle \underset{j}{}}A_{\mathrm{}_j\mathrm{}_j^{}}^{(N_j)}={\displaystyle \underset{j}{}}\sqrt{{\displaystyle \frac{2}{N_j}}}\mathrm{sin}\pi {\displaystyle \frac{(\mathrm{}_j+1)(\mathrm{}_j^{}+1)}{N_j}},`$
and we use here and the rest of this paper the notation $`𝐞\left[x\right]=\mathrm{exp}(2\pi ix)`$.
Let us note that if a vector $`\underset{\stackrel{~}{}}{\mu }`$ satisfies the beta condition (3.9), the vector $`\underset{\stackrel{~}{}}{\mu }+b_0\underset{\stackrel{~}{}}{\beta _0}+_jb_j\underset{\stackrel{~}{}}{\beta _j}`$ for $`b_0,b_j,(j=1,\mathrm{},R)`$ also satisfies the beta condition by virtue of (3.8). Using this fact, we define the function $`F_{\underset{\stackrel{~}{}}{\mu }}^\lambda `$ for $`(\lambda ,\underset{\stackrel{~}{}}{\mu })`$ which satisfies the beta conditions (3.9) as a sum of $`f_{\underset{\stackrel{~}{}}{\mu }+b_0\underset{\stackrel{~}{}}{\beta _0}+_jb_j\underset{\stackrel{~}{}}{\beta _j}}^\lambda `$’s as
$`F_\mu ^\lambda (\tau )={\displaystyle \underset{b_0,b_j}{}}(1)^{s_0+b_0}f_{\underset{\stackrel{~}{}}{\mu }+b_0\underset{\stackrel{~}{}}{\beta _0}+_jb_j\underset{\stackrel{~}{}}{\beta _j}}^\lambda (\tau ),`$
where the sum is taken for $`b_0_{2K}`$ and $`b_j_2`$. The sign $`(1)^{s_0+b_0}`$ is $`(1)`$ for the Ramond sector.
These functions have very good modular properties. Especially by S transformation, the functions are mixed among those which satisfy the beta condition;
$`F_{\underset{\stackrel{~}{}}{\mu }}^\lambda (\tau +1)=𝐞\left[{\displaystyle \underset{j}{}}{\displaystyle \frac{\mathrm{}_j(\mathrm{}_j+1)}{4N_j}}{\displaystyle \frac{1}{2}}\underset{\stackrel{~}{}}{\mu }\underset{\stackrel{~}{}}{\mu }{\displaystyle \frac{1}{24}}\left({\displaystyle \underset{j}{}}{\displaystyle \frac{N_j2}{N_j}}+{\displaystyle \frac{d}{2}}+1\right)\right]F_{\underset{\stackrel{~}{}}{\mu }}^\lambda (\tau ),`$
$`F_{\underset{\stackrel{~}{}}{\mu }}^\lambda (1/\tau )={\displaystyle \underset{\lambda ^{},\underset{\stackrel{~}{}}{\mu }^{}}{\overset{\mathrm{even},\mathrm{beta}}{}}}A_{\lambda \lambda ^{}}\left({\displaystyle \underset{j}{}}{\displaystyle \frac{1}{\sqrt{8N_j}}}\right){\displaystyle \frac{1}{\sqrt{8KJ}}}𝐞\left[\underset{\stackrel{~}{}}{\mu }\underset{\stackrel{~}{}}{\mu }^{}\right](1)^{s_0+s_0^{}}F_{\underset{\stackrel{~}{}}{\mu }^{}}^\lambda ^{}(\tau ),`$
where the sums of $`\lambda ^{},\underset{\stackrel{~}{}}{\mu }^{}`$ is taken for restricted subclass that satisfies the conditions (3.2) and the beta condition (3.9).
With this function $`F_\mu ^\lambda `$, we obtain the modular invariant $`Z_{GSO}`$ as
$`Z_{GSO}(\tau ,\overline{\tau })={\displaystyle \frac{1}{4^R\times 2K}}{\displaystyle \underset{\lambda ,\overline{\lambda },\underset{\stackrel{~}{}}{\mu }}{\overset{\mathrm{even},\mathrm{beta}}{}}}L_{\lambda \overline{\lambda }}F_{\underset{\stackrel{~}{}}{\mu }}^\lambda (\tau )\overline{F}_{\underset{\stackrel{~}{}}{\mu }}^{\overline{\lambda }}(\overline{\tau }),`$
where $`L_{\lambda \overline{\lambda }}=_jL_{\mathrm{}_j\overline{\mathrm{}}_j}^{(G_j,N_j2)}`$ is the product of $`G_j=A,D,E`$ type modular invariants of $`\widehat{SU(2)}_{N_j2}`$ .
We can check the modular invariance of the above partition function.
We expect from spacetime supersymmetry that the partition function vanishes, or equivalently $`F_{\underset{\stackrel{~}{}}{\mu }}^\lambda (\tau )=0`$. It is a future work to check that it actually vanishes.
Here, we find a solution, but it may not be the only solution, and there can be some variety of modular invariant partition function. Actually, for $`R=1`$ case, there are many other solutions associated with the other modular invariants of the theta system .
### 3.3 GSO dependent part of the partition function : $`d=4`$ case
In this subsection, we comment on the $`d=4`$ case. To construct the modular invariant partition function in the $`d=4`$ case, we combine the four fermions to construct the affine currents $`\widehat{SO(2)}_1\times \widehat{SO(2)}_1`$ and label the the Verma module by indices $`s_1`$ and $`s_0`$. Then, the modular invariant partition function can be constructed in almost the same way as the $`d=2,6`$ cases.
First we define the vectors $`\underset{\stackrel{~}{}}{\mu }`$’s and the inner product between them as
$`\underset{\stackrel{~}{}}{\mu }:=(s_1,s_0;s_1,\mathrm{},s_R;m_1,\mathrm{},m_R;M),`$
$`\underset{\stackrel{~}{}}{\mu }\underset{\stackrel{~}{}}{\mu }^{}:={\displaystyle \frac{s_1s_1^{}}{4}}{\displaystyle \frac{s_0s_0^{}}{4}}{\displaystyle \underset{j}{}}{\displaystyle \frac{s_js_j^{}}{4}}+{\displaystyle \underset{j}{}}{\displaystyle \frac{m_jm_j^{}}{2N_j}}{\displaystyle \frac{MM^{}}{2KJ}}.`$
It is convenient to introduce special vectors $`\underset{\stackrel{~}{}}{\beta _0}`$, $`\underset{\stackrel{~}{}}{\beta _j}`$ and $`\underset{\stackrel{~}{}}{\beta _1}`$
$`\underset{\stackrel{~}{}}{\beta _0}=(1,1;1,\mathrm{},1;1,\mathrm{},1;M),`$
$`\underset{\stackrel{~}{}}{\beta _j}=(0,2;0,\mathrm{},0,\underset{\underset{S_j}{}}{2},0,\mathrm{},0;0,\mathrm{},0;0),(j=1,\mathrm{},R),`$
$`\underset{\stackrel{~}{}}{\beta _1}=(2,2;0,\mathrm{},0;0,\mathrm{},0;0).`$
Using these vectors, we can construct the building blocks $`f_{\underset{\stackrel{~}{}}{\mu }}^\lambda (\tau )`$ as
$`f_{\underset{\stackrel{~}{}}{\mu }}^\lambda (\tau ):=\chi _{s_1}(\tau )\chi _{s_0}(\tau )\chi _m^{\mathrm{}_1,s_1}(\tau )\mathrm{}\chi _m^{\mathrm{}_R,s_R}(\tau )\mathrm{\Theta }_{M,KJ}(\tau )/\eta (\tau ),`$
where $`\chi _{s_1}(\tau )`$ and $`\chi _{s_0}(\tau )`$ are $`\widehat{SO(2)}_1`$ characters. Then the GSO conditions and the condition of fermionic sectors are
$`2\underset{\stackrel{~}{}}{\beta _0}\underset{\stackrel{~}{}}{\mu }2+1,\underset{\stackrel{~}{}}{\beta _j}\underset{\stackrel{~}{}}{\mu },\underset{\stackrel{~}{}}{\beta _1}\underset{\stackrel{~}{}}{\mu },`$ (3.10)
and we can construct the modular invariant partition function by the beta method in this case. Next we introduce the function $`F_{\underset{\stackrel{~}{}}{\mu }}^\lambda (\tau )`$ as
$`F_{\underset{\stackrel{~}{}}{\mu }}^\lambda (\tau )={\displaystyle \underset{b_0_{2K},b_j_2,b_1_2}{}}(1)^{b_0+s_0}f_{\underset{\stackrel{~}{}}{\mu }+b_0\underset{\stackrel{~}{}}{\beta _0}+_jb_j\underset{\stackrel{~}{}}{\beta _j}+b_1\underset{\stackrel{~}{}}{\beta _1}}^\lambda (\tau ),`$
then we obtain the GSO dependent part of the modular invariant partition function $`Z_{GSO}`$
$`Z_{GSO}(\tau ,\overline{\tau })={\displaystyle \frac{1}{4^R\times 4K}}{\displaystyle \underset{\lambda ,\overline{\lambda },\underset{\stackrel{~}{}}{\mu }}{\overset{\mathrm{even},\mathrm{beta}}{}}}L_{\lambda \overline{\lambda }}F_{\underset{\stackrel{~}{}}{\mu }}^\lambda (\tau )\overline{F}_{\underset{\stackrel{~}{}}{\mu }}^{\overline{\lambda }}(\overline{\tau }).`$
We can check the modular invariance of the above partition function.
## 4 Elliptic genus
In this section, we calculate the elliptic genus of the theory . The definition of the elliptic genus is
$`Z(\tau ,\overline{\tau },z):=\mathrm{Tr}_{RR}(1)^Fq^{L_0c/24}\overline{q}^{\overline{L}_0c/24}y^{J_0},`$
where the trace is taken for the RR states, and $`y=\mathrm{exp}(2\pi iz)`$. This elliptic genus has the following modular properties;
$`Z(\tau +1,\overline{\tau }+1,z)=Z(\tau ,\overline{\tau },z)=Z(\tau ,\overline{\tau },z),`$
$`Z(1/\tau ,1/\overline{\tau },z/\tau )=𝐞\left[{\displaystyle \frac{\widehat{c}}{2}}{\displaystyle \frac{z^2}{\tau }}\right]Z(\tau ,\overline{\tau },z).`$ (4.1)
Here, we omit the contribution from the flat space time and consider only the internal part describing the Calabi-Yau $`n`$-fold $`X`$. We calculate its elliptic genus and the Witten index, and discuss its geometrical interpretation.
Let us consider again “the criticality condition”, in other words “the Calabi-Yau condition” of $`X`$. Here the total $`\widehat{c}`$ should be $`n`$ because we want the theory that describes a Calabi-Yau $`n`$-fold. Therefore, the total $`\widehat{c}`$ of the $`𝒩=2`$ Liouville and the minimal models should satisfy the relations
$`n=\widehat{c}=1+Q^2+{\displaystyle \underset{j}{}}{\displaystyle \frac{N_j2}{N_j}}.`$ (4.2)
We introduce the following vectors with $`R`$ components $`\{m_j\}`$, and the inner product between them as
$`\nu :=(m_1,\mathrm{},m_R),`$
$`\nu \nu ^{}:={\displaystyle \underset{j}{}}{\displaystyle \frac{m_jm_j^{}}{2N_j}}.`$
We also introduce the special vector $`\gamma _0`$ with all components $`2`$
$`\gamma _0:=(2,\mathrm{},2).`$
With these notations, the condition (4.2) becomes
$`Q^2=n1R+\gamma _0\gamma _0.`$
Next we let $`N:=\mathrm{lcm}(N_j)`$, and define $`J`$ as
$`{\displaystyle \frac{2J}{N}}:=Q^2.`$ (4.3)
In this paper, we concentrate only the case that $`(n1R)`$ is even, then in this case, $`J`$ is an integer. In terms of $`J`$, the finite distance condition $`Q^2>0`$ can be written as $`J>0`$.
Because we want a Calabi-Yau CFT, we have to pick up only the states with integral $`U(1)`$ charges. This condition is realized as the condition
$`\gamma _0\nu +pQ.`$ (4.4)
From this, $`p`$ can be written with an arbitrary integer $`u`$
$`p={\displaystyle \frac{1}{Q}}\left(u\gamma _0\nu \right).`$
Following the same manner as in the previous section, we let $`u=2Jv+w`$ and sum up for $`v`$. It leads to the theta function
$`{\displaystyle \underset{v}{}}q^{\frac{1}{2}p^2}y^{pQ}=\mathrm{\Theta }_{N(w\gamma _0\nu ),NJ}(\tau ,2z/N).`$
Note that $`N(w\gamma _0\nu )`$ is an integer and $`\mathrm{\Theta }_{N(w\gamma _0\nu ),NJ}(\tau ,2z/N)`$ has good modular properties.
Collecting these, we define the building blocks $`g_{\underset{\stackrel{~}{}}{\nu }}^\lambda `$ as
$`g_{\underset{\stackrel{~}{}}{\nu }}^\lambda (\tau ,z):={\displaystyle \underset{s_0,s_j=1,3}{}}\chi _\mu ^\lambda (\tau ,z){\displaystyle \frac{\mathrm{\Theta }_{M,NJ}(\tau ,2z/N)}{\eta (\tau )}}(1)^{\frac{s_0}{2}_j\frac{s_j}{2}+\gamma _0\nu +\frac{M}{N}},`$
where $`\underset{\stackrel{~}{}}{\nu }:=(\nu ,M)`$. In the sign $`(1)^{J_0}=(1)^{\frac{s_0}{2}_j\frac{s_j}{2}+\gamma _0\nu +\frac{M}{N}}`$ , the part $`(\frac{s_0}{2}_j\frac{s_j}{2})`$ represents contributions of ordinary $`U(1)`$ charges from the indices $`s_0,s_j`$, and the rest $`\gamma _0\nu +\frac{M}{N}=w=u2Jv`$ reflects contributions from the indices $`m_j`$ and $`S^1`$ momentum.
Let us define the inner product between $`\underset{\stackrel{~}{}}{\nu }`$ and $`\underset{\stackrel{~}{}}{\nu }^{}`$ as
$`\underset{\stackrel{~}{}}{\nu }\underset{\stackrel{~}{}}{\nu }^{}:=\nu \nu ^{}{\displaystyle \frac{MM^{}}{2NJ}},`$
and the special vector
$`\underset{\stackrel{~}{}}{\gamma _0}=(\gamma _0,2J).`$
We also introduce the functions $`I_m^{\mathrm{}}`$ and $`I_\nu ^\lambda `$
$`I_m^{\mathrm{}}(\tau ,z):=\chi _m^{\mathrm{},1}(\tau ,z)\chi _m^{\mathrm{},3}(\tau ,z),`$
$`I_\nu ^\lambda (\tau ,z):=I_{m_1}^\mathrm{}_1(\tau ,z)\mathrm{}I_{m_R}^\mathrm{}_R(\tau ,z).`$
With these notations, the building blocks $`g_{\underset{\stackrel{~}{}}{\nu }}^\lambda `$ can be written as
$`g_{\underset{\stackrel{~}{}}{\nu }}^\lambda (\tau ,z)={\displaystyle \frac{\theta _1(\tau ,z)}{\eta (\tau )}}I_\nu ^\lambda (\tau ,z){\displaystyle \frac{\mathrm{\Theta }_{M,NJ}(\tau ,2z/N)}{\eta (\tau )}}(1)^{\underset{\stackrel{~}{}}{\gamma _0}\underset{\stackrel{~}{}}{\nu }},`$
where we omit the overall irrelevant phase. The condition (4.4) can be rewritten as
$`\underset{\stackrel{~}{}}{\gamma _0}\underset{\stackrel{~}{}}{\nu },`$ (4.5)
and again we call this condition “the beta condition”.
Now, we construct elliptic genus using the above building blocks $`g_{\underset{\stackrel{~}{}}{\nu }}^\lambda `$ which satisfy the condition (4.5).
Note that if $`\underset{\stackrel{~}{}}{\nu }`$ satisfies the beta condition, then $`\underset{\stackrel{~}{}}{\nu }+b_0\underset{\stackrel{~}{}}{\gamma _0}`$ for $`b_0`$ also satisfies the beta condition, because
$`\underset{\stackrel{~}{}}{\gamma _0}\underset{\stackrel{~}{}}{\gamma _0}=\gamma _0\gamma _0{\displaystyle \frac{2J}{N}}=(n1R),`$
is an integer. <sup>1</sup><sup>1</sup>1 Actually, it is an even integer. Remember that we concentrate the case in which $`(n1R)`$ is an even integer. Here we used the definition of $`J`$ (4.3). So let us define the new functions $`G_{\underset{\stackrel{~}{}}{\nu }}^{\mathrm{}}`$ as follows;
$`G_{\underset{\stackrel{~}{}}{\nu }}^{\mathrm{}}(\tau ,z)={\displaystyle \underset{b_0_N}{}}g_{\underset{\stackrel{~}{}}{\nu }+b_0\underset{\stackrel{~}{}}{\gamma _0}}^\lambda (\tau ,z),`$
where $`\underset{\stackrel{~}{}}{\nu }`$ satisfies the beta condition (4.5). Then, from the modular properties of $`g_{\underset{\stackrel{~}{}}{\nu }}^\lambda `$
$`g_{\underset{\stackrel{~}{}}{\nu }}^\lambda (\tau +1,z)=𝐞\left[{\displaystyle \underset{j}{}}{\displaystyle \frac{\mathrm{}_j(\mathrm{}_j+1)}{4N_j}}{\displaystyle \frac{1}{2}}\underset{\stackrel{~}{}}{\nu }\underset{\stackrel{~}{}}{\nu }+{\displaystyle \frac{R+1}{8}}{\displaystyle \frac{1}{24}}\left({\displaystyle \underset{j}{}}{\displaystyle \frac{N_j2}{N_j}}+2\right)\right]g_{\underset{\stackrel{~}{}}{\nu }}^\lambda (\tau ,z),`$
$`g_{\underset{\stackrel{~}{}}{\nu }}^\lambda (1/\tau ,z/\tau )=(i)^R𝐞\left[{\displaystyle \frac{n}{2}}{\displaystyle \frac{z^2}{\tau }}\right]`$
$`\times {\displaystyle \underset{\lambda ^{},\underset{\stackrel{~}{}}{\nu }^{}}{\overset{\mathrm{even}}{}}}A_{\lambda \lambda ^{}}{\displaystyle \frac{1}{_j\sqrt{2N_j}}}{\displaystyle \frac{1}{\sqrt{2NJ}}}𝐞[\underset{\stackrel{~}{}}{\nu }\underset{\stackrel{~}{}}{\nu }^{}](1)^{\underset{\stackrel{~}{}}{\gamma _0}(\underset{\stackrel{~}{}}{\nu }\underset{\stackrel{~}{}}{\nu }^{})}g_{\underset{\stackrel{~}{}}{\nu }^{}}^\lambda ^{}(\tau ,z),`$
$`G_{\underset{\stackrel{~}{}}{\nu }}^{\mathrm{}}`$ have very good modular properties;
$`G_{\underset{\stackrel{~}{}}{\nu }}^\lambda (\tau +1,z)=𝐞\left[{\displaystyle \underset{j}{}}{\displaystyle \frac{\mathrm{}_j(\mathrm{}_j+1)}{4N_j}}{\displaystyle \frac{1}{2}}\underset{\stackrel{~}{}}{\nu }\underset{\stackrel{~}{}}{\nu }+{\displaystyle \frac{R+1}{8}}{\displaystyle \frac{1}{24}}\left({\displaystyle \underset{j}{}}{\displaystyle \frac{N_j2}{N_j}}+2\right)\right]G_{\underset{\stackrel{~}{}}{\nu }}^\lambda (\tau ,z),`$
$`G_{\underset{\stackrel{~}{}}{\nu }}^\lambda (1/\tau ,z/\tau )=(i)^R𝐞\left[{\displaystyle \frac{n}{2}}{\displaystyle \frac{z^2}{\tau }}\right]`$
$`\times {\displaystyle \underset{\lambda ^{},\underset{\stackrel{~}{}}{\nu }^{}}{\overset{\mathrm{even},\mathrm{beta}}{}}}A_{\lambda \lambda ^{}}{\displaystyle \frac{1}{_j\sqrt{2N_j}}}{\displaystyle \frac{1}{\sqrt{2NJ}}}𝐞[\underset{\stackrel{~}{}}{\nu }\underset{\stackrel{~}{}}{\nu }^{}](1)^{\underset{\stackrel{~}{}}{\gamma _0}(\underset{\stackrel{~}{}}{\nu }\underset{\stackrel{~}{}}{\nu }^{})}G_{\underset{\stackrel{~}{}}{\nu }^{}}^\lambda ^{}(\tau ,z).`$
Using these functions, we obtain the elliptic genus in the following form;
$`Z(\tau ,\overline{\tau },z)={\displaystyle \frac{1}{2^RN}}{\displaystyle \frac{1}{\sqrt{\tau _2}|\eta (\tau )|^2}}{\displaystyle \underset{\lambda ,\overline{\lambda },\underset{\stackrel{~}{}}{\nu }}{\overset{\mathrm{even},\mathrm{beta}}{}}}L_{\lambda \overline{\lambda }}G_{\underset{\stackrel{~}{}}{\nu }}^\lambda (\tau ,z)\overline{G}_{\underset{\stackrel{~}{}}{\nu }}^{\overline{\lambda }}(\overline{\tau },0).`$
Here $`L_{\lambda \overline{\lambda }}`$ is the product of $`\widehat{SU(2)}`$ modular invariants, and the factor $`1/\sqrt{\tau _2}|\eta (\tau )|^2`$ is contribution of $`\varphi `$. We can check that the above elliptic genus has the right modular properties (4.1) with $`\widehat{c}=n`$.
Actually, this elliptic genus is $`0`$ because it has an overall factor $`\overline{\theta }_1(\overline{\tau },0)=0`$.
### 4.1 Hodge number and Witten index
To get some nontrivial information from the above elliptic genus, we factor out the trivial parts and define $`\widehat{Z}`$ by the equations
$`Z(\tau ,\overline{\tau },z)={\displaystyle \frac{\theta _1(\tau ,z)\overline{\theta }_1(\tau ,0)}{\sqrt{\tau _2}|\eta (\tau )|^6}}\widehat{Z}(\tau ,\overline{\tau },z),`$
$`\widehat{Z}(\tau ,\overline{\tau },z)={\displaystyle \frac{1}{2^RN}}{\displaystyle \underset{\lambda ,\overline{\lambda },\underset{\stackrel{~}{}}{\nu }}{\overset{\mathrm{even},\mathrm{beta}}{}}}L_{\lambda \overline{\lambda }}\widehat{G}_{\underset{\stackrel{~}{}}{\nu }}^\lambda (\tau ,z)\overline{\widehat{G}}_{\underset{\stackrel{~}{}}{\nu }}^{\overline{\lambda }}(\overline{\tau },0),`$
$`\widehat{G}_{\underset{\stackrel{~}{}}{\nu }}^\lambda (\tau ,z)={\displaystyle \underset{b_0_N}{}}\widehat{g}_{\underset{\stackrel{~}{}}{\nu }+b_0\underset{\stackrel{~}{}}{\gamma _0}}^\lambda (\tau ,z),`$
$`\widehat{g}_{\underset{\stackrel{~}{}}{\nu }}^\lambda (\tau ,z)=I_\nu ^\lambda (\tau ,z)\mathrm{\Theta }_{M,NJ}(\tau ,2z/N)(1)^{\underset{\stackrel{~}{}}{\gamma _0}\underset{\stackrel{~}{}}{\nu }},`$
Now, we take the limit $`\tau i\mathrm{}`$ and consider the ground states. In this limit, $`\mathrm{\Theta }_{M,NJ}`$ becomes
$`\mathrm{\Theta }_{M,NJ}(i\mathrm{},z)=\delta _M^{\mathrm{mod}\mathrm{NJ}},`$
so, the $`\widehat{G}`$’s can be evaluated as
$`\widehat{G}_{\underset{\stackrel{~}{}}{\nu }}^\lambda =\{\begin{array}{cc}I_{\nu +\frac{M}{2J}\gamma _0}^\lambda (i\mathrm{},z)\hfill & (M0mod2J),\hfill \\ 0\hfill & (\text{others}).\hfill \end{array}`$
Then, $`\widehat{Z}`$ is expressed in the formula
$`\underset{\tau i\mathrm{}}{lim}\widehat{Z}={\displaystyle \frac{1}{2^RN}}{\displaystyle \underset{\lambda ,\overline{\lambda },\underset{\stackrel{~}{}}{\nu }}{\overset{\mathrm{even},\mathrm{beta}}{}}}\delta _M^{\mathrm{mod}\mathrm{\hspace{0.33em}2}\mathrm{J}}L_{\lambda \overline{\lambda }}I_{\nu +\frac{M}{2J}\gamma _0}^\lambda (i\mathrm{},z)\overline{I}_{\nu +\frac{M}{2J}\gamma _0}^{\overline{\lambda }}(i\mathrm{},0).`$
When we replace the $`\nu +\frac{M}{2J}\gamma _0`$ by $`\nu `$, then we can perform the sum of $`M_{2NJ}`$. Moreover, from the fact
$`I_{m_j}^\mathrm{}_j(i\mathrm{},z)=\delta _{m_j\mathrm{}_j1}^{\mathrm{mod}\mathrm{\hspace{0.33em}2}\mathrm{N}_\mathrm{j}}y^{\frac{\mathrm{}+1}{N}\frac{1}{2}}\delta _{m_j+\mathrm{}_j+1}^{\mathrm{mod}\mathrm{\hspace{0.33em}2}\mathrm{N}_\mathrm{j}}y^{\frac{\mathrm{}+1}{N}+\frac{1}{2}},`$
it can be seen that the even condition $`\mathrm{}_j+m_j1mod2`$ is included in this factor. We obtain a formula of the $`\widehat{Z}`$ in this limit
$`\underset{\tau i\mathrm{}}{lim}\widehat{Z}={\displaystyle \frac{1}{2^R}}{\displaystyle \underset{\nu }{\overset{\mathrm{beta}}{}}}{\displaystyle \underset{j}{}}\left[{\displaystyle \underset{\mathrm{}_j,\overline{\mathrm{}}_j}{}}L_{\mathrm{}_j\overline{\mathrm{}}_j}^{(N_j)}I_{m_j}^\mathrm{}_j(i\mathrm{},z)I_{m_j}^{\overline{\mathrm{}}_j}(i\mathrm{},0)\right].`$
So far, we treat rather general cases, but from now, we take an example and restrict ourselves to calculations in the example. We consider the example which satisfies all the following conditions.
* All minimal models are A type. So, $`L_{\lambda \overline{\lambda }}=\delta _{\lambda \overline{\lambda }}`$.
* $`R=n+1`$.
* $`N_1=N_2=\mathrm{}=N_R=N`$.
In other words, this example is the case where the associated Calabi-Yau manifold $`X`$ is the hypersurface of the form
$`z_1^N+z_2^N+\mathrm{}+z_{n+1}^N=0\text{ in }^{n+1}.`$ (4.6)
We can write the finite distance condition as $`N<n+1`$, which is equivalent to the condition that first Chern number of $`X/^\times `$ is positive. In this case, nontrivial factor of the elliptic genus can be calculated as
$`\underset{\tau i\mathrm{}}{lim}\widehat{Z}={\displaystyle \frac{1}{2^R}}{\displaystyle \underset{\nu }{\overset{\mathrm{beta}}{}}}{\displaystyle \underset{j}{}}\left[{\displaystyle \underset{\mathrm{}_j}{}}\left(\delta _{m_j\mathrm{}_j1}^{\mathrm{mod}\mathrm{\hspace{0.33em}2}\mathrm{N}_\mathrm{j}}y^{\frac{\mathrm{}+1}{N}\frac{1}{2}}\delta _{m_j+\mathrm{}_j+1}^{\mathrm{mod}\mathrm{\hspace{0.33em}2}\mathrm{N}_\mathrm{j}}y^{\frac{\mathrm{}+1}{N}+\frac{1}{2}}\right)\left(\delta _{m_j+\mathrm{}_j+1}^{\mathrm{mod}\mathrm{\hspace{0.33em}2}\mathrm{N}_\mathrm{j}}\delta _{m_j\mathrm{}_j1}^{\mathrm{mod}\mathrm{\hspace{0.33em}2}\mathrm{N}_\mathrm{j}}\right)\right]`$
$`={\displaystyle \frac{y^{\frac{n+1}{2}}}{2^R}}{\displaystyle \underset{\nu }{\overset{\mathrm{beta}}{}}}{\displaystyle \underset{j}{}}\left[{\displaystyle \underset{\mathrm{}_j}{}}\left(\delta _{m_j\mathrm{}_j1}^{\mathrm{mod}\mathrm{\hspace{0.33em}2}\mathrm{N}_\mathrm{j}}y^{\frac{\mathrm{}+1}{N}}+\delta _{m_j+\mathrm{}_j+1}^{\mathrm{mod}\mathrm{\hspace{0.33em}2}\mathrm{N}_\mathrm{j}}y^{\frac{\mathrm{}+1}{N}+1}\right)\right].`$
When we put $`m_j=a_j+Nb_j,(b_j=0,1;a_j=0,1,\mathrm{}N1)`$, then beta condition becomes
$`{\displaystyle \underset{j}{}}a_j0modN.`$
We obtain the $`\widehat{Z}`$ in this limit
$`\underset{\tau i\mathrm{}}{lim}\widehat{Z}={\displaystyle \underset{p=1}{\overset{n}{}}}h_py^{p\frac{n+1}{2}},`$
$`h_p:={\displaystyle \underset{\genfrac{}{}{0.0pt}{}{a_j=1,\mathrm{},N1,}{_ja_j=pN}}{}}1={\displaystyle \underset{i=0}{\overset{p}{}}}(1)^i\left(n+1{\displaystyle \genfrac{}{}{0.0pt}{}{i}{}}\right)\left((pi)(N1)+p1{\displaystyle \genfrac{}{}{0.0pt}{}{n}{}}\right).`$
We show several examples of $`h_p`$ for lower $`n,N`$ in Table 1.
These coefficients $`h_p`$ seem to coincide with the middle dimensional Hodge numbers of $`X/^\times `$ except for the cohomology elements generated by cup products of a Kähler form of $`X/^\times `$, on which we mention below.
In our model, the Witten index can also be calculated as
$`\underset{\tau i\mathrm{},z0}{lim}\widehat{Z}={\displaystyle \underset{p=1}{\overset{n}{}}}h_p`$
$`=(1)^{n+1}\left[1+{\displaystyle \frac{(1N)^{n+1}1}{N}}\right]`$
$`=(1)^{n+1}\left[n+1+{\displaystyle \frac{(1N)^{n+1}1}{N}}n\right].`$
On the other hand, the Euler number of the $`(n1)`$-dimensional manifold $`X/^\times `$ is expressed in the formula
$`\chi _{X/^\times }=n+1+{\displaystyle \frac{(1N)^{n+1}1}{N}}.`$
The Witten index of the CFT almost coincide with the Euler number $`\chi _{X/^\times }`$ of $`X/^\times `$. One of the differences of the two is the sign $`(1)^{n+1}`$, but this is not relevant. Except this difference of the sign, the Witten index is smaller by $`n`$ than the Euler number of $`X/^\times `$ in our case. This difference may correspond to the cohomology elements generated by cup products of the Kähler form of $`X/^\times `$. On $`X`$, these cohomology elements are absent because they appear when we take the quotient of $`X`$ by $`^\times `$. This seems the reason why the Witten index is smaller by $`n`$ than the Euler number of $`X/^\times `$.
## 5 Conclusion and discussion
We construct the toroidal partition function of the string theory described by the combination of an $`𝒩=2`$ Liouville theory and multiple $`𝒩=2`$ minimal models. This partition function is actually modular invariant, and we can conclude that the theory exists consistently.
This string theory is thought to describe the string on a noncompact singular Calabi-Yau manifold. To check this proposition, we also calculate the elliptic genus of this theory and the Witten index.
The Euler number defined from non-trivial factor of the Witten index in the CFT seems to be that of the non-vanishing elements of the cohomology. In the case of a singular manifold, there are vanishing elements of the cohomology, which are supported on the singular point and reflect the structure of singularities.
The fact that the vanishing elements of the cohomology cannot be seen, is probably related to our method of construction in which we include only the states in “principal continuous series” of the $`SL(2)`$ theory. If we can include some “discrete series” (but it is difficult) , the structure of the singularities might be seen in the CFT.
Another reason is that we treat the $`𝒩=2`$ Liouville theory as free field theory in this paper. It is mentioned in that if we treat appropriately the effect of Liouville potential, the Witten index does not vanish and gives the Euler number including the vanishing elements of the cohomology. In this paper, since we treat the case of $`\mu =0`$ and not deformed singularity, the vanishing elements of the cohomology actually vanish and it is consistent with the vanishing Witten index.
We may not be able to use our result to analyze the structure of the singularities, but we can use it to analyze the string on the positively curved manifold $`X/^\times `$. Especially it is interesting to analyze the D-branes wrapped on infinite cycle in this noncompact Calabi-Yau manifold through the recipes of boundary states in the CFT as the case of the ordinary Gepner models .
### Acknowledgement
I would like to thank Tsuneo Uematsu and Katsuyuki Sugiyama for useful discussions and encouragements. I would also like to thank to the organizers of Summer Institute 2000 at Yamanashi, Japan, 7-21 August, 2000, and the participants of it, especially, Michihiro Naka, Masatoshi Nozaki, Yuji Sato and Yuji Sugawara for useful discussions.
This work is supported in part by JSPS Research Fellowships for Young Scientists.
## Appendix A. Theta functions and characters
We use the following notations in this paper.
$`𝐞\left[x\right]:=\mathrm{exp}(2\pi ix),`$
$`\delta _m^{\mathrm{mod}\mathrm{N}}:=\{\begin{array}{cc}1\hfill & (m0modN),\hfill \\ 0\hfill & (\mathrm{others}),\hfill \end{array}`$
where $`m`$ and $`N`$ are integers. The useful formula is
$`{\displaystyle \underset{j_N}{}}𝐞\left[{\displaystyle \frac{jm}{N}}\right]=N\delta _m^{\mathrm{mod}\mathrm{N}},`$
where $`m`$ and $`N`$ are integers.
The SU(2) classical theta functions are defined as
$`\mathrm{\Theta }_{m,k}(\tau ,z)={\displaystyle \underset{n}{}}q^{k\left(n+\frac{m}{2k}\right)^2}y^{k\left(n+\frac{m}{2k}\right)},`$
where $`q:=𝐞\left[\tau \right],y:=𝐞\left[z\right]`$. The Jacobi’s theta functions are also defined as
$`\theta _1(\tau ,z):=i{\displaystyle \underset{n}{}}(1)^nq^{\left(n\frac{1}{2}\right)^2}y^{\left(n\frac{1}{2}\right)},\theta _2(\tau ,z):={\displaystyle \underset{n}{}}q^{\left(n\frac{1}{2}\right)^2}y^{\left(n\frac{1}{2}\right)},`$
$`\theta _3(\tau ,z):={\displaystyle \underset{n}{}}q^{n^2}y^n,\theta _4(\tau ,z):={\displaystyle \underset{n}{}}(1)^nq^{n^2}y^n.`$
Two kinds of theta functions are related by equations
$`2\mathrm{\Theta }_{0,2}=\theta _3+\theta _4,2\mathrm{\Theta }_{1,2}=\theta _2+i\theta _1,`$
$`2\mathrm{\Theta }_{2,2}=\theta _3\theta _4,2\mathrm{\Theta }_{3,2}=\theta _2i\theta _1.`$
The Dedekind $`\eta `$ function is represented as an infinite product
$`\eta (\tau ):=q^{\frac{1}{24}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1q^n).`$
The character $`\chi _s(\tau ,z),s=0,1,2,3`$ of $`\widehat{SO(d)}_1`$ for $`d/22+1`$ can be expressed as
$`\chi _0(\tau ,z)={\displaystyle \frac{\theta _3(\tau ,z)^{d/2}+\theta _3(\tau ,z)^{d/2}}{2\eta (\tau )^{d/2}}},`$
$`\chi _1(\tau ,z)={\displaystyle \frac{\theta _2(\tau ,z)^{d/2}+(i\theta _1(\tau ,z))^{d/2}}{2\eta (\tau )^{d/2}}},`$
$`\chi _2(\tau ,z)={\displaystyle \frac{\theta _3(\tau ,z)^{d/2}\theta _3(\tau ,z)^{d/2}}{2\eta (\tau )^{d/2}}},`$
$`\chi _3(\tau ,z)={\displaystyle \frac{\theta _2(\tau ,z)^{d/2}(i\theta _1(\tau ,z))^{d/2}}{2\eta (\tau )^{d/2}}}.`$
Let us denote the character of a Verma module $`(\mathrm{},m,s)`$ in the level $`(N2)`$ minimal model as $`\chi _m^{\mathrm{},s}(\tau ,z)`$. This character satisfies equivalence relations
$`\chi _m^{\mathrm{},s}=\chi _{m+2N}^{\mathrm{},s}=\chi _m^{\mathrm{},s+4}=\chi _{m+N}^{N2\mathrm{},s+2}.`$
The explicit form of this character is written in .
We collect the modular properties of these functions. Under the T transformations, they behave as
$`\mathrm{\Theta }_{m,k}(\tau +1,z)=𝐞\left[{\displaystyle \frac{m^2}{4k}}\right]\mathrm{\Theta }_{m,k}(\tau ,z),`$
$`\theta _1(\tau +1,z)=𝐞\left[\frac{1}{8}\right]\theta _1(\tau ,z),\theta _2(\tau +1,z)=𝐞\left[\frac{1}{8}\right]\theta _2(\tau ,z),`$
$`\theta _3(\tau +1,z)=\theta _4(\tau ,z),\theta _4(\tau +1,z)=\theta _3(\tau ,z),`$
$`\eta (\tau +1)=𝐞\left[1/24\right]\eta (\tau ),`$
$`\chi _s(\tau +1,z)=𝐞\left[{\displaystyle \frac{s^2}{8}}{\displaystyle \frac{d}{48}}\right]\chi _s(\tau ,z),`$
$`\chi _m^{\mathrm{},s}(\tau +1,z)=𝐞\left[{\displaystyle \frac{\mathrm{}(\mathrm{}+2)}{4N}}{\displaystyle \frac{m^2}{4N}}+{\displaystyle \frac{s^2}{8}}{\displaystyle \frac{N2}{8N}}\right]\chi _m^{\mathrm{},s}(\tau ,z),`$
and for S transformations, they have modular properties
$`\mathrm{\Theta }_{m,k}(1/\tau ,z/\tau )=\sqrt{i\tau }𝐞\left[{\displaystyle \frac{k}{4}}{\displaystyle \frac{z^2}{\tau }}\right]{\displaystyle \underset{m^{}_{2k}}{}}{\displaystyle \frac{1}{\sqrt{2k}}}𝐞\left[{\displaystyle \frac{mm^{}}{2k}}\right]\mathrm{\Theta }_{m^{},k}(\tau ,z),`$
$`\theta _1(1/\tau ,z/\tau )=i\sqrt{i\tau }𝐞\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{z^2}{\tau }}\right]\theta _1(\tau ,z),\theta _2(1/\tau ,z/\tau )=\sqrt{i\tau }𝐞\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{z^2}{\tau }}\right]\theta _4(\tau ,z),`$
$`\theta _3(1/\tau ,z/\tau )=\sqrt{i\tau }𝐞\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{z^2}{\tau }}\right]\theta _3(\tau ,z),\theta _4(1/\tau ,z/\tau )=\sqrt{i\tau }𝐞\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{z^2}{\tau }}\right]\theta _2(\tau ,z),`$
$`\eta (1/\tau )=\sqrt{i\tau }\eta (\tau ),`$
$`\chi _s(1/\tau ,z/\tau )=𝐞\left[{\displaystyle \frac{d}{4}}{\displaystyle \frac{z^2}{\tau }}\right]{\displaystyle \underset{s^{}=0}{\overset{3}{}}}{\displaystyle \frac{1}{2}}𝐞\left[{\displaystyle \frac{d}{2}}{\displaystyle \frac{ss^{}}{4}}\right]\chi _s^{}(\tau ,z),`$
$`\chi _m^{\mathrm{},s}(1/\tau ,z/\tau )=𝐞\left[{\displaystyle \frac{N2}{2N}}{\displaystyle \frac{z^2}{\tau }}\right]{\displaystyle \frac{1}{\sqrt{8N}}}{\displaystyle \underset{\mathrm{},m,s}{\overset{\mathrm{even}}{}}}A_{\mathrm{}\mathrm{}^{}}𝐞\left[{\displaystyle \frac{ss^{}}{4}}+{\displaystyle \frac{mm^{}}{2N}}\right]\chi _m^{}^{\mathrm{}^{},s^{}}(\tau ,z),`$
$`A_{\mathrm{}\mathrm{}^{}}=\sqrt{{\displaystyle \frac{2}{N}}}\mathrm{sin}\left[\pi {\displaystyle \frac{(\mathrm{}+1)(\mathrm{}^{}+1)}{N}}\right],`$
where the sum $`_{\mathrm{},m,s}^{\mathrm{even}}`$ means that $`\mathrm{}+m+s0mod2`$ for $`(\mathrm{},m,s)`$.
We use the notation $`f(\tau )`$ for a function $`f(\tau ,z)`$ of $`\tau ,z`$ with substituting $`z=0`$
$`f(\tau ):=f(\tau ,z=0).`$
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# Relativistic 𝐽-matrix method
## I Introduction
The J-matrix method, introduced by Heller and Yamani and developed by Yamani and Fishman , is an example of an algebraic method in quantum scattering theory. Comparing with the algebraic variational theories the method has been shown to be free of the false resonances problem . It has been used in construction of the Gauss quadrature of the continuum (see also ), the definition and analysis of a reproducing kernel in the context of Harris eigenvalues . Quite recently, it has been used in formulation of complex-scaling method and in development of the multi-channel Green’s functions by means of a complete $`L^2`$ basis.
The crux of the method is representation of the Hamiltonian in a suitable non-orthogonal basis changing the differential scattering problem into the purely algebraic one. Thus far only the non-relativistic version of the method has been formulated. The aim of this paper is to develop the simple relativistic formulation of the method in its theoretical framework for potentials sufficiently regular at the origin and vanishing at infinity faster than the Coulomb one.
## II Non-relativistic J-matrix - radial kinetic energy case
First we briefly review the non-relativistic Jacobi matrix approach introduced in Refs. and extended in . We recall only the case when the potential vanishing faster than the Coulomb one is involved, as we shall formulate the relativistic formalism for this kind of potential. The Coulomb case is much more complicated and it will be considered elswere. Let $`\{\varphi _n^l\}_{n=0}^{\mathrm{}}`$ be either Laguerre or Gaussian (Hermite) basis set. The explicit forms of both bases as well as some other formulas concerning non-relativistic problem (see Ref. ) are collected in table I. Only the second basis, i.e., the Gaussian one, forms an orthogonal set, hence, in general, the notion of biorthonormality is needed. The set $`\{\overline{\varphi }_n^l\}_{n=0}^{\mathrm{}}`$ is biorthonormal to $`\{\varphi _n^l\}_{n=0}^{\mathrm{}}`$ with respect to the unitary scalar product if $`\overline{\varphi }_m^l|\varphi _n^l_0^{\mathrm{}}\overline{\varphi }_m^l(\lambda r)\varphi _n^l(\lambda r)𝑑r=\delta _{mn}`$. Biorthonormal basis functions $`\{\overline{\varphi }_n^l\}`$ are also given in table I. The important feature of the sets $`\{\varphi _n^l\}`$ is that the radial kinetic energy operator:
$$H_0\frac{k^2}{2}\frac{1}{2}\frac{\mathrm{d}^2}{\mathrm{d}r^2}+\frac{l(l+1)}{2r^2}\frac{k^2}{2}$$
(1)
if expanded in any of them, takes the tridiagonal or Jacobi form:
$$J_{mn}\varphi _m^l|(H_0k^2/2)\varphi _n^l,J_{mn}0\text{only for}m=n,n\pm 1.$$
(2)
In the above $`k`$ is a wave number related to the energy $``$ and mass $`m`$ of the projectile
$$k^2=\frac{2m}{\mathrm{}^2}.$$
(3)
It must be stressed here that the matrix elements $`J_{mn}`$ are functions of $`k`$, i.e. $`J_{mn}=J_{mn}(k)`$. The regular solution $`S(k,r)`$ of the equation
$$(H_0k^2/2)S(k,r)=0$$
(4)
is simply proportional to the Riccati-Bessel function, satisfying $`S(k,r)r^{l+1}`$ as $`r0`$ and $`S(k,r)\stackrel{r\mathrm{}}{}\mathrm{sin}(kr\frac{\pi l}{2})`$. Using an expansion of $`S(k,r)`$ in the basis $`\{\varphi _n^l\}`$, i.e. $`S(k,r)=_{n=0}^{\mathrm{}}s_m^l\varphi _n^l(\lambda r)`$, one can write equation (4) in the form
$$\underset{n=0}{\overset{\mathrm{}}{}}J_{mn}s_n^l=0.$$
(5)
As shown in Ref. , using the explicit form of the matrix elements $`J_{mn}`$ one can find the expansion coefficients $`s_n^l`$ in terms of Gegenbauer polynomials (see table I). Again we have $`s_n^l=s_n^l(k)`$. In the J-matrix method to solve a scattering problem one introduces the second, cosine-like function $`C(k,r)`$, which is required to satisfy $`C(k,r)r^{l+1}`$ as $`r0`$ and $`C(k,r)\stackrel{r\mathrm{}}{}\mathrm{cos}(kr\frac{\pi l}{2})`$. It cannot be the second solution of the original, homogeneous problem as this solution, proportional to the Riccati-Neumann function, is singular at the origin.
The required $`C(k,r)`$ function has been found , in another way, namely by solving an inhomogeneous equation :
$$(H_0k^2/2)C(k,r)=\beta \overline{\varphi }_0^l(\lambda r),\beta =\frac{k}{2s_0^l}$$
(6)
with $`s_0^l`$ being the first expansion coefficient of sine solution. Then the expansion coefficients of $`C(k,r)`$ satisfy the equation
$$\underset{n=0}{\overset{\mathrm{}}{}}J_{mn}c_n^l=\beta \overline{\varphi }_0^l.$$
(7)
The corresponding coefficients $`c_n^l=c_n^l(k)`$ (see table I) have been also found by some differential technique . The calculated expansions $`S(k,r)=_{n=0}^{\mathrm{}}s_n^l\varphi _n^l(\lambda r)`$ and $`C(k,r)=_{n=0}^{\mathrm{}}c_n^l\varphi _n^l(\lambda r)`$ have been used in an approximate solution of the original scattering problem on the radial potential $`V=V(r)`$ vanishing faster then the Coulomb potential:
$$(H_0+V\frac{k^2}{2})\psi _E=0.$$
(8)
Namely, the potential $`V`$ has been replaced by a truncated potential operator
$$V^N=P_N^{}VP_N,$$
(9)
where $`P_N`$ is the generalised projection operation:
$$P_N=\underset{n=0}{\overset{N1}{}}|\varphi _n^l\overline{\varphi }_n^l|.$$
(10)
The new potential operator can be written in the basis $`\{\varphi _n^l\}`$ as an $`N\times N`$ matrix with the matrix elements $`V_{mn}^N=\varphi _n^l|V\varphi _n^l`$. Then the exact solution $`\psi _E^N`$ of the new problem:
$$(H_0+V^N\frac{k^2}{2})\psi _E^N=0$$
(11)
has been expanded in the basis $`\{\varphi _n^l\}`$ as
$$\psi _E^N(r)=\underset{n=0}{\overset{N1}{}}a_n^l\varphi _n^l+\underset{n=N}{\overset{\mathrm{}}{}}(s_n^l+\mathrm{tan}\delta _Nc_n^l)\varphi _n^l$$
(12)
to satisfy the boundary requirement $`\psi _E^N(r)\stackrel{r\mathrm{}}{}\mathrm{sin}(kr\frac{\pi l}{2})+\mathrm{tan}\delta _N\mathrm{cos}(kr\frac{\pi l}{2})`$. The $`\mathrm{tan}\delta _N`$ is an approximation of the tangent of the sought phase shift $`\delta `$ of the exact solution $`\psi _E`$ of the problem (8). The left-hand side projection of (11) onto the basis $`\{\varphi _n^l\}`$ gives then infinitely many equations depending on $`n`$. However all equations for $`nN+1`$ are satisfied automatically as coefficients $`s_n^l`$, $`c_n^l`$ satisfy the same recursion relation (5) for any $`m>0`$. The remaining finite set on equations involve $`N+1`$ unknowns $`\mathrm{tan}\delta _N`$, $`\{a_{nl}\}_{n=0}^{N1}`$. Those equations can be easily solved . In particular, using the recursion relation for matrix elements $`J_{nm}`$ the tangent can be calculated giving
$$\mathrm{tan}\delta _N=\frac{s_{N1}^l+g_{N1,N1}()J_{N,N1}s_N^l}{c_{N1}^l+g_{N1,N1}()J_{N,N1}c_N^l}$$
(13)
where $`g_{N1,N1}()=_{n=0}^{N1}\mathrm{\Gamma }_{N1,m}^2/(_m)`$ with the matrix $`\mathrm{\Gamma }`$ diagonalising the finite-dimensional problem $`(\mathrm{\Gamma }^{}P^{}(H_0+V\frac{k^2}{2})P\mathrm{\Gamma })_{mn}=(_n)\delta _{mn}`$. Here the energy dependent quantity $`g_{N1,N1}()`$ can be viewed as the matrix element of the inverse of the truncated operator $`P^{}(H_0+V^N\frac{k^2}{2})P`$ if restricted to the $`N`$-dimensional space where it does not vanish. The quantities $`_n`$ (see also and references therein).
## III Relativistic Jacobi-matrix problem
Now we shall turn to the relativistic problem. Before the formulation of the method we shall find the relativistic counterparts of $`S(k,r)`$ and $`C(k,r)`$ in some suitable basis. We shall also calculate the relativistic Jacobi matrix elements in this basis. For this purpose consider the free Dirac equation:
$`(_0E/c\mathrm{})\mathrm{\Psi }\left(\begin{array}{cc}(mc^2E)/c\mathrm{}& d/dr+\kappa /r\\ d/dr+\kappa /r& (mc^2E)/c\mathrm{}\end{array}\right)\left(\begin{array}{c}F(r)\\ G(r)\end{array}\right)=\left(\begin{array}{c}0\\ 0\end{array}\right)`$ (20)
In the above the total energy $`E`$ is related to the rest energy $``$ as $`E=+mc^2`$. Let $`l(\kappa )`$ be the non-negative solution of the equation $`l(l+1)=\kappa (\kappa +1)`$, i.e. $`l(\kappa )=\kappa `$ and $`l(\kappa )=\kappa 1`$ for positive and negative $`\kappa `$, respectively. We shall usually omit the symbol $`\kappa `$ in the notation throughout the text and write $`l`$ only, remembering that the latter depends on $`\kappa `$. Then equation (20) has two independent solutions. The first one, regular at the origin
$`\mathrm{\Psi }_{reg}(r)=\left(\begin{array}{c}F_{reg}(r)\\ G_{reg}(r)\end{array}\right)\left(\begin{array}{c}\widehat{ȷ}_l(\stackrel{~}{k}r)\\ \pm ϵ\widehat{ȷ}_{l1}(\stackrel{~}{k}r)\end{array}\right).`$ (25)
is constituted by the Riccati-Bessel functions with boundary behaviour
$$\widehat{ȷ}_l(x)\stackrel{x0}{}\frac{x^{l+1}}{(2l+1)!!}\text{and}\widehat{ȷ}_l(x)\stackrel{x\mathrm{}}{}\mathrm{sin}(x\frac{\pi l}{2})$$
(26)
The numbers $`ϵ`$ and $`\stackrel{~}{k}`$ in (25) are standard abbreviations
$`ϵ\sqrt{{\displaystyle \frac{Emc^2}{E+mc^2}}},\stackrel{~}{k}{\displaystyle \frac{\sqrt{(Emc^2)(E+mc^2)}}{c\mathrm{}}}.`$ (27)
The quantity $`\stackrel{~}{k}`$ converges in the non-relativistic limit $`c\mathrm{}`$ to the number $`k=\sqrt{\frac{2m}{\mathrm{}^2}}`$. The second solution of (20), irregular at zero is given by
$`\mathrm{\Psi }_{irr}(r)=\left(\begin{array}{c}F_{irr}(r)\\ G_{irr}(r)\end{array}\right)\left(\begin{array}{c}\widehat{n}_l(\stackrel{~}{k}r)\\ \pm ϵ\widehat{n}_{l1}(\stackrel{~}{k}r)\end{array}\right).`$ (32)
Here we have the Ricatti-Neumann functions with properties:
$$\widehat{n}_l(x)\stackrel{x0}{}\frac{(2l1)!!}{x^l}\text{and}\widehat{n}_l(x)\stackrel{x\mathrm{}}{}\mathrm{cos}(x\frac{\pi l}{2})$$
(33)
In both solutions (25) and (32) the upper and lower signs in the small components correspond to negative and positive $`\kappa `$, respectively. From the above it can be immediately seen that the regular solution $`\mathrm{\Psi }_{reg}`$ is the relativistic counterpart of non-relativistic function $`S(k,r)`$. For the sake of consistency with the non-relativistic case, hereafter we shall denote the relativistic sine-like solution $`\mathrm{\Psi }_{reg}`$ by $`\mathrm{\Psi }_S`$. To develop the Jacobi matrix analysis we have to introduce a suitable basis set.
### A The basis set
in the Hilbert space $`L^2(0,\mathrm{})C^2`$ on which the Dirac operator from (20) is defined. Let again $`\{\varphi _n^l(x)\}`$ be either the Laguerre or the Gaussian basis set and let $`\psi _n^l(\lambda r)=(\kappa /r+\mathrm{d}/\mathrm{d}r)\varphi _n^l(\lambda r)`$<sup>*</sup><sup>*</sup>* One should keep in mind that here dependence on $`\kappa `$ is not only present via $`l`$ coefficient, but via operator $`(\kappa /\mathrm{r}+\mathrm{d}/\mathrm{dr})`$. Then the basis set defined for our purposes is
$$\mathrm{\Phi }_{n\kappa }^+(r)\left(\begin{array}{c}\varphi _n^l(\lambda r)\\ 0\end{array}\right),\mathrm{\Phi }_{n\kappa }^{}(r)\left(\begin{array}{c}0\\ \psi _n^l(\lambda r)\end{array}\right)$$
(34)
The above set depends on the positive reals number $`\lambda `$ which can be treated as a nonlinear variational parameter (see, for instance, . Note that the set (34) satisfies the “kinetic balance condition”. The latter condition is generally defined as a requirement that, if the funcitions $`\{\gamma _i\}`$ are used to expand large component of solution of Dirac equation, then the basis $`\{\omega _i\}`$ used for expansion of small component should consist of linear combinations of functions $`\{(\kappa /r+\mathrm{d}/\mathrm{d}r)\gamma _i\}`$. Use of such a basis is the simplest way to omit the problem of so called “finite basis set disease” (see ) in estimation of bound states of the atomic system. It seems that it would be also interesting in future to consider the relativistic J-matrix problem in the context of the relativistic Sturmian basis (see and references therein) as it is known that the relativistic free particle Green function takes particularily simple form in this basis.
The biorthonormal elements to the functions (34) obviously are $`\overline{\mathrm{\Phi }}_{n\kappa }^+(r)=(\overline{\varphi }_n^l(\lambda r),0)^T`$, $`\overline{\mathrm{\Phi }}_{n\kappa }^{}(r)=(0,\overline{\psi }_n^l(\lambda r))^T`$. As usual, we denote by $`\overline{f_n}`$ the element biorthonormal to $`f_n`$. The elements $`\overline{\varphi }_n^{l(\kappa )}`$ are recalled in table I. Here we shall calculate the biorthonormal elements $`\overline{\psi }_n^{l(\kappa )}`$.
It is easy to show integrating by parts, that biorthonormal elements $`\{\overline{\psi }_n^l(x)\}_{n=0}^{\mathrm{}}`$ should satisfy the equation
$$_0^{\mathrm{}}\left[\left(\frac{\kappa }{r}\frac{\mathrm{d}}{\mathrm{d}r}\right)\overline{\psi }_m^l(\lambda r)\right]\varphi _n^l(\lambda r)d(r)=\delta _{mn}.$$
(35)
Hence it suffices only to solve the following inhomogeneous differential equation
$$\left(\frac{\kappa }{x}\frac{\mathrm{d}}{\mathrm{d}x}\right)\overline{\psi }_n^l(x)=\overline{\varphi }_n^l(x),x=\lambda r.$$
(36)
The resulting functions are given in Table II. They all belong to the space $`L^2(0,\mathrm{})`$. This fact is obvious apart from the case of negative $`\kappa `$ for the Gaussian set. This case needs more careful analysis as here it is not possible to give the functions by explicit formula. For all $`\kappa <0`$ the functions $`\{\overline{\psi }_n^l\}`$ due to Gaussian set behave as $`r^{l+2}`$ at the origin, and vanish at infinity not slower than $`r^{(l+1)}`$ as the limit of the occurring integral is finite. Then $`(\psi _n^l)^2`$ behaves for $`r\mathrm{}`$ as $`r^{2(l+1)}`$ and, as $`l`$ is nonnegative, $`\psi _n^l`$ exists. Thus in both cases, when $`\{\varphi _n^l\}`$ is either the Laguerre or the Gaussian basis set, all elements $`\{\overline{\psi }_n^l\}`$ biorthonormal to new functions $`\psi _n^l=(\kappa /r+d/dr)\varphi _n^l`$ belong to $`L^2(0,\mathrm{})`$. Then obviously biorthonormal elements $`\overline{\mathrm{\Phi }}_{n\kappa }`$ due to the relativistic case belong to the Hilbert space $`L^2(0,\mathrm{})C^2`$. Note that we do not need the explicit forms of biorthonormal functions in our considerations.
### B Expansions of relativistic sine and cosine solutions
Now we are in the position to find the expansions of sine-like $`\mathrm{\Psi }_S=\mathrm{\Psi }_{reg}`$ and cosine-like $`\mathrm{\Psi }_C`$ solutions. For the latter we demand to satisfy three requirements:
(1) $`\mathrm{\Psi }_C`$ should have $`\mathrm{\Psi }_{irr}`$ type asymptotic form,
(2) $`\mathrm{\Psi }_C`$ should exhibit regular behaviour at the origin,
(3) coefficients of $`\mathrm{\Psi }_C`$ expansion should satisfy (apart from at most the few first ones) the same recurrence equations like the ones of $`\mathrm{\Psi }_{reg}=\mathrm{\Psi }_S`$.
Consider first the solution $`\mathrm{\Psi }_U(\stackrel{~}{k},r)=\left(\begin{array}{cc}F_U(\stackrel{~}{k},r),& G_U(\stackrel{~}{k},r)\end{array}\right)^T`$ of the inhomogeneous equation of type (20):
$$(_0\frac{E}{c\mathrm{}})\mathrm{\Psi }_U\mathrm{\Phi }_{inh}.$$
(37)
In the above the index $`U=S,C`$ corresponds to sine-line and cosine-like solution. The inhomogeneity is chosen as $`\mathrm{\Phi }_{inh}=\mathrm{\Phi }_{0\kappa }^+=\left(\begin{array}{cc}\mathrm{\Omega }_U\overline{\varphi }_0^l,& 0\end{array}\right)^T`$ and the coefficients $`\mathrm{\Omega }_U,U=S,C`$ are $`\mathrm{\Omega }_S=0`$, $`\mathrm{\Omega }_C=ϵ/s_0^l`$.
Equation (37) can be also written as
$`\begin{array}{c}(\kappa /r\mathrm{d}/\mathrm{d}r)G_U\stackrel{~}{k}ϵF_U=\mathrm{\Omega }_U\overline{\psi }_0^l\\ (\kappa /r+\mathrm{d}/\mathrm{d}r)F_U\frac{\stackrel{~}{k}}{ϵ}G_U=0.\end{array}`$ (40)
We can introduce the relativistic counterpart of the Jacobi matrix :
$$𝒥_{mn}^{ss^{}}\mathrm{\Phi }_{m\kappa }^s|(_0E/c\mathrm{})\mathrm{\Phi }_{n\kappa }^s,s,s^{}=\pm ,m,n=0,1,2,\mathrm{}$$
(41)
The matrix elements of $`𝒥`$ can be expressed in an extremely simple form. To see this, let us define the $`2\times 2`$ matrices $`𝒥_{mn}`$ defined by their matrix elements as $`\{𝒥_{mn}\}_{ss^{}}𝒥_{mn}^{ss^{}}`$. Then it can be easily seen that in the spinor basis the new matrix takes the particularly simple form:
$$𝒥_{mn}=\left(\begin{array}{cc}\stackrel{~}{k}ϵ\varphi _m^l|\varphi _n^l& \psi _m^l|\psi _n^l\\ \psi _m^l|\psi _n^l& \frac{\stackrel{~}{k}}{ϵ}\psi _m^l|\psi _n^l\end{array}\right).$$
(42)
The explicit forms of the integrals constituting elements of the above matrix are given in table III. They are simply related to the non-relativistic J-matrix elements (2)(c.f. ):
$$J_{mn}=\frac{1}{2}\psi _m^l|\psi _n^l\frac{k^2}{2}\varphi _m^l|\varphi _n^l.$$
(43)
Now we shall predict the expansions of the two solutions in basis (34) in the following form
$$\mathrm{\Psi }_U=\underset{s=\pm }{}\underset{n=0}{\overset{\mathrm{}}{}}u_{n\kappa }^s\mathrm{\Phi }_{n\kappa }^s\underset{n=0}{\overset{\mathrm{}}{}}u_n^l(\stackrel{~}{k})\left(\begin{array}{c}\varphi _n^l\\ (ϵ/\stackrel{~}{k})\psi _n^l\end{array}\right),U=S,C;u=s,c$$
(44)
i.e. we predict that large components of sine-like and cosine-like solutions are given by the same expansion coefficients $`s_n^l,c_n^l`$ as in the non-relativistic case, only taken in the modified point $`\stackrel{~}{k}`$ and that the small components coefficients are only rescaled by $`ϵ/\stackrel{~}{k}`$.
It can be easily verified that (44) really solves the equation (20). Namely putting the above expansion into the equation and using the definition of matrix elements (41) we get the infinite set of equations:
$$\underset{s^{}=\pm }{}\underset{n=0}{\overset{\mathrm{}}{}}𝒥_{mn}^{ss^{}}u_{n\kappa }^s^{}=\mathrm{\Omega }_U\overline{\varphi }_0^l\delta _{m0}\delta _{s,+},s=\pm ,m=0,1,2,\mathrm{}$$
(45)
as for any pair $`m,n`$ fixed the second element in lower row of the matrix (42) is rescaled by $`\frac{\stackrel{~}{k}}{ϵ}`$ we obtain immediately that all equations (45) with a negative “index” $`s=`$“-” are satisfied trivially. Recalling the definition of $`l=l(\kappa )`$ (cf. the remark following equation (20)) after integration by parts one gets $`\psi _m^l|\psi _n^l=(\frac{\kappa }{r}+\frac{d}{dr})\varphi _m^l|(\frac{\kappa }{r}+\frac{d}{dr})\varphi _n^l=\varphi _m^l|(\frac{d^2}{dr^2}+\frac{l(l+1)}{r^2})\varphi _n^l`$. Taking into account the form of the upper row of the matrix (42) and the identity (43) we obtain immediately that equations (45) with the “index” $`s=`$“+” have the identical form with the sets of equations (5), (7) if the latter are evaluated at $`\stackrel{~}{k}`$ instead of $`k`$. Thus we have shown that the expansions (44) are in fact solutions of equation (40). From the non-relativistic case we see that their large components have the desired behaviour at the origin and at infinity. Moreover, as the equations of the type (20) are coupled and the behaviour of the one component determines the behaviour of the other one. Hence both the components of the solutions $`\mathrm{\Psi }_S`$, $`\mathrm{\Psi }_C`$ have the asymptotic behaviour we need for purposes of our method, i.e., $`\mathrm{\Psi }_S`$ is simply the regular solution, $`\mathrm{\Psi }_C`$ behaves as $`\mathrm{\Psi }_{irr}`$ at infinity and as $`\mathrm{\Psi }_{reg}`$ at the origin. As the inhomogeneity involves only one biorthonormal element $`\mathrm{\Psi }_{0\kappa }^+`$ both functions satisfy the same set of equations apart from the first one (see formula (42)).
## IV Potential scattering
Now we shall consider the central problem of our paper which is the approximate solution within the $`𝒥`$-matrix formalism. Consider again the radial part of the scattering problem of a projectile on a target described by a sufficiently regular potential $`V=V(r)`$ vanishing at infinity faster then the Coulomb potential. To solve the problem one has to find the solution of the following equation
$$\left(_0+\frac{V}{c\mathrm{}}\frac{E}{c\mathrm{}}\right)\mathrm{\Psi }_E=0.$$
(46)
A solution $`\mathrm{\Psi }_E`$ of the above equation is required to satisfy the boundary condition $`\mathrm{\Psi }_E(\stackrel{~}{k},r)\stackrel{r\mathrm{}}{}\mathrm{\Psi }_S(\stackrel{~}{k},r)+\stackrel{~}{t}\mathrm{\Psi }_C(\stackrel{~}{k},r)`$ where the tangent of the phase shift, $`\stackrel{~}{t}=\mathrm{tan}\stackrel{~}{\delta }`$, is to be found.
To develop the formalism of the relativistic J-matrix (we shall denote it by $`𝒥`$-matrix to distinguish from the non-relativistic case) we use the generalised projection operators:
$$𝒫_N=\underset{s=\pm }{}\underset{n=0}{\overset{N1}{}}|\mathrm{\Phi }_{n\kappa }^s\overline{\mathrm{\Phi }}_{n\kappa }^s|,$$
(47)
and introduce the truncated potential:
$$𝒱^N=𝒫_N^{}\frac{V}{c\mathrm{}}𝒫_N.$$
(48)
where $`𝒫_N^{}`$ corresponds to hermitian conjugate of $`𝒫_N`$.
Now one can seek the exact solution of the equation with truncated potential:
$$\left(_0+𝒱^N\frac{E}{c\mathrm{}}\right)\mathrm{\Psi }_E^N(r)=0.$$
(49)
Note that for any $`\psi L^2(0,\mathrm{})C^2`$ the function $`𝒱^N\psi `$ vanishes at infinity faster then $`\frac{1}{r}`$. Recall that we assumed that our original potentials vanish at infinity faster then $`\frac{1}{r^2}`$. Thus although equation (49) has not a standard Dirac equation form with the same scalar potential in its large and small part, still its solution asymptotically satisfies free Dirac equation. Hence the solution $`\mathrm{\Psi }_E^N`$ satisfies the boundary condition
$$\mathrm{\Psi }_E^N(\stackrel{~}{k},r)\mathrm{\Psi }_S(\stackrel{~}{k},r)+\stackrel{~}{t}_N\mathrm{\Psi }_C(\stackrel{~}{k},r),$$
(50)
where $`\stackrel{~}{t}_N`$ is an approximated tangent of phase shift. As the potential operator $`𝒱^N\stackrel{N\mathrm{}}{}\frac{V}{c\mathrm{}}`$ we expect that for $`N\mathrm{}`$, $`\stackrel{~}{t}_N`$ converges to correct value $`\stackrel{~}{t}=\mathrm{tan}\stackrel{~}{\delta }`$.
Now we shall find more details about the form of the solution $`\mathrm{\Psi }_E^N`$. The most general formula is
$$\mathrm{\Psi }_E^N=\underset{s=\pm }{}\underset{m=0}{\overset{\mathrm{}}{}}d_{m\kappa }^s|\mathrm{\Phi }_{m\kappa }^s.$$
(51)
Consider the matrix representation of equation (49). Putting the expansion of the function $`\mathrm{\Psi }_E^N`$ in the basis $`\{\mathrm{\Psi }_{n\kappa }^\pm \}`$ we get the infinite set of equations:
$$\underset{s^{}=\pm }{}\underset{n=0}{\overset{\mathrm{}}{}}(𝒥+𝒱^N)_{mn}^{ss^{}}d_{n\kappa }^s^{}=0,s=\pm ,m=0,1,2,\mathrm{}$$
(52)
It can be easily seen by the right-hand side projection of equations (40) onto the basis $`\{\mathrm{\Psi }_{n\kappa }^\pm \}`$.
According to analysis following the formula (45) the expansion coefficients $`\{d_{n\kappa }^\pm \}`$ of $`\mathrm{\Psi }_E^N`$ must satisfy: (a) for the large component $`_{n=0}^{\mathrm{}}J_{mn}(\stackrel{~}{k})d_{n\kappa }^+=0`$, $`m>N`$ with elements $`J_{mn}()`$ given by the non-relativistic formula, (b) for the small component $`d_{n\kappa }^{}=\frac{ϵ}{\stackrel{~}{k}}d_{n\kappa }^+`$. Moreover we impose the additional condition (c) $`F_E^NS(\stackrel{~}{k},r)+\stackrel{~}{t}_NC(\stackrel{~}{k},r)`$ (see condition (50)). This gives us, together with the condition (b), the following required form of the sought solution $`\mathrm{\Psi }_E^N`$ of equation (49) (c.f. ):
$$\mathrm{\Psi }_E^N=\underset{m=0}{\overset{N1}{}}\left(\begin{array}{c}d_{m\kappa }^+\varphi _m^l\\ d_{m\kappa }^{}\frac{ϵ}{\stackrel{~}{k}}\psi _m^l\end{array}\right)+\underset{m=N}{\overset{\mathrm{}}{}}\left(\begin{array}{c}(s_{\kappa m}^++\stackrel{~}{t}_Nc_{\kappa m}^+)\varphi _n^l\\ (s_{\kappa m}^{}+\stackrel{~}{t}_Nc_{\kappa m}^{})\psi _n^l\end{array}\right),$$
(53)
where the abbreviations $`s_{\kappa m}^\pm `$, $`c_{\kappa m}^\pm `$ has been used according to (44). After adding and subtracting the term $`_{m=0}^{N1}\left(\begin{array}{cc}(s_{\kappa m}^++\stackrel{~}{t}_Nc_{\kappa m}^+)\varphi _n^l,& (s_{\kappa m}^{}+\stackrel{~}{t}_Nc_{\kappa m}^{})\psi _n^l\end{array}\right)^T`$ to the left hand side of the above equation it is straightforward to see that the above function satisfies the asymptotic condition (50).
Let us turn back to equations (52). In general, in analogy to the non-relativistic case, they can be schematically represented as follows:
$$\left(\begin{array}{cccccccccccccccccc}& & & & & & & & & & & & & & & & & \\ & X& X& X& & & & & 0& & & & & & X& X& X& \\ & & X& X& X& & & & & & & & & X& X& X& & \\ & & & X& X& X& X& X& X& & & & X& X& X& & & \\ & & & & X& X& X& X& X& & & X& X& X& & & & \\ & & & & X& X& X& X& X& & X& X& X& & & & & \\ & & & & X& X& X& X& X& X& X& X& & & & & & \\ & & & & X& X& X& X& X& X& X& & & & & & 0& \\ & 0& & & & & & X& X& X& X& X& X& X& & & & \\ & & & & & & X& X& X& X& X& X& X& X& & & & \\ & & & & & X& X& X& & X& X& X& X& X& & & & \\ & & & & X& X& X& & & X& X& X& X& X& & & & \\ & & & X& X& X& & & & X& X& X& X& X& X& & & \\ & & X& X& X& & & & & & & & & X& X& X& & \\ & X& X& X& & & & & & 0& & & & & X& X& X& \\ & & & & & & & & & & & & & & & & & \end{array}\right)\left(\begin{array}{c}\\ s_{N+1,l}^+\\ s_{N,l}^+\\ d_{N1,l}^+\\ \\ \\ d_{1,l}^+\\ d_{0,l}^+\\ d_{0,l}^{}\\ d_{1,l}^{}\\ \\ \\ d_{N1,l}^{}\\ s_{N,l}^{}\\ s_{N+1,l}^{}\\ \end{array}\right)=\left(\begin{array}{c}\\ 0\\ 0\\ 0\\ \\ \\ 0\\ 0\\ 0\\ 0\\ \\ \\ 0\\ 0\\ 0\\ \end{array}\right)$$
From the construction of the required form (53) we see that all the equations for $`m>N`$ are satisfied automatically. Thus one has to solve the remaining $`2N+2`$ equations with the unknowns $`\stackrel{~}{t}_N`$, $`d_{0\kappa }^+,d_{1\kappa }^+,\mathrm{},d_{N1,\kappa }^+;d_{0\kappa }^{},d_{1\kappa }^{},\mathrm{},d_{N1,\kappa }^{}`$. Note that here the number of equations is greater then the number of sought quantities ($`2N+1`$), so in general the set of equations of such a form can have no solution. But in our particular case the solution certainly exists as the general theory of differential equations assures the existence of $`\mathrm{\Psi }_E^N`$ and, according to the previous analysis, (53) represents the most general required form of $`\mathrm{\Psi }_E^N`$.
Using equations (52) one obtains the following form of the remaining equations:
$$\left(\begin{array}{ccccccc}𝒥_{N,N1}^{++}c_{N1}^+& 𝒥_{N,N1}^{++}& 0& \mathrm{}& 0& 𝒥_{N,N1}^+& 𝒥_{N,N1}^+c_{N1}^{}\\ 𝒥_{N1,N}^{++}c_N^+& (𝒥+𝒱^N)_{N1,N1}^{++}& (𝒥+𝒱^N)_{N1,N2}^{++}& \mathrm{}& 𝒥_{N1,N2}^+& 𝒥_{N1,N1}^+& 𝒥_{N1,N}^+c_N^{}\\ 0& (𝒥+𝒱^N)_{N2,N1}^{++}& (𝒥+𝒱^N)_{N2,N2}^{++}& \mathrm{}& 𝒥_{N2,N2}^+& 𝒥_{N2,N1}^+& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& (𝒥+𝒱^N)_{0,N1}^{++}& (𝒥+𝒱^N)_{0,N2}^{++}& \mathrm{}& (𝒥+𝒱^N)_{0,N2}^+& (𝒥+𝒱^N)_{0,N1}^+& 0\\ 0& (𝒥+𝒱^N)_{0,N1}^+& (𝒥+𝒱^N)_{0,N2}^+& \mathrm{}& (𝒥+𝒱^N)_{0,N2}^{}& (𝒥+𝒱^N)_{0,N1}^+& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 𝒥_{N2,N1}^+& 𝒥_{N2,N2}^+& \mathrm{}& (𝒥+𝒱^N)_{N2,N2}^{}& (𝒥+𝒱^N)_{N2,N1}^{}& 0\\ 𝒥_{N1,N}^+c_N^+& 𝒥_{N1,N1}^+& 𝒥_{N1,N2}^+& \mathrm{}& (𝒥+𝒱^N)_{N1,N2}^{}& (𝒥+𝒱^N)_{N1,N1}^{}& 𝒥_{N1,N}^{}c_N^{}\\ 𝒥_{N,N1}^+c_{N1}^+& 𝒥_{N,N1}^+& 0& \mathrm{}& 0& 𝒥_{N,N1}^{}& 𝒥_{N,N1}^{}c_{N1}^{}\end{array}\right)$$
$$\times \left(\begin{array}{c}\stackrel{~}{t}_N\\ d_{N1}^+\\ d_{N2}^+\\ \\ d_0^+\\ d_0^{}\\ \\ d_{N2}^{}\\ d_{N1}^{}\\ \stackrel{~}{t}_N\end{array}\right)=\left(\begin{array}{c}𝒥_{N,N1}^{++}s_{N1}^++𝒥_{N,N1}^+s_{N1}^{}\\ 𝒥_{N1,N}^{++}s_N^+𝒥_{N1,N}^+s_N^{}\\ 0\\ \\ 0\\ 0\\ \\ 0\\ 𝒥_{N1,N}^+s_N^+𝒥_{N1,N}^{}s_N^{}\\ 𝒥_{N,N1}^+s_{N1}^++𝒥_{N,N1}^{}s_{N1}^{}\end{array}\right)$$
Keeping in mind that the inner $`2N\times 2N`$ matrix $`(𝒥+𝒱^N)_{nm}^{ss^{}}`$, $`s=\pm `$, $`m,n=0,1,\mathrm{},N1`$ is Hermitian and real, hence symmetric, and recalling the definition of $`𝒱^N`$ we can solve the above equations by some orthogonal matrix $`\mathrm{\Gamma }`$ (cf. the non-relativistic case):
$$(\mathrm{\Gamma }^{}𝒫_n^{}(_0+\frac{V}{c\mathrm{}}\frac{E}{c\mathrm{}})𝒫_N\mathrm{\Gamma })_{mn}^{ss^{}}=\frac{1}{c\mathrm{}}(E_n^sE)\delta _{nm}\delta _{ss^{}}$$
(54)
$`2N\times 2N`$ matrix $`𝒢(E)`$ with elements defined as
$$𝒢_{mn}^{ss^{}}(E)=\underset{p=\pm }{}\underset{i=0}{\overset{N1}{}}c\mathrm{}\frac{\mathrm{\Gamma }_{m,i}^{sp}\mathrm{\Gamma }_{n,i}^{s^{}p}}{E_i^pE}$$
(55)
is an inverse of the $`2N\times 2N`$ matrix representation of the truncated operator $`𝒫_N^{}(_0+\frac{V}{c\mathrm{}}\frac{E}{c\mathrm{}})𝒫_N`$. It can be viewed as the approximation of the relativistic Green function in the basis (34). The numbers $`E_n^p`$ are the relativistic counterparts of the Harris eigenvalues . They represent a finite approximation of the spectrum of the relativistic Hamiltonian $`_0+\frac{V}{c\mathrm{}}`$. In particular, they include positive approximations of the first $`N`$ energy levels due to the potential $`V`$ and the $`N`$ negative pseudo-energies due to the continuous spectrum. As our basis (34) satisfies the kinetic balance condition there is a hope that $`E_i^p`$ satisfy the generalised form of the Hylleraas-Undheim theorem (see, for instance and references therein). It means, in particular, that (i) N positive values among set $`\{E_i^p\}`$ approximate the exact eigenenergies $`_0+\frac{V}{c\mathrm{}}`$ from the above and that (ii) the remaining N eigenvalues have values below $`mc^2`$.
We can introduce now the $`(2N+2)\times (2N+2)`$ block-diagonal matrix :
$$\stackrel{~}{\mathrm{\Gamma }}_{(2N+2)\times (2N+2)}=diag(1,\mathrm{\Gamma }_{2N\times 2N},1)$$
(56)
and act with it on the left-hand side of the above set of $`2N+2`$ equations. Using the fact that the matrix $`𝒥_{N,N1}`$ given by (42) is nonsingular the set of $`(2N+2)\times (2N+2)`$ the equations can be solved with respect to the approximate tangent of phase shift. Using the properties of the coefficients of the matrix $`𝒥`$ one can derive the tangent of the appoximated phase shift in the form similar to the non-relativistic formula:
$$\stackrel{~}{t}_N=\frac{s_{N1}^l(\stackrel{~}{k})+(2ϵ/\stackrel{~}{k})𝒢_{N1,N1}^{++}(E)J_{N,N1}(\stackrel{~}{k})s_N^l(\stackrel{~}{k})}{c_{N1}^l(\stackrel{~}{k})+(2ϵ/\stackrel{~}{k})𝒢_{N1,N1}^{++}(E)J_{N,N1}(\stackrel{~}{k})c_N^l(\stackrel{~}{k})}.$$
(57)
Note that in the above the $`J_{N,N1}`$ stands for the non-relativistic $`J`$-matrix element (see (43)). The fact that we have $`2N+2`$ equations and $`2N+1`$ unknows results in second, very similar formula for $`\stackrel{~}{t}_N`$ with $`(\stackrel{~}{k}/ϵ)𝒢_{N1,N1}^+(E)`$ instead of $`𝒢_{N1,N1}^{++}(E)`$. From the previous analysis we know that both equations must give the same $`\stackrel{~}{t}_N`$ which means that one has $`𝒢_{N1,N1}^+(E)=(ϵ/\stackrel{~}{k})𝒢_{N1,N1}^{++}(E)`$.
## V Discussion
Comparing equation (57) with the non-relativistic formula (13) one can see that apart from the quantity $`(2ϵ/\stackrel{~}{k})𝒢_{N1,N1}^{++}(E)`$, all elements of the expression for tangent of the phase shift have the same form as in (13), they are only evaluated in relativistic wave number $`\stackrel{~}{k}`$.
Now let us note that for any $`N`$ the above formula for tangent shift converges to the non-relativistic limit as the speed of light $`c`$ approaches infinity. Indeed, the used basis (34) ensures (see ) that in the limit of infinite $`c`$ the large component satisfies the correct Schrödinger equation (11) with the wave number $`k=lim_c\mathrm{}\stackrel{~}{k}`$. This means that the related tangent of the phase shift must also satisfy a correct limit, i.e.
$$\underset{c\mathrm{}}{lim}\stackrel{~}{t}_N=t_N.$$
(58)
From the above we get immediately $`lim_c\mathrm{}(2ϵ/\stackrel{~}{k})𝒢_{N1,N1}^{++}(E)=g_{N1,N1}()`$. Moreover $`(2ϵ/\stackrel{~}{k})𝒢_{N1,N1}^{++}(E)`$ plays the analogous role as $`g_{N1,N1}()`$. In fact, the matrices $`𝒢(E)`$ and $`g()`$ can be viewed as the finite approximations of the Green functions of the relativistic and non-relativistic Hamiltonians with the potential $`V`$, respectively. The form of the factor $`2ϵ/\stackrel{~}{k}`$ is simply connected with the normalisations of the Green functions in both cases. It can be seen from the simple analysis of the set of second order equations derived in a standard way from the Dirac equation.
From the practical point of view, the convergence can be improved with the help of additional parameter $`\lambda `$. As we mentioned before, the latter can be treated as an additional variational parameter. In particular its optimal value will depend on the range of the potential. It can be simply seen that potentials of long range should be treated with small $`\lambda `$ while potentials with support located close to the origin will require large values of the parameter.
In conclusion, we have provided the relativistic version of Jacobi matrix method for well defined class of potentials. The usage of the basis satisfying the “kinetic balance condition” allowed for a simple formulation of the method. In particular, the derived expression for the tangent of the phase shift is similar to its non-relativistic counterpart and reproduces the latter as a correct non-relativistic limit.
The author is especially grateful to R. Szmytkowski for suggesting the problem, many helpful discussions, comments and remarks. He also thanks J. E. Sienkiewicz for discussion on kinetic balance condition and P. Syty for remarks on the manuscript. The work is supported by the Committee for Scientific Research (Poland) under project No. 2P03B 000912. The support from Foundation for Polish Science is also gratefully acknowledged.
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# PERSPECTIVES IN HIGH-ENERGY PHYSICS
## 1 The Standard Model
The Standard Model continued to rule accelerator experiments during 1999, even as the heroic efforts of the CERN accelerator engineeers pushed the LEP centre-of-mass energy to 202 GeV, and briefly to 204 GeV. There were no surprises in fermion-pair production or in the bread-and-butter reaction of LEP2, $`e^+e^{}W^+W^{}`$. Both the $`\gamma W^+W^{}`$ and $`Z^0W^+W^{}`$ triple-gauge-boson vertices are there, as seen in Fig. 1, with magnitudes close to the Standard Model values . Looking into the final states, there are no confirmed interferences in $`(W^+\overline{q}q)(W^{}\overline{q}q)`$ final states, due to either colour rearrangement or Bose-Einstein correlations: the difference between the $`W`$ mass measured in purely hadronic and other final states is 15 $`\pm `$ 55 MeV . Combining all the LEP measurements, one finds
$$m_W=80.401\pm 0.048\mathrm{GeV}(\mathrm{LEP}),$$
(1)
contributing with the hadron colliders ($`M_W=80.448\pm 0.062`$ GeV) to a global average
$$m_W=80.419\pm 0.038\mathrm{GeV}(\mathrm{world}).$$
(2)
This error is now comparable with the value estimated indirectly from precision electroweak measurements: $`m_W=80.382\pm 0.026`$ GeV, provides new, independent evidence for a light Higgs boson, and begins to impact significantly the radiative-correction estimate
$$m_H=77_{39}^{+69}\mathrm{GeV}$$
(3)
when $`\alpha _{em}(m_Z)^1=128.878\pm 0.090`$ is assumed (or log$`(M_H/\mathrm{GeV})=1.96_{0.23}^{+0.21}`$ if the estimate $`128.905\pm 0.036`$, with the error reduced by theory, is assumed). The Higgs boson probably weighs less than 200 GeV.
The plan is to raise the LEP energy as high as possible during 2000, with the primary aim of searching for the Higgs boson. An integrated luminosity of 50/pb per experiment at 206 GeV would increase the sensitivity of the Higgs reach from the current lower limit of 107.9 GeV to about 114 GeV . A small numerical increase, but in the most interesting range, also from the point of view of supersymmetry . The LEP energies attained so far range up to 208.7 GeV, with a total luminosity (so far) of 109 pb<sup>-1</sup> at an average energy above 205 GeV. It seems that the target sensitivity to $`m_H=114`$ GeV is well within reach. At the time of writing, the current sensitivity is to $`m_H113.4`$ GeV, and the latest update may be obtained from .
Then, in Autumn 1999, LEP must be shut down and dismantled to make way for the LHC excavations and installation. It will be the end of an era of precision electroweak measurements. The search for the Higgs boson will then pass to Fermilab, where the Tevatron has a chance of exploring higher Higgs masses if it gathers more than 10 fb<sup>-1</sup> of luminosity, as seen in Fig. 2 .
## 2 CP Violation
For 24 years, it was possible to think that the CP violation seen in the $`K^0\overline{K}^0`$ system might be due entirely to a superweak force beyond the Standard Model, inducing CP violation in the $`K^0\overline{K}^0`$ mass matrix: Im$`M_{12}0`$. Thanks to the new measurements of $`ϵ^{}/ϵ`$ in 1999 , signalling direct CP violation in the $`K^02\pi `$ decay amplitudes, we now know the superweak theory cannot be the whole story. The Standard Model with 6 quarks permits CP violation in the $`W^\pm `$ couplings , but there are no other sources if there is just one Higgs doublet. The MSSM contains many possible sources of CP violation, of which the most important for hadron decays may be the phases of the trilinear soft supersymmetry-breaking couplings $`A_{t,b}`$ of the third generation and of the gluino mass, relative to the Higgs mixing parameter $`\mu `$.
In the Standard Model, one may estimate
$`\mathrm{Re}\left({\displaystyle \frac{ϵ^{}}{ϵ}}\right)13\mathrm{Im}\lambda _t\left({\displaystyle \frac{130\mathrm{MeV}}{m_s(m_c)}}\right)^2\left({\displaystyle \frac{\mathrm{\Lambda }_{\overline{MS}}^{(4)}}{340\mathrm{MeV}}}\right)`$
$`\times \left[B_6(1\mathrm{\Omega })0.4B_8\left({\displaystyle \frac{m_t(m_t)}{165\mathrm{GeV}}}\right)^{2.5}\right]`$ (4)
where Im$`\lambda _t`$ is a Cabibbo-Kobayashi-Maskawa angle factor and $`\mathrm{\Omega }`$ is an electroweak penguin effect. Calculating Re$`(ϵ^{}/ϵ)`$ accurately is difficult because of potential cancellations in (4), and the fact that the strong-interaction matrix elements $`B_6`$ and $`B_8`$ are relatively poorly known. A conservative estimate in the Standard Model would be that
$$1\times 10^4<\mathrm{Re}(ϵ^{}/ϵ)<3\times 10^3$$
(5)
The world data, including the new NA48 measurement reported here , average to
$$\mathrm{Re}(ϵ^{}/ϵ)=(19.3\pm 2.4)\times 10^4$$
(6)
as seen in Fig. 3.
This is compatible with the Standard Model range (5), although probably requiring a relatively big penguin matrix element $`B_6`$: an emperor maybe? The measurement of $`(ϵ^{}/ϵ)`$ will be refined by further KTeV and NA48 data, and by KLEO at DAFNE.
If the relatively large value of $`(ϵ^{}/ϵ)`$ has a large supersymmetric contribution, this could show up in $`K_L^0\pi ^0\overline{\nu }\nu `$, $`\pi ^0e^+e^{}`$ and $`\pi ^0\mu ^+\mu ^{}`$ decays . The present upper limits on these decays from KTeV are far above the Standard Model predictions, so here is a physics opportunity worth pursuing in parallel with $`B`$ experiments.
The main target of these experiments will be the Standard Model unitarity triangle shown in Fig. 4 , with the hope of finding a discrepancy: perhaps it will turn out to be a quadrilateral? The poster child for the dawning CP age is the measurement of $`\mathrm{sin}2\beta `$ via $`B^0J/\psi K_S`$ decay. The theory is gold-plated – with penguin pollution expected only at the $`10^3`$ level – and the experiment is clean. Indeed, between them, OPAL , CDF and ALEPH have almost measured it:
$$\mathrm{sin}2\beta =0.91\pm 0.35$$
(7)
A new era of precision flavour physics is now dawning, with the $`B`$ factories now starting to take data. They should reduce the error in $`\mathrm{sin}2\beta `$ below 0.1, and the LHC experiments aim at an error $``$ 0.01.
That is the good news: the bad news concerns the measurement of $`\mathrm{sin}2\alpha `$ via $`B^0\pi ^+\pi ^{}`$ decay. This mode has now been seen by CLEO , but with a relatively low branching ratio
$$B(\pi ^+\pi ^{})=(0.43_{0.14}^{+0.16}\pm 0.05)\times 10^5$$
(8)
whereas
$$B(K^+\pi ^{})=(1.72_{0.24}^{+0.25}\pm 0.12)\times 10^5$$
(9)
This is bad news because it suggests that penguin pollution is severe, and difficult to disentangle, as well as threatening a large background. This worry has revived interest in the study of $`\mathrm{sin}2\alpha `$ via $`B^0\rho ^\pm \pi ^{}`$ and other decays. There is no time here to review proposed strategies for measured $`\gamma `$: suffice to say that within a decade we may know the sum $`\alpha +\beta +\gamma `$ with a precision $`10^0`$. Let us hope that it turns out to be inconsistent with 180<sup>0</sup>!
## 3 Beyond the Standard Model
The Standard Model has many defects, the most severe being that it agrees with all accelerator data. Nevertheless, it is very unsatisfactory, in that it provides no explanations for the particle quantum numbers $`(Q_{em},I,Y`$, colour), and contains 19 arbitrary parameters: 3 independent gauge couplings, 1 CP-violating strong-interaction phase, 6 quark masses, 3 charged-lepton masses, 4 Cabibbo-Kobayashi-Maskawa mixing parameters, the $`W`$ mass and the Higgs mass.
As if this was not enough, the indications of neutrino oscillations introduce at lest 9 more parameters: 3 neutrino masses, 3 neutrino mixing angles and 3 CP-violating phases, without even talking about the mechanism for neutrino mass generation. Also, we should not forget gravity, with Newton’s constant and the cosmological constant as two new parameters, at least one more parameter to generate the baryon asymmetry of the Universe, at least one more to describe cosmological inflation, etc.
It is common to group the open problems beyond the Standard Model into three categories: Mass – do the particle masses indeed originate from a Higgs boson, and is this accompanied by supersymmetric particles? Flavour – why are there so many flavours of quarks and leptons, and what explains the ratios of their masses and weak mixings? and Unification – is there a simple group structure containing the strong, weak and electromagnetic interactions? Beyond these beyonds lurks the quest for a Theory of Everything that should also include gravity, reconcile it with quantum mechanics, explain the origin of space-time and why it has 4 dimensions, and make Colombian coffee. This is the ambition of string theory and its latest incarnation as $`M`$ theory.
Subsequent sections of this talk deal with these ideas and how they may be tested.
## 4 Neutrino Masses and Oscillations
Why not? Although the Standard Model predicts $`m_\nu =0`$ if one ignores possible non-renormaliz- able interactions, there is no deep reason why this should be so. Theoretically, we expect masses to vanish only if there is a good asymmetry reason, in the form of an exact gauge symmetry. For example, the $`U(1)`$ gauge symmetry of QED guarantees the conservation of electric charge and the masslessness of the photon. However, there is no exact gauge symmetry or massless gauge boson coupled to lepton number $`L`$, which we therefore expect to be violated. Neutrino masses are possible if there is an effective $`\mathrm{\Delta }L=2`$ interaction of the Majorana form $`m_\nu \nu \nu `$. Such interactions are generic in Grand Unified Theories (GUTs) with their extra particles, but could even be fabricated from Standard Model particles alone , if one allows non-renormalizable interactions of the form $`\frac{1}{M}\nu H\nu Hm_\nu =<0|H|0>^2/M`$, where $`M`$ is some heavy mass scale: $`Mm_W`$.
Nevertheless, generic neutrino mass terms, as they arise in typical GUTs, have the seesaw form
$$(\nu _L,\nu _R)\left(\begin{array}{cc}0& m\\ m^T& M\end{array}\right)\left(\begin{array}{c}\nu _L\\ \nu _R\end{array}\right)$$
(10)
where the $`\nu _R`$ are singlet “right-handed” neutrinos. After diagonalization, (10) yields
$$m_\nu =m\frac{1}{M}m^T$$
(11)
which is naturally small: $`m_\nu m_{q,\mathrm{}}`$, if $`mm_{q,\mathrm{}}`$ and $`M=𝒪(M_{GUT})`$. For example, if $`m`$ 10 GeV and $`M10^{13}`$ GeV, one finds $`m_\nu 10^2`$ eV, in the range indicated by experiments on solar and atmospheric neutrinos. Each of the fields $`(\nu _L,\nu _R)`$ in (10) should be regarded as a three-dimensional vector in generation space, so $`m_\nu `$ (11) is a 3$`\times `$3 matrix. its diagonalization $`V_\nu `$ relative to that of the charged leptons $`V_{\mathrm{}}`$ yields the Maki-Nakagawa-Sakata neutrino mixing matrix $`V_{MNS}=V_{\mathrm{}}V_\nu ^+`$ between the interaction eigenstates $`\nu _{e,\mu ,\tau }`$ .
Could there be additional light neutrinos? The LEP neutrino counting measurements tell us that these could only be sterile neutrinos $`\nu _s`$. But what would prevent them from acquiring large masses $`m_s\nu _s\nu _s:m_sm_W`$, since they have no electroweak quantum numbers to forbid them via selection rules? This is exactly what happens to the $`\nu _R`$ in (10). Most theorists expect the observed neutrinos $`\nu _{light}\nu _L`$, and their effective mass term to be of the Majorana type. Before the advent of the atmospheric neutrino data, most theorists might have favoured small neutrino mixing angles, by analogy with the CKM mixing of quarks. But this is not necessarily the case, and many plausible models have now been constructed in which the neutrino mixing is large, because of large mixing in the light-lepton sector $`V_{\mathrm{}}`$ and/or in the heavy Majorana mass matrix $`M`$ and hence $`V_\nu `$.
In hierarchical models of neutrino models, as often arise in GUTs: $`m_1^2m_2^2m_3^2`$, one may interpret $`\mathrm{\Delta }m^2`$ as $`m_{\nu _H}^2`$, the squared mass of the heavier eigenstate in the oscillation. Cosmological data exclude $`m_\nu >`$ 3 eV, but even $`m_\nu `$ 0.03 eV may be of cosmological importance . Thus, although the unconfirmed LSND signal would certainly be of cosmological interest, so also are the confirmed atmospheric neutrino data.
Even the solar neutrino data become of interest to cosmology if the three neutrino flavous are almost degenerate, with $`\mathrm{\Delta }m^2m`$, a possibility that cannot be excluded by the oscillation data. Such a degeneracy is constrained by the upper limit on neutrinoless $`\beta \beta `$ decay :
$$<m_\nu >_e<0.2\mathrm{eV}$$
(12)
where the expectation value is weighted by the neutrinos’ electron couplings. The limit (12) is compatible with the neutrino masses close to the cosmological upper limit (which nearly coincides with the direct upper limit on $`m_{\nu _e}`$ from the end-point of Tritium $`\beta `$ decay) if there is almost maximal neutrino mixing. There is an issue whether this can be maintained in the presence of the mass renormalization expected in GUTs, which tend to break the mass degeneracy and cause mixing to become non-maximal . However, these difficulties may be avoided in some models of neutrino flavour symmetries .
As was discussed here by Smirnov and is seen in Fig. 5, the best fit to the Super-Kamiokande atmospheric-neutrino data is with near-maximal mixing and $`\mathrm{\Delta }m^23.7\times 10^3`$ eV, the 90 % confidence-level range being (2 to 7)$`\times 10^3`$ eV<sup>2</sup> . This is compatible with the 90 % confidence-level ranges favoured by the other atmospheric-neutrino experiments (Kamiokande, MACRO and Soudan). As for the solar-neutrino data, there have been four favoured regions, three of which shown in Fig. 6: the large- and small-mixing-angle MSW solutions (LMA and SMA) with $`\mathrm{\Delta }m^210^5`$ eV<sup>2</sup>, another MSW solution with lower $`\mathrm{\Delta }m^210^7`$ eV<sup>2</sup> (LOW), and vacuum solutions (VAC) with $`\mathrm{\Delta }m^210^9`$ to $`10^{11}`$ eV<sup>2</sup> .
How may one discriminate between the different oscillation scenarios? We know already from Super-Kamiokande that $`\nu _\mu \nu _e`$ atmospheric oscillations are not dominant, and more stringent upper limits are imposed by the Chooz and Palo Verde reactor experiments . However, $`\nu _\mu \nu _e`$ oscillations may be present at a lower level: if so, they may enable the sign of $`\mathrm{\Delta }m^2`$ to be determined via matter effects. There are several ways to distinguish between dominance by $`\nu _\mu \nu _\tau `$ or $`\nu _\mu \nu _s`$. The zenith-angle distributions in various categories of Super-Kamiolande events (high-energy partially-contained events, throughgoing muons and neutral-current enriched events) favour $`\nu _\mu \nu _\tau `$ oscillations, now at the 99 % confidence level . Another way to distinguish $`\nu _\mu \nu _\tau `$ from $`\nu _\mu \nu _s`$ is $`\pi ^0`$ production in Super-Kamiokande, which is impossible in $`\nu _s`$ interactions. The present data appear to favour $`\nu _\mu \nu _\tau `$, but, before reaching a conclusion, it will be necessary to reduce the systematic error in the production rate, which should be possible using data from the K2K near detector.
As for solar neutrinos, no significant distortion in the energy spectrum is now observed by Super-Kamiokande , and there is an upper limit on the flux of hep neutrinos that excludes the possibility that they might influence oscillation interpretations. Moreover, the day-night effect is now apparent only at the 1.3-$`\sigma `$ level, and the seasonal variation that seems to be emerging is completely consistent with the expected geometric effect of the Earth’s orbital eccentricity. These observations are all consistent with the LMA interpretation if $`\mathrm{\Delta }m^2>2\times 10^3`$ eV<sup>2</sup>, but the constraints from the day and night spectra now disfavour the SMA and VAC interpretations at the 95 % confidence level . Moreover, $`\nu _e\nu _s`$ is also disfavoured at the same level. Super-Kamiokande seems to be pushing us towards the LMA $`\nu _\mu \nu _\tau `$ interpretation, but has not yet provided a ‘smoking gun’ for this interpretation.
In the near future, SNO will be providing information on the neutral-current/charged-curr- ent ratio, discriminating between the different solar-neutrino scenarios and telling us whether the $`\nu _e`$ have oscillated into $`\nu _\mu /\nu _\tau `$ or $`\nu _s`$ . In the longer term, BOREXINO will check the disappearance of the <sup>7</sup>Be solar neutrinos .
Many of the most important prospective developments may be provided by long-baseline terrestrial neutrino experiments. The first of these is K2K, which has already released some preliminary results . They see 17 fully-contained events in their fiducial region, whereas 29.2$`{}_{3.3}{}^{}{}_{}{}^{+3.5}`$ would have been expected in the absence of oscillations, versus 19.3$`{}_{2.4}{}^{}{}_{}{}^{+2.5}`$ if $`\delta m^2=3\times 10^3`$. Overall, they see a total of 44 events, whereas about 74 (50) would have been expected in all event categories in the absence of oscillations (if $`\delta m^2=3\times 10^3`$), A detailed fit including data from the run currently underway and energy-spectrum information is now being prepared, but the present data already appear to be very promising! In a few years’ time, KamLAND will use reactor neutrinos to check the LMA MSW solution to the solar-neutrino deficit . Starting probably in 2003, MINOS will be looking for $`\nu _\mu `$ disappearance, measuring the neutral-current/charged-current ratio, and looking for $`\nu _e`$ appearance .
Then, starting in 2005, the CERN-Gran Sasso project will provide the opportunity to look for $`\nu _\mu \nu _\tau `$ oscillations directly via $`\tau `$ production . This seems to me a key experiment: “If you have not seen the body, you have not proven the crime.” The beam energy has been optimized for $`\tau `$ production, and either of the proposed experiments (OPERA and ICANOE ) should be able to discover $`\tau `$ production at the 4-$`\sigma `$ level if the atmospheric-neutrino parameters are in the range favoured by Super-Kamiokande, as exemplified in Fig. 7. Additionally, ICANOE may be able to probe the LMA MSW solar solution via low-energy atmospheric-neutrino events.
## 5 Supersymmetry
As you know, the motivation for supersymmetry at accessible energies is the hierarchy problem , namely why $`m_Wm_P`$ , or equivalently why $`G_FG_N`$, or equivalently why the Coulomb potential dominates over the Newton potential in atoms. Even if the hierarchy is set by hand, an issue is raised by the quantum corrections:
$$\delta m_{H,W}^2=𝒪\left(\frac{\alpha }{\pi }\right)\mathrm{\Lambda }^2m_W^2$$
(13)
where $`\mathrm{\Lambda }`$ is a cutoff reflecting the appearance of new physics beyond the Standard Model. Supersymmetry is a favoured example, making the corrections (13) naturally small:
$$\delta m_{H,W}^2=𝒪\left(\frac{\alpha }{\pi }\right)\left|m_B^2m_F^2\right|$$
(14)
which is $`<m_{H,W}^2`$ if
$$|m_B^2m_F^2|<1\mathrm{TeV}^2$$
(15)
for the difference in mass-squared between spartners. Circumstantial evidence for supersymmetry is provided by the concordance between the gauge coupling strengths measured at LEP and elsewhere with supersymmetric GUTs , and the indirect LEP indications for a light Higgs boson , as predicted by supersymmetry and discussed shortly. However, neither these nor the naturalness/fine-tuning argument set rigorous upper bounds on the sparticle mass scale .
Another argument favouring low sparticle masses is provided by cold dark matter . The lightest sparticle is commonly expected to be the lightest eigenstate of the neutralino mass matrix:
$$\left(\begin{array}{cccc}M_2& 0& \frac{g_2v_2}{\sqrt{2}}& \frac{g_2v_1}{\sqrt{2}}\\ & & & \\ 0& M_1& \frac{g^{}v_2}{\sqrt{2}}& \frac{g^{}v_1}{\sqrt{2}}\\ & & & \\ \frac{g_2v_2}{\sqrt{2}}& \frac{g^{}v_1}{\sqrt{2}}& 0& \mu \\ & & & \\ \frac{g_2v_1}{\sqrt{2}}& \frac{g^{}v_1}{\sqrt{2}}& \mu & 0\end{array}\right)$$
(16)
The gaugino masses $`M_2,M_1`$ are commonly assumed to be equal at the GUT scale:
$$M_2=M_1m_{1/2}$$
(17)
and are then renormalized: $`M_2/M_1\alpha _2/\alpha _1`$ at lower scales. The scalar masses may also be universal at the GUT scale, in which case they are also renormalized: $`m_{0_i}^2=m_0^2+C_im_{1/2}`$, and ratio of Higgs v.e.v.’s is commonly denoted by $`\mathrm{tan}\beta v_2/v_1`$.
The LEP limits on neutralinos and charginos exclude the possibility that the lightest neutralino $`\chi `$ is an almost pure photino or Higgsino. If univerality is assumed, or if neutralinos constitute the bulk of the cold dark matter , a dominant $`U(1)`$ gaugino component is favoured, as shown in Fig. 8 . As discussed in , one must take into account the various constraints on the universal parameters $`(m_{1/2},m_0)`$ imposed by LEP, the absence of charged $`\stackrel{~}{\tau }`$ dark matter, the observed rate for $`bs\gamma `$ decay and the absence of a charge- and colour-breaking (CCB) vacuum. The allowed dark-matter region is stretched to large $`m_{1/2}`$ by neutralino-slepton coannihilation, which allow $`m_\chi <`$ 600 GeV .
The most direct experimental constraints on the supersymmetric parameter space are those for charginos, neutralinos and sleptons at LEP, which impose $`m_{\chi ^\pm }>`$ 100 GeV and $`m_{\stackrel{~}{e}}>`$ 90 GeV generically. The Tevatron constraints on squarks and gluinos are of less direct importance if universality is assumed. However, stop searches at LEP and the Tevatron are important for constraining the radiative corrections to the lightest supersymmetric Higgs mass, whose leading terms contribute
$$\delta m_h^2=𝒪(\alpha )\frac{m_t^4}{m_W^2}\mathrm{ln}\left(\frac{m_{\stackrel{~}{t}}^2}{m_t^2}\right)$$
(18)
These are relevant to the constraints on $`m_0`$ and $`m_{1/2}`$ imposed by the absence of a Higgs boson at LEP .
Fig. 9 displays the lower limits on $`m_\chi `$ imposed by all these constraints, either if scalar-mass universality is (UHM) or is not (nUHM) assumed . In the UHM<sub>min</sub> scenario in Fig. 9, the absence of CCB vacua is not required. Fig. 9 also displays the expected impact of LEP searches in 2000, assuming pessimistically that they are unsuccessful. We see that
$$m_\chi >50\mathrm{GeV}\mathrm{and}\mathrm{tan}\beta >3,$$
(19)
with the precise values depending on the scenario adopted.
The above analysis assumed that CP violation could be ignored. However, as emphasized here by Kane , CP violation in the soft supersymmetry-breaking parameters may be important. Indeed, such CP violation is essential in electroweak baryogenesis <sup>1</sup><sup>1</sup>1The popular alternative of leptogenesis was discussed here by Ma ., which also requires a first-order phase transition, and hence a relatively light Higgs boson and stop squark, as seen in Fig. 10 ! LEP had been thought almost to exclude such a scenario because of its lower limit on $`m_h`$. However, we have recently emphasized that the LEP lower limit on $`m_h`$ may be greatly relaxed in the presence of CP violation. For example, as seen in Fig. 11, the h-H-A mixing induced by CP violation may suppress the hZZ coupling, and the $`h\overline{b}b`$ coupling may also be suppressed, as seen in Fig. 12 . We are currently re-evaluating the LEP lower limit on $`m_h`$, taking these effects into account.
## 6 Opportunities @ Future Accelerators
As discussed here by Fernandez , the LHC is under active construction, and is scheduled to start operating in 2005. In this talk, I concentrate on two selected LHC physics topics, namely the quest for the Holy Higgs, and the search for supersymmetry. At low masses, the $`H\overline{b}b`$ and $`\gamma \gamma `$ decay signatures look the most promising, with $`H4l^\pm `$ over a large range of intermediate masses, and $`HW^+W^{}l^+\nu l^{}\overline{\nu },l^\pm \nu jj`$ and $`HZZl^+l^{}\overline{\nu }\nu `$ interesting for high masses. As seen in Fig. 13 , there are no holes in the mass coverage, a couple of decay modes can normally be observed for any mass, and the Higgs mass can typically be measured with a precision $`10^3<\mathrm{\Delta }m_H/m_H<10^2`$. The LHC will also be able to discover supersymmetric Higgs bosons in two or more channels, over all the supersymmetric parameter space.
The LHC will produce principally strongly-interacting sparticles, squarks $`\stackrel{~}{q}`$ and gluinos $`\stackrel{~}{g}`$, and they sould be detectable if they weigh $`<`$ 2 TeV, as seen in Fig. 14 . This will enable the LHC to cover the parameter range allowed if the lightest supersymmetric particle provides the cold dark matter. The $`\stackrel{~}{g}`$ and $`\stackrel{~}{q}`$ often decay via complicated cascades, e.g., $`\stackrel{~}{g}\stackrel{~}{b}\overline{b},\stackrel{~}{b}\chi ^{}b,\chi ^{}\chi l^+l^{}`$, which may be reconstructed to provide some detailed mass measurements, as indicated in Fig. 15 .
A plausible scenario for physics after the LHC is that the Higgs will have been discovered, and one or two decay modes observed, and that several sparticles will have been discovered, but that heavier Higgses, charginos and sleptons may have escaped detection.
These lacunae provide some of the motivation for $`e^+e^{}`$ linear-collider (LC) physics. The very clean experimental environment, the egalitarian production of new weakly-interacting particles and the prospective availability of polarization make such a LC rather complementary to the LHC . One of the big issues is what energy to choose for a first-generation LC: we know there is the $`\overline{t}t`$ threshold at $`E_{cm}`$ 350 GeV and we believe there should be a ZH threshold at $`E_{cm}=m_Z+m_H<`$ 300 GeV. If one is above threshold, detailed studies of many Higgs decay modes (as seen in Fig. 16 ), or measurements of sparticle masses, become possible. However, we do not know what the supersymmetry threshold might be (assuming there is one!). For this reason, I think it is essential to retain as much flexibility as possible in the LC running energy.
In the meantime, the suggestion that the cold dark matter might consist of supersymmetric particles can be used to guess the likelihood that a LC of given energy might find supersymmetry. As seen in Fig. 17, we find that all the dark matter parameter space can be explored if $`E_{cm}`$ 1.25 TeV, and about 90 % if $`E_{cm}`$ = 1 TeV, but that an LC with $`E_{cm}`$ = 0.5 TeV would only cover about 60 % of the dark matter parameter region.
My opinion is that physics will demand an LC in the TeV energy range, and I hope that the world can converge on a (single) project in this energy range. In the rest of this talk, I assume that an LC with $`E_{cm}`$ 1 TeV will be built somewhere in the world, and ask what other accelerator projects might be interesting .
One suggestion is a future larger hadron collider with 100 TeV $`<E_{cm}<`$ 200 TeV, that could explore the 10 TeV mass region for the first time, if its luminosity rises to $`10^{35}`$ cm<sup>-2</sup> s<sup>-1</sup> or so. Such a machine is probably technically feasible, but it would be enormous, with a circumference of 100 to 500 km. The principal challenge will be reducing the unit cost by an order of magnitude compared to the LHC. At present, we cannot formulate the physics questions for such a machine with great clarity.
Another possibility is a higher-energy LC with $`E_{cm}>`$ 2 TeV, capable, e.g., of making precise and complete studies of sparticle spectroscopy, or of any other electroweakly-interacting sparticles weighing $`<`$ 1 TeV. CERN is developing a potential technology for such an LC, called CLIC , in which an intense low-energy drive beam is used to provide RF to accelerate a more energetic but less intense colliding beam. Accelerating gradients in the range 100 to 200 MeV $`m^1`$ appear possible, enabling an LC with $`E_{cm}<`$ 5 TeV to be accommodated in a tunnel $``$ 35 km long. The first physics study for such a machine was made at La Thuile in 1987 , and a new physics study has now been initiated . Fig. 18 is a first result from this new study, showing what a $`Z^{}`$ resonance might look like at CLIC .
A third type of accelerator currently attracting much interest is a complex of muon storage rings, as illustrated in Fig. 19. These could be developed in three steps : first a neutrino factory in which muons are simply allowed to decay without colliding, secondly a Higgs factory colliding $`\mu ^+\mu ^{}`$ at $`E_{cm}m_H`$, and thirdly a high-energy collider which might be compared with CLIC as a device to probe the high-energy frontier. The chief advantages of neutrino beams from $`\mu `$ decays, as opposed to conventional beams from hadron decays, are their precisely calculable fluxes and spectra, and the facts that equal numbers of $`\nu _\mu `$ and $`\overline{\nu }_e`$ (or $`\overline{\nu }_\mu `$ and $`\nu _e`$) are produced. A $`\mu ^+\mu ^{}`$ collider used as a Higgs factory can measure very precisely the mass and width of any Higgs boson with mass around 100 GeV, distinguishing between the Standard Model and a superymmetric extension, and measurements of supersymmetric Higgs peaks could provide a unique window on CP violation. A $`\mu ^+\mu ^{}`$ collider at the high-energy frontier has advantages over an $`e^+e^{}`$ LC with similar energy, conferred by the more precise energy calibration and reduced energy spread, but the ultimate energy is limited by the neutrino radiation hazard. However, although it is very attractive, many technical problems need to be solved before the feasibility of such a muon storage rign complex can be established.
The basic concept for a neutrino factory involves a proton driver with beam power 1 to 20 MW, provided by either a linac or a rapid-cycling synchrotron. This is used to produce pions, which decay into muons, of which about 0.1/proton are cooled, accelerated to 10 to 50 GeV, and stored in a ring. This need not be circular, and may look more like a bent paper-clip, as seen in Fig. 20 , with several straight sections sending $`(10^{20}`$ to $`10^{21})\overline{\nu }_\mu ,\nu _e`$ per year each towards detectors at different distances.
In long-baseline neutrino experiments with a neutrino factory, the sensitivities to mixing angles and $`\mathrm{\Delta }m^2`$ depend on $`E_\mu `$ and the distance $`L`$ . As we see in Fig. 21, a $`\nu _\mu \nu _e`$ appearance experiment would be much more sensitive than the present Super-Kamiokande upper limit , or what may be achieved with MINOS. Moreover, as seen in Fig. 22, with sufficiently many $`\mu `$ decays one may be sensitive to CP violation and matter effects (which depend on the sign of $`\mathrm{\Delta }m^2`$) in neutrino oscillations. For this, a detector at a distance of 2000 to 5000 km would be particularly advantageous. Ultimately, one could imagine a “World-Wide Neutrino Web” consisting of a $`\nu `$ factory in one region of the world feeding detectors in the same and other regions, as illustrated in Fig. 23.
Turning now to a Higgs factory, in the absence of a beam energy spread, the line shape (see Fig. 24) should be
$$\sigma _H(s)=\frac{4\pi \mathrm{\Gamma }(H\mu ^+\mu ^{})\mathrm{\Gamma }(HX)}{(sm_H^2)^2+m_H^2\mathrm{\Gamma }_H^2}$$
(20)
It seems that it might be possible to reduce the beam energy spread to $``$ 0.01 % or 10 MeV, comparable to the natural width of 3 MeV for a Standard Model Higgs weighing about 100 GeV. Calibrating the beam energy via the $`\mu ^\pm `$ polarization, it should then be possible to measure $`m_H`$ with a precision of $`\pm `$ 0.1 MeV, and the width to within 0.5 MeV, sufficient to distinguish a Standard Model Higgs boson from the lightest supersymmetric Higgs boson, over a large range of parameter space . In the supersymmetric case, a second-generation Higgs factory able to explore the “Twin Peaks” of the $`H`$ and $`A`$ shown in Fig. 25 might be even more interesting, providing tests of CP symmetry analogous to those in the $`K^0\overline{K}^0`$ system .
At the high-energy frontier, $`\mu ^+\mu ^{}`$ colliders would benefit from the absence of beamstrahlung and reduced initial-state radiation, as compared to an $`e^+e^{}`$ LC such as CLIC. However, the latter offers controllable beam polarization, $`e\gamma `$ and $`\gamma \gamma `$ colliders “for free”, and avoids the problems presented by $`\mu `$ decays. Moreover, an $`e^+e^{}`$ LC demonstrator, namely the SLC, has been built, whereas many of the technologies needed for a $`\mu ^+\mu ^{}`$ collider are at best glints in the eye, at present.
## 7 Towards a Theory of Everything?
The job description of a Theory of Everything (TOE) is to unify all the fundamental interactions, including gravity, and to solve all the problems that arise in attempts to quantize gravitation. A possible solution is to replace point-like elementary particles by extended objects, and the first incarnation of this idea used one-dimensional closed loops of string. It was soon realized that this scenario requires extra space-time dimensions 10 = 4 + 6 in the supersymmetric case, and/or extra interactions. The current reincarnation of this idea as $`M`$ theory includes other extended objects such as two-dimensional membranes, solids, etc. .
The key question is how to test these ideas. A popular suggestion is that the 6 surplus space dimensions are compactified on a Calabi-Yau manifold, but which one? 473 800 776 are known ! Recently we have embarked on a systematic study of Calabi-Yau (CY) spaces, constructed as zeroes of polynomials in weighted projective spaces, a technique which enables one to explore some of their internal properties and focus on these with desirable features . Our harvest so far comprises 182 737 CY spaces, of which 211 have 3 generations and K3 fibrations as desired in some approaches to $`M`$ theory .
In the face of this ambiguity, how can one speak of phenomenological predictions from string theory? In fact, it has told us correctly that there cannot be more than 10 dimensions, that the gauge group cannot be very large, that matter representations cannot be very big, and that the top quark should not weigh more than about 190 GeV. It has also provided a first-principles estimate of the unification scale, $`M_U`$ few $`\times 10^{17}`$ GeV , which is not so far from the phenomenological bottom-up estimate of (1 or 2) $`\times 10^{16}`$ GeV.
Moreover, we now understand that all the different string theories are related by dualities in the general framework known as $`M`$ theory . The question then becomes, in which part of its parameter space do we live? The GUT mass-scale calculation suggests that the string coupling may be strong , corresponding to one large dimension: $`L1/M_{GUT}1/10^{16}`$ GeV $`1/m_Pl_P`$.
Adventurous souls have then gone on to propose that one or more “small” dimensions might actually be rather large, perhaps $`L1`$ TeV<sup>-1</sup> or even $``$ 1 mm ? In such models, there may be observable modifications of Newton’s law: $`G_N/rG_N(1/r)(L/r)^{\mathrm{\Delta }D}`$ at short distances, as well as possible new accelerator signatures. This suggestion offers plenty of phenomenological fun, but why should the scale of gravity sink so low? Are there any advantages for the hierarchy problem? So far, I have seen it reformulated, but not yet solved.
Before closing, I would like to mention a couple of radical possibilities for string phenomenology. As we heard here, surprisingly many ultra-high-energy (UHE) cosmic rays have been observed above $`E5\times 10^{19}`$ eV, more than expected above the GZK cutoff due to the reaction $`p+\gamma _{CMB}\mathrm{\Delta }^+`$ . Unless one modifies relativistic kinematics (see later) these UHE cosmic rays should have originated nearby, at distances $`d<`$ 100 Mpc for $`E10^{20}`$ eV, but no discrete sources have yet been confirmed.
Might they originate from the decays of supermassive dark matter particles? It has recently been realized that such particles weighing 10<sup>10</sup> GeV or more might have been produced by non-thermal mechanisms early in the history of the Universe . Possible candidates for these unstable heavy relics can be found in string theory, particularly as bound states in the hidden sector, which we have termed cryptons . They could well have masses $`10^{13}`$ GeV and be metastable, decaying via higher-dimensional operators into multiple leptons and quarks. Simulations indicate that they could well produce the observed UHE tail of the cosmic-ray spectrum, as seen in Fig. 26 . The Auger project should be able to tell us whether this mechanism is tenable .
Even more speculative is the suggestion of quantum-gravitational phenomenology. Might the space-time foam of quantum-gravitational fluctuations in the fabric of space induce quantum decoherence and/or CPT violation at the microscopic scale ? Here the most sensitive probe may be the $`K^0\overline{K}^0`$ system . Might the velocity of light (or a neutrino ) depend on its energy, because of recoil effects on the space-time vacuum? Here the most sensitive direct tests may be provided by distant, energetic sources with short time-scales, such as gamma-ray bursters (GRBs) , active galactic nuclei and pulsars , and some such models are also constrained by the kinematics of UHE cosmic rays . Fig. 27 shows fits to BATSE data on GRB 970508 in different energy channels. A regression analysis of fits to GRBs with measured redshifts has been used to constrain any possible energy dependence of the velocity of light: $`\delta c/cE/M:M>10^{15}`$ GeV .
## 8 Final Comments
The history of physics reveals many ways in which it may advance, being driven either by pure theoretical thought or by experimental breakthroughs. Pure theoretical speculation must in any case be tempered by experimental reality: we can never forget that in physics, as any other science, experiment is the ultimate arbiter. Particle physics is currently fortunate. On the one hand, experiments at LEP and elsewhere have shown that the Standard Model is a solid rock on which to build. On the other hand, experiments on neutrinos strongly indicate oscillations, and hence physics beyond the Standard Model. There are exciting new experimental programmes underway to explore the flavour problem, to pin down models of neutrino oscillations, and to explore the TeV energy range.
Beyond our daily concerns, we have the great unsolved problem of twentieth-century physics left to stimulate us: reconcile gravity and quantum mechanics. Great theoretical advances towards this goal have been made during recent years, but we do not know how far we are from this goal. In particular, we have not yet defined a clear experimental test that will confirm or refute string theory. Finding it is our key phenomenological challenge.
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# ■("DSF 25/2000, astro-ph/0007419") Testing Standard and Degenerate Big Bang Nucleosynthesis with BOOMERanG and MAXIMA-1
\[
## Abstract
We test different Big Bang Nucleosynthesis scenarios using the recent results on the Cosmic Microwave Background Anisotropies provided by the BOOMERanG and MAXIMA-1 experiments versus the observed abundances of $`{}_{}{}^{4}He`$, $`D`$ and $`{}_{}{}^{7}Li`$. The likelihood analysis, based on Bayesian approach, shows that, in the case of high deuterium abundance, $`Y_D=(2.0\pm 0.5)10^4`$, both standard and degenerate BBN are inconsistent with the CMBR measurement at more than $`3\sigma `$. Assuming low deuterium abundance, $`Y_D=(3.4\pm 0.3)10^5`$, the standard BBN model is still inconsistent with present observations at $`2\sigma `$ level, while the degenerate BBN results to be compatible. Unless systematics effects will be found in nuclide abundances and/or in CMBR data analysis this result may be a signal in favour of new physics like a large chemical potential of the relic neutrino-antineutrino background.
\] One of the main goals of modern cosmology is the knowledge of the energy density content of the universe. The four parameters $`\mathrm{\Omega }_B`$, $`\mathrm{\Omega }_{CDM}`$, $`\mathrm{\Omega }_\nu `$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$, giving, respectively, the baryon, cold dark matter, neutrino and cosmological constant contributions to the total energy density, in unit of the critical density, and the Hubble constant, $`H_0=100hKms^1Mpc^1`$, enter as crucial parameters in several cosmological observables. A well known example is provided by the baryon density parameter which plays an essential role in determining the abundances of light nuclides produced in the early universe. An increasing interest has been also devoted to other issues, as structure formation and the anisotropy of the Cosmic Microwave Background Radiation (CMBR). As for our theoretical understanding of the Big Bang Nucleosynthesis (BBN), to test the theoretical models which describe these aspects of the hot Big Bang model, it is essential to have precise measurements of the several $`\mathrm{\Omega }_ih^2`$. Furthermore, combining different cosmological observables, and comparing the way they are able to constrain the $`\mathrm{\Omega }_ih^2`$, allows for a check of the consistency of our present understanding of the evolution of the universe. This can provide new hints on phenomena which took place at the macroscopic cosmological level, or rather related with the very microscopic structure of fundamental interactions.
This letter represents a contribution in this direction. In particular, we do perform a combined analysis of the dependence on the energy fractions $`\mathrm{\Omega }_Bh^2`$ and $`\mathrm{\Omega }_\nu h^2=N_\nu \mathrm{\Omega }_\nu ^0h^2`$ for massless neutrinos ($`N_\nu `$ standing for the effective neutrino number and $`\mathrm{\Omega }_\nu ^0h^2`$ for the energy contribution of a single $`\nu \overline{\nu }`$ specie) of CMBR anisotropies and BBN. This is aimed to test the standard and degenerate BBN scenario, using the recent results of the BOOMERanG and MAXIMA-1 CMBR experiments and the measurements of $`{}_{}{}^{4}He`$, $`D`$ and $`{}_{}{}^{7}Li`$ primordial abundances.
The theoretical tools necessary to achieve this goal are nowadays rather robust and the accuracy in the predictions of the Big Bang model is remarkably improved. In fact, recently, a new generation of BBN codes have been developed , which give the $`{}_{}{}^{4}He`$ mass fraction $`Y_p`$ with a theoretical error of few per mille. This effort is justified in view of the small statistical error which is now quoted in the $`Y_p`$ measurements. On the other hand, the theoretical predictions on the CMBR anisotropies angular power spectrum, in the case of primordial passive and coherent perturbations, expected in the most general inflationary model, have also recently reached a $`1\%`$ level accuracy, with a refined treatment of H, He I, and He II recombination.
Concerning the CMBR experimental data, after the COBE satellite first detection of CMBR anisotropies, at scales larger than $`5^o`$, and more than $`20`$ independent detections at different frequencies and scales, see e.g. , an important new insight is represented by the recent results obtained by the BOOMERanG Collaboration . For the first time, in fact, multifrequency maps of the microwave background anisotropies were realized over a significant part of the sky, with $`10^{}`$ resolution and high signal to noise ratio. The anisotropy power spectrum, $`C_{\mathrm{}}`$, was measured in a wide range of angular scales from multipole $`\mathrm{}50`$ up to $`\mathrm{}600`$, with error bars of the order of $`10\%`$, showing a peak at $`\mathrm{}_{peak}=(197\pm 6)`$ with an amplitude $`DT_{200}=(69\pm 8)\mu K`$. While the presence of such peak, compatible with inflationary scenario, was already suggested by previous measurements , the absence of secondary peaks after $`\mathrm{}300`$ with a flat spectrum with an amplitude of $`40\mu K`$ up to $`\mathrm{}625`$ was a new and unexpected result. This result obtained then an impressive confirmation by the MAXIMA-1 experiment up to $`\mathrm{}800`$.
As far as BBN is concerned, in the last few years many results have been also obtained on light element primordial abundances. The $`{}_{}{}^{4}He`$ mass fraction $`Y_p`$, has been measured with a $`0.1\%`$ precision in two independent surveys, from regression to zero metallicity in Blue Compact Galaxies, giving a low value $`Y_p^{(l)}=0.234\pm 0.003`$ , and a high one $`Y_p^{(h)}=0.244\pm 0.002`$ , which are compatible at $`2\sigma `$ level only, may be due to large systematic errors. As in , in our analysis we adopt the more conservative value $`Y_p=0.238\pm 0.005`$.
A similar controversy holds in $`D`$ measurements, where observations in different Quasars Absorption line Systems (QAS) lead to the incompatible results $`Y_D^{(l)}=\left(3.4\pm 0.3\right)10^5`$ , and $`Y_D^{(h)}=\left(2.0\pm 0.5\right)10^4`$ . We will perform our analysis for both low and high $`D`$ data. Finally, the most recent estimate for $`{}_{}{}^{7}Li`$ primordial abundance, from the Spite plateau observed in the halo of POP II stars, gives $`Y_{{}_{}{}^{7}Li}=\left(1.73\pm 0.21\right)10^{10}`$ . The light nuclide yields strongly depend on the baryon matter content of the universe, $`\mathrm{\Omega }_Bh^2`$. High values for this parameter result in a larger $`{}_{}{}^{4}He`$ mass fraction and a lower deuterium number density. In particular, assuming a standard BBN scenario, i.e. vanishing neutrino chemical potentials, the likelihood analysis gives, at $`95\%`$ C.L. ,
$$\begin{array}{ccc}\text{low D}\hfill & \mathrm{\Omega }_Bh^2=.017\pm .003& 1.7N_\nu 3.3,\\ \text{high D}\hfill & \mathrm{\Omega }_Bh^2=.007_{.002}^{+.007}& 2.3N_\nu 4.4.\end{array}$$
(1)
As already pointed out by several authors , these values for $`\mathrm{\Omega }_Bh^2`$, though in the correct order of magnitude, are however somehow smaller than the baryon fraction which more easily fit the CMBR data. In fact, as mentioned, while the narrow first peak around $`l200`$ is a confirmation of the inflationary paradigma, a flat universe with adiabatic perturbations, the lack of observation of a secondary peak at smaller scales raises new intriguing questions about the values of the cosmological parameters. In particular, this result may be a signal in favour of a larger $`\mathrm{\Omega }_Bh^20.03`$, since increasing the baryon fraction enhances the odd peaks only. In this respect the measurement of the third peak at larger multipole moments is extremely important. Though it is fair to say that a further analysis of the BOOMERanG and MAXIMA-1 data, as well as new data from future experiment, are needed to clarify this issue, it is however timely to study how our theoretical understanding of BBN can be reconciled with the larger values of $`\mathrm{\Omega }_Bh^2`$ suggested by CMBR data.
In figure 1 we start addressing the problem quantitatively, by reporting our theoretical predictions on the several abundances, together with the corresponding experimental values, and the CMBR 68$`\%`$ C.L. bound on $`\mathrm{\Omega }_Bh^2=0.030_{0.004}^{+0.007}`$, obtained by our joint analysis of both BOOMERanG and MAXIMA-1 data (see below for the details). The BBN predictions are obtained in the standard scenario for $`N_\nu =3`$.
Unless systematics effects will be found in nuclide abundances and/or in CMBR data analysis, we see that there is a rather sensible disagreement between the values of the baryon energy density required by standard BBN and by CMBR data. In particular, from the 68$`\%`$ C.L. CMBR result for $`\mathrm{\Omega }_Bh^2`$, we can derive the following constraints on the $`Y_i`$,
$`0.2487<`$ $`Y_{{}_{}{}^{4}He}`$ $`<0.2536`$ (2)
$`9.2310^6<`$ $`Y_D`$ $`<2.3210^5`$ (3)
$`7.9210^6<`$ $`Y_{{}_{}{}^{3}He}`$ $`<9.8510^6`$ (4)
$`5.4510^{10}<`$ $`Y_{{}_{}{}^{7}Li}`$ $`<1.2510^9,`$ (5)
which are systematically different from the available measured values. We stress once again that our theoretical analysis is accurate, at worst, at the level of few percent<sup>*</sup><sup>*</sup>*It is important to note that this result is obtained following a Bayesian approach in the analysis. Assuming the standard BBN, the best fit model for the CMBR spectrum has a normalized $`\chi ^250`$, over $`40`$ degrees of freedom. Being the $`\chi ^2`$ slightly greater than one, a CMB analysis using the BBN prior as in is still perfectly consistent..
It has already been stressed that a simple way to improve the agreement of observed nuclide abundances with $`\mathrm{\Omega }_Bh^20.02`$ is to assume non vanishing neutrino chemical potentials at the BBN epoch, a scenario already extensively studied in the past . The effect of neutrino chemical potentials $`\mu _\alpha `$, with $`\alpha `$ the neutrino specie, is twofold. A non-vanishing $`\xi _\alpha =\mu _{\nu _\alpha }/T_\nu `$, contribute to $`N_\nu `$ as
$$N_\nu =3+\mathrm{\Sigma }_\alpha \left[\frac{30}{7}\left(\frac{\xi _\alpha }{\pi }\right)^2+\frac{15}{7}\left(\frac{\xi _\alpha }{\pi }\right)^4\right],$$
(6)
implying a larger expansion rate of the universe with respect to the non-degenerate scenario, and a higher value for the neutron to proton density ratio at the freeze-out. Furthermore, a positive (negative) value for $`\xi _e`$ means a larger (smaller) number of $`\nu _e`$ with respect to $`\overline{\nu }_e`$, thus enhancing (lowering) $`np`$ processes. Notice that extra relativistic degrees of freedom, like light sterile neutrinos, would contribute to $`N_\nu `$ as well, and in this respect BBN cannot distinguish between their contribution to the total universe expansion rate and the one due to neutrino degeneracy. Therefore our estimates for $`N_\nu `$ can only represent an upper bound for the total neutrino chemical potentials. However, it is also worth stressing that a large $`N_\nu >3`$, $`does`$ indeed require a positive $`\xi _e`$, since to end up with the observed values for the $`Y_i`$, a fine tuned balance between $`\xi _e`$ and $`N_\nu `$ is required .
Increasing $`N_\nu `$ also weakly affects the CMBR anisotropy spectrum in two ways. The growth of perturbations inside the horizon is in fact lowered, resulting in a decay of the gravitational potential and hence in an increase of the anisotropy near the first peak. Moreover, the size of horizon and sound horizon at the last scattering surface is changed, and this, with additional effects in the damping, varies the amplitude and position of the other peaks, see e.g. .
To test the degenerate BBN scenario we have performed a likelihood analysis of the data. First, to constrain the values of the parameter set $`(\xi _e`$, $`N_\nu `$, $`\mathrm{\Omega }_Bh^2)`$ from the data on $`{}_{}{}^{4}He`$, $`D`$ and $`{}_{}{}^{7}Li`$ we define a total likelihood function,
$$_{Nucl}(N_\nu ,\mathrm{\Omega }_Bh^2,\xi _e)=L_DL_{{}_{}{}^{4}He}L_{{}_{}{}^{7}Li},$$
(7)
as described in Ref.. For a fully degenerate BBN, since the effect of a positive $`\xi _e`$ can be compensated by larger $`N_\nu `$, $`_{Nucl}(N_\nu ,\mathrm{\Omega }_Bh^2,\xi _e)`$ may sensibly differ from zero in a region with rather large values of $`N_\nu `$. We have chosen to constrain this parameter to be $`N_\nu <16`$. This upper limit has been considered after checking that it is well outside the $`95\%`$ upper limit on $`N_\nu `$ from the BOOMERanG and MAXIMA-1 data (again, see below). The other two parameters are chosen in the following ranges, $`1\xi _e1`$ and $`0.004\mathrm{\Omega }_Bh^20.110`$.
Since CMBR spectrum is not sensible to $`\xi _e`$ alone, we have marginalized over $`\xi _e`$. All nuclide abundances have been evaluated using the new BBN code described in , while the theoretical uncertainties $`\sigma _i^{th\mathrm{\hspace{0.17em}2}}(N_\nu ,\mathrm{\Omega }_Bh^2)`$ are found by linear propagation of the errors affecting the various nuclear rates entering in the nucleosynthesis reaction network .
The CMBR data analysis methods have been already extensively described in various papers . The anisotropy power spectrum, was estimated in $`12`$ orthogonalized (independent) bins between $`\mathrm{}=50`$ and $`\mathrm{}=650`$ by the BOOMERanG experiment and in $`10`$ bins, from $`\mathrm{}=70`$ to $`\mathrm{}=750`$, by the MAXIMA-1 experiment. The likelihood function of the CMBR signal $`C_B`$ inside the bins is well approximated by an offset lognormal distribution, such that the quantity $`D_B=ln(C_B+x_B)`$ (where $`x_B`$ is the offset correction) is a gaussian variable. The likelihood for a given cosmological model is then defined by $`2\mathrm{l}\mathrm{n}L=(D_B^{th}D_B^{ex})M_{BB^{}}(D_B^{}^{th}D_B^{}^{ex})`$, where $`M_{BB^{}}`$ is a matrix that defines the noise correlations and $`D_B^{th}`$ is the offset lognormal theoretical band power. The theoretical $`D_B^{th}`$ were generated using fast and accurate Boltzmann solvers . Our database of models is sampled in physical variables as in , but we only consider flat models, $`h=0.65\pm 0.2`$ and the effective number of neutrinos is allowed to vary up to $`N_\nu =20`$. Following we can therefore constrain the parameter of interest, $`\mathrm{\Omega }_Bh^2`$ and $`N_\nu `$, by finding the remaining “nuisance” parameters which maximize them. As mentioned, at 95 and 99$`\%`$ C.L. we found $`N_\nu 13`$ and $`N_\nu 16`$, respectively. A similar bound, using BOOMERanG data only, has been obtained in .
In figure 2 we have summarized the second main result of our analysis. In the $`\mathrm{\Omega }_Bh^2N_\nu `$ plane we show the 95$`\%`$ C.L. likelihood regions for both the high and low $`D`$ measurements, as well as the analogous contours for standard BBN, obtained running our code with $`\xi _e=0`$. In the same plot we show the 68 and 95 $`\%`$ C.L. regions obtained by CMBR data.
As a first comment, the standard BBN, $`\xi _e=0`$, and CMBR data analysis lead to quite different values for $`\mathrm{\Omega }_Bh^2`$. This can be clearly seen from the reported 95$`\%`$ results, but we have verified that the 99$`\%`$ C.L. contour for high $`D`$ has no overlap with the region picked up by BOOMERanG and MAXIMA-1 data, and a very marginal one for low $`D`$.
For the degenerate scenario, increasing $`N_\nu `$, the allowed intervals for $`\mathrm{\Omega }_Bh^2`$ shift towards larger values. However the high $`D`$ values require a baryon content of the universe energy density which is still too low, $`\mathrm{\Omega }_Bh^20.018`$, to be in agreement with CMBR results. A large overlap is instead obtained for the low $`D`$ case, whose preferred $`\mathrm{\Omega }_Bh^2`$ span the range $`0.012\mathrm{\Omega }_Bh^20.036`$. As expected, a larger $`N_\nu `$ helps in improving the agreement with the high CMBR $`\mathrm{\Omega }_Bh^2`$ value, but is important to stress that a large value for $`N_\nu `$ is not preferred by the CMBR data alone, being, in this case, the best fit $`N_\nu 3`$. If we only consider the $`95\%`$ overlap region we get the following conservative bounds:
$$4N_\nu 13,0.024\mathrm{\Omega }_Bh^20.034.$$
(8)
In this region $`\xi _e`$ varies in the range $`0.07\xi _e0.43`$. As we said, values $`N_\nu 3`$, as suggested from our analysis, can be either due to weak interacting neutrino degeneracy, or rather to other unknown relativistic degrees of freedom.
In conclusion, we have shown how a precision analysis of BBN and CMBR data is able to tell us about possible new features of the cosmological model describing the evolution of the universe. We have quantitatively discussed how larger values for the baryonic matter content $`\mathrm{\Omega }_Bh^2`$ may be reconciled with BBN predictions in a degenerate scenario. In this respect the observed low value of deuterium more easily fits with the constraints given by BOOMERanG and MAXIMA-1 data. The crucial aspect of our analysis is that this agreement is realized with an effective neutrino degrees of freedom larger than 4 at 95$`\%`$. This represents an indication in favour of a degenerate neutrino background and/or new particle species contributing to relativistic matter in the universe. As final remark, we should note that our CMB analysis was restricted on a specific class of models with a limited numbers of parameters. Including curvature, a gravity waves background, or removing our priors on $`h`$ would not move the peak of our likelihood on $`\mathrm{\Omega }_Bh^2`$, but would enlarge our C.L. more on the high value side . Including drastical deviations from the standard model like topological defects , decoherence in the primordial fluctuations or assuming a less general inflationary model would drastically change the conclusions of our work. Each of these effect leaves a characteristic imprint on CMB, so hopefully with new data available in the near future, as well as further analysis of the BOOMERanG and MAXIMA results, it will be possible to severely scrutinize this result.
One of the author, A. Melchiorri would like to thank P. de Bernardis, P. Ferreira, J. Silk and N. Vittorio for valuable comments and discussions; G. Miele would like to thank CERN TH-Division for support.
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# Collapsing Perfect Fluid in Higher Dimensional Spherical Spacetimes
## I Introduction
Recently, we have studied the gravitational collapse of perfect fluid in four-dimensional spacetimes , and found that when solutions have continuous self-similarity, the formation of black holes always starts with zero-mass, while when solutions have no such a symmetry it starts with a mass gap. The solutions with zero masses actually represent naked singularities. Thus, if the Cosmic Censorship Conjecture is correct , it seems that Nature prohibits the existence of solutions of the Einstein field equations with self-similarity.
Quite recently, there has been renewed interest in studying higher dimensional spacetimes from the point of view of both Cosmology and gravitational collapse. In particular, it was found that the exponent $`\gamma `$, appearing in the mass scaling form of black holes, depends on the dimensions of the spacetimes .
In this paper, we shall generalize our previous studies to the case of perfect fluid in N-dimensional spherical spacetimes and will show that our previous results for gravitational collapse, obtained in four-dimensional spacetimes , are also valid in N-dimensional spacetimes. The solutions to be presented below can be also considered as representing cosmological models. However, in this paper we shall not consider this possibility and leave it to another occasion. The rest of the paper is organized as follows: In Sec.$`II`$ we shall derive the general form of metric for conformally flat spherical spacetimes with N-dimensions. As an application, all the homogeneous and isotropic Friedmann-Robertson-Walker-like solutions for a perfect fluid with the equation of state $`p=\alpha \rho `$ are found for spacetimes with any dimensions. The solutions can be classified into three different families, flat, open, and close, depending on the curvature of the $`(N1)`$-dimensional spatial part of the metric, as that in four-dimensional case. In Sec. $`III`$, two families of these solutions are studied in the context of gravitational collapse, one has continuous self-similarity, and the other has neither continuous self-similarity nor discrete self-similarity. In order to study gravitational collapse in N-dimensional spacetimes, in this section a definition for the mass function is first given, whereby the location of the apparent horizons can be read off directly. This definition gives the correct results for the static case and reduces to the four-dimensional one when $`N=4`$ . The paper is ended with Sec. $`IV`$, where our main conclusions are derived.
It should be noted that multidimensional spacetimes where all dimensions are in the equal foot, like the ones to be considered here, are not so realistic, as we are living in an effectively 4-dimensional spacetime, so in principle one might expect that by dimensional reduction the multidimensional spacetimes should reduce to our 4-dimensional world. To have such reduction possible, the dimensions should have different weights in any realistic model. In this sense, the models considered in this paper are very ideal.
## II The General Homogeneous and Isotropic Solutions in N-dimensional Spacetimes
To start with, let us consider the general metric for N-dimensional spacetimes with spherical symmetry ,
$$ds^2=G(t,r)dt^2K(t,r)\left(dr^2+r^2d\mathrm{\Omega }_{N2}^2\right),$$
(1)
where $`\{x^\mu \}\{t,r,\theta ^2,\theta ^3,\mathrm{},\theta ^{N1}\}(\mu =0,1,2,\mathrm{},N1)`$ are the usual N-dimensional spherical coordinates, $`d\mathrm{\Omega }_{N2}^2`$ is the line element on the unit (N-2)-sphere, given by
$`d\mathrm{\Omega }_{N2}^2`$ $`=`$ $`\left(d\theta ^2\right)^2+\mathrm{sin}^2(\theta ^2)\left(d\theta ^3\right)^2+\mathrm{sin}^2(\theta ^2)\mathrm{sin}^2(\theta ^3)\left(d\theta ^4\right)^2`$ (3)
$`+\mathrm{}+\mathrm{sin}^2(\theta ^2)\mathrm{sin}^2(\theta ^3)\mathrm{}\mathrm{sin}^2(\theta ^{N2})\left(d\theta ^{N1}\right)^2`$
$`=`$ $`{\displaystyle \underset{i=2}{\overset{N1}{}}}\left[{\displaystyle \underset{j=2}{\overset{i1}{}}}\mathrm{sin}^2(\theta ^j)\right]\left(d\theta ^i\right)^2.`$ (4)
Then, it can be shown that the conformally flat condition $`C_{\mu \nu \lambda \delta }=0`$, where $`C_{\mu \nu \lambda \delta }`$ denotes the Weyl tensor, reduces to a single equation
$$D,_{rr}\frac{D_{,r}}{r}=0,$$
(5)
where $`D(G/K)^{1/2}`$, and $`(),_r()/r`$, etc. The general solution of the above equation is given by
$$D(t,r)=f_1(t)+f_2(t)r^2,$$
(6)
where $`f_1`$ and $`f_2`$ are two arbitrary functions of $`t`$ only. Using the freedom of coordinate transformations, it can be shown that there are essentially only two different cases, one is $`f_1(t)=1,f_2(t)=0`$, and the other is $`f_1(t)0,f_2(t)=1`$. Thus, it is concluded that the general conformally flat N-dimensional metric with spherically symmetry takes the form
$$ds^2=G(t,r)\left[dt^2h^2(t,r)\left(dr^2+r^2d\mathrm{\Omega }_{N2}^2\right)\right],$$
(7)
where
$$h(t,r)=\{\{\begin{array}{cc}1,\hfill & \\ \left[f_1(t)+r^2\right]^1,cr\hfill & \end{array}$$
(8)
with $`f_1(t)0`$. In the following we shall refer solutions with $`h(t,r)=1`$ as Type $`A`$ solutions, and solutions with $`h(t,r)=[f_1(t)+r^2]^1`$ as Type $`B`$ solutions.
When $`f_1(t)=Const.`$, say, $`f_1`$, we can introduce a new radial coordinate $`\overline{r}`$ via the relation
$$\overline{r}=\frac{r}{f_1+r^2},$$
(9)
then the metric (3) becomes
$$ds^2=G(t,r)(dt^2\frac{d\overline{r}^2}{14f_1\overline{r}^2}\overline{r}^2d\mathrm{\Omega }_{N2}^2),(f_1(t)=Const.).$$
(10)
If we further set $`G(t,r)=G(t)`$, the above metric becomes the Friedmann-Robertson-Walker (FRW) metric but in N-dimensional spacetimes. Solving the corresponding Einstein field equations for a perfect fluid with the equation of state $`p=\alpha \rho `$, we find two classes of solutions, where $`\rho `$ denotes the energy density of the fluid, $`p`$ the pressure, and $`\alpha `$ is an arbitrary constant. The details of the derivation of these solutions were given in , so in the following we should only present the final form of the solutions.
Type $`A`$ solutions. In this case, the general solutions are given by
$$h(t,r)=1,G(t,r)=\left(1Pt\right)^{2\xi },$$
(11)
where $`P`$ is a constant, which characterizes the strength of the spacetime curvature. In particular, when $`P=0`$, the spacetime becomes Minkowski. The constant $`\xi `$ is a function of $`\alpha `$ and the spacetime dimension $`N`$, given by
$$\xi \frac{2}{(N3)+\alpha (N1)}.$$
(12)
The corresponding energy density and velocity of the fluid are given, respectively, by
$`p`$ $`=`$ $`\alpha \rho =3\kappa ^1(N1)\alpha \xi ^2P^2\left(1Pt\right)^{2(\xi +1)},`$ (13)
$`u_0`$ $`=`$ $`\left(1Pt\right)^\xi ,u_i=0,(i=1,2,\mathrm{},N1).`$ (14)
This class of solutions belongs to the ones studied by Ivashchuk and Melnikov in the context of higher dimensional Cosmology . As a matter of fact, setting $`n=0`$ and $`N_0=N1`$ in their general solutions, we shall obtain the above ones. When $`\alpha =0,(N1)^1`$, the above solutions reduce, respectively, to the one for a dust and radiation fluid, which were also studied recently by Chatterjee and Bhui . The $`\alpha =0`$ case was studied in the context of gravitational collapse, too . It can be shown that the curvature of the (N-1)-dimensional spatial part of the metric $`t=Const.`$ in this case is zero, and the spacetimes correspond to the spatially flat FRW model.
Type $`B`$ solutions. In this case, the general solutions are given by,
$$h(t,r)=\frac{1}{f_1+r^2},G(t)=\left[A\mathrm{cosh}(\omega t)+B\mathrm{sinh}(\omega t)\right]^{2\xi },$$
(15)
where $`\omega 2\sqrt{f_1}/\xi `$, $`A`$ and $`B`$ are integration constants, and $`\xi `$ is defined by Eq.(12). The energy density and velocity of the fluid now are given by
$`p`$ $`=`$ $`\alpha \rho =12\kappa ^1N(N1)\alpha f_1(A^2B^2)\left[A\mathrm{cosh}(\omega t)+B\mathrm{sinh}(\omega t)\right]^{2(1+\xi )},`$ (16)
$`u_0`$ $`=`$ $`\left[A\mathrm{cosh}(\omega t)+B\mathrm{sinh}(\omega t)\right]^\xi ,u_i=0,(i=1,2,\mathrm{}N1).`$ (17)
It can be shown that the curvature of the (N-1)-dimensional spatial part of the metric $`t=Const.`$ in the present case is different from zero. In fact, when $`f_1>0`$, the curvature is positive, and the spacetime is spatially closed, and when $`f_1<0`$, the curvature is negative, and the spacetime is spatially open. The particular solution with $`\alpha =(N3)/(N1)`$ was found in , given by Eqs.(49) and (50), but its physical studies were excluded. As far as we know, the rest of the solutions are new.
It should be noted that the above solutions are valid for any constant $`\alpha `$. However, in the rest of the paper we shall consider only the case where $`(N2)^1\alpha 1`$, so that the three energy conditions, weak, dominant, and strong, are satisfied . When $`N=4`$, these solutions reduce to the FRW solutions, which have been studied in the context of gravitational collapse in . Therefore, in the following we shall assume that $`N4`$.
## III Gravitational Collapse of Perfect Fluid in N-dimensional Spacetimes
To study the above solutions in the context of gravitational collapse, we need first to define the local mass function. Recently, Chatterjee and Bhui gave a definition, in generalizing the Cahill and Macvittie mass function in four-dimensional spacetimes to N-dimensional spacetimes . However, in this paper we shall use the following definition for the mass function,
$$1\frac{2m(t,r)}{B_Nr_{ph}^{N3}}=g^{\mu \nu }r_{ph}^{}{}_{,\mu }{}^{}r_{ph}^{}{}_{,\nu }{}^{},$$
(18)
where
$$B_N=\frac{\kappa \mathrm{\Gamma }\left(\frac{N1}{2}\right)}{2(N2)\pi ^{(N1)/2}},$$
(19)
with $`\mathrm{\Gamma }`$ denoting the gamma function, and $`r_{ph}`$ the geometric radius of the (N-2)-unit sphere. It can be shown that this definition in general is different from that given by Chatterjee and Bhui , and reduces to the one usually used in four-dimensional spacetimes when $`N=4`$ , and yields the correct mass of black holes in N-dimensions for the static spherical spacetimes .
In the study of gravitational collapse, another important conception is the apparent horizon, the formation of which indicates the formation of black holes. The apparent horizon in the present case is defined as the outmost boundary of the trapped $`(N2)`$-spheres . The location of the trapped $`(N2)`$-spheres is the place where the outward normal of the surface, $`r_{ph}=Const.`$, is null, i.e.,
$$g^{\mu \nu }r_{ph}^{}{}_{,\mu }{}^{}r_{ph}^{}{}_{,\nu }{}^{}=0.$$
(20)
Then, the mass function on the apparent horizon is given by
$$M_{AH}=\frac{B_N}{2}r_{ph}^{N3}|_{r=r_{AH}},$$
(21)
where $`r=r_{AH}`$ is a solution of Eq.(20), which corresponds to the outmost trapped surface. In gravitational collapse $`M_{AH}`$ is usually taken as the mass of black holes . With the above definition for the mass function, let us study the main properties of the above two types of solutions separately.
### A Type A solutions
The mass function defined by Eq.(18) in this case takes the form
$$m(t,r)=\frac{B_N}{2}\frac{\xi ^2P^2r^{N1}}{(1Pt)^{2\xi (N3)}},$$
(22)
while Eq.(20) has the solution
$$r_{AH}=\frac{1}{\xi }\frac{|1Pt|}{|P|},$$
(23)
which represents the location of the apparent horizon of the solutions. When $`\xi =1`$, the apparent horizon represents a null surface in the $`(t,r)`$-plane, and when $`0\xi <1`$, the apparent horizon is spacelike, and when $`\xi =1`$, it is null, while when $`\xi >1`$, it is timelike. The spacetime is singular when $`t=1/P`$. This can be seen, for example, from the Kretschmann scalar, which now is given by
$$R^{\alpha \beta \gamma \delta }R_{\alpha \beta \gamma \delta }=6\xi ^2P^4\left\{N1+\xi ^2\left[N2+\underset{A=1}{\overset{N3}{}}(N2A)\right]\right\}(1Pt)^{4(1+\xi )}.$$
(24)
When $`P>0`$, it can be shown that the singularity always hides behind the apparent horizon, and when $`P<0`$, the singularity is naked. In the latter case, the solutions can be considered as representing cosmological models, while in the former the solutions as representing the formation of black holes due to the gravitational collapse of the perfect fluid. Substituting Eq.(23) into Eq.(21) we find that, as $`t+\mathrm{}`$, the mass of the black hole becomes infinitely large. To remend this shortage, one may follow to cut the spacetime along a timelike hypersurface, say, $`r=r_0(t)`$, and then join the part with $`r<r_0(t)`$ with an asymptotically flat N-dimensional spacetimes. From Eqs.(13) we can see that the fluid is comoving with the coordinates. Thus, the timelike hypersurface now can be chosen as $`r=r_0=Const.`$ Then, it can be seen that at the moment $`t_c=(P\xi r_01)/P`$, the whole ball collapses inside the apparent horizon, so the contribution of the collapsing ball to the total mass of such a formed black hole is given by
$$M_{BH}^F=m(t_c,r_0)=\frac{B_N}{2}\xi ^{\xi (N3)}r_0^{(N3)(1+\xi )}P^{\xi (N3)}.$$
(25)
From the above expression we can see that the mass is proportional to $`P`$, the parameter that characterizes the strength of the initial data of the collapsing ball. Thus, when the initial data is very weak ($`P0`$), the mass of the formed black hole is very small ($`M_{BH}0`$). In principle, by properly tuning the parameter $`P`$ we can make it as small as wanted.
It is interesting to note that this class of solutions admits a homothetic Killing vector,
$$\zeta _0=\frac{1Pt}{(1+\zeta )P},\zeta _1=\frac{r}{1+\zeta },\zeta _i=0,(i=2,3,\mathrm{},N1),$$
(26)
which satisfies the conformal Killing equation,
$$\zeta _{\mu ;\nu }+\zeta _{\nu ;\mu }=2g_{\mu \nu }.$$
(27)
Introducing two new coordinates via the relations,
$$\stackrel{~}{t}=\frac{(1Pt)^{\xi +1}}{(1+\xi )P},\stackrel{~}{r}=r^{1+\xi },$$
(28)
we find that the metric can be written in an explicit self-similar form,
$$ds^2=d\stackrel{~}{t}^2\left[(\xi +1)^{1/\xi }Px\right]^{\frac{2\xi }{\xi +1}}d\stackrel{~}{r}^2\left[(\xi +1)Px\right]^{\frac{2\xi }{\xi +1}}\stackrel{~}{r}^2d\mathrm{\Omega }_{N2}^2,$$
(29)
where $`x\stackrel{~}{t}/\stackrel{~}{r}`$ is the self-similar variable.
It is well-known that an irrotational “stiff” fluid ($`\alpha =1`$) in four-dimensional spacetimes is energetically equal to a massless scalar field . It can be shown that this also the case for N-dimensional spacetimes. In particular, for the above solutions with $`\alpha =1`$, the corresponding massless scalar field $`\varphi `$ is given by
$$\varphi =\pm \left[\frac{N1}{\kappa (N2)}\right]^{1/2}\mathrm{ln}\left(1Pt\right)+\varphi _0,$$
(30)
where $`\varphi _0`$ is a constant.
### B Type B solutions
The solutions in this case are given by Eq.(15). When $`f_1>0`$, the spacetime is close, and to have the metric be real the constant $`B`$ has to be imaginary. The spacetimes are singular when,
$$t|_{f_1>0}=\frac{1}{|\omega |}\text{arctan}\left(\frac{|B|}{A}\right)+2n\pi ,$$
(31)
where $`n`$ is an integer. When $`f_1<0`$, the spacetime is singular only when
$$t|_{f_1<0}=\frac{1}{\omega }\text{arctanh}\left(\frac{B}{A}\right).$$
(32)
Therefore, in the following we shall consider only the case where $`f_1<0`$. In this case, to have the energy density of the fluid be non-negative, we need to impose the condition $`B^2A^2`$. Then, the metric coefficient $`G(t)`$ can be written as
$$G=\left(B^2A^2\right)^\xi \mathrm{sinh}^{2\xi }[\omega (t_0ϵt)],$$
(33)
where $`ϵ=sign(B)`$, and $`t_0`$ is defined as
$$\mathrm{sinh}(\omega t_0)=\frac{A}{(B^2A^2)^{1/2}}.$$
Clearly, the conformal factor $`(B^2A^2)^\xi `$ does not play any significant role, without loss of generality, in the following we shall set it to be one. If we further introduce a new radial coordinate via the relation,
$$\overline{r}=h(t,r)𝑑r=\frac{1}{a}\mathrm{ln}\left|\frac{a+r}{ar}\right|,$$
(34)
where $`a(f_1)^{1/2}`$, the corresponding metric takes the form,
$$ds^2=\mathrm{sinh}^{2\xi }\left[\frac{2}{\xi }\left(t_0ϵt\right)\right]\left\{dt^2dr^2\frac{\mathrm{sinh}^2(2r)}{4}d\mathrm{\Omega }_{N2}^2\right\}.$$
(35)
Note that in writing the above expression, we had set, without loss of generality, $`a=1`$, and dropped the bars from $`r`$. From this metric, it can be shown that it is not self-similar. In the following we shall use this form of metric for the study of the Type B solutions. The corresponding mass function and Kretschmann scalar are given, respectively, by
$`m(r,t)`$ $`=`$ $`{\displaystyle \frac{B_N}{2^{N2}}}\mathrm{sinh}^{N1}(2r)\mathrm{sinh}^{\xi (N3)2}\left[2\xi ^1(t_0ϵt)\right],`$ (36)
$``$ $`=`$ $`{\displaystyle \frac{96}{\xi ^2}}\left\{(N1)+\xi ^2\left[N2+{\displaystyle \underset{A=1}{\overset{N3}{}}}(N2A)\right]\right\}\times \mathrm{sinh}^{4(\xi +1)}\left[2\xi ^1\left(t_0ϵt\right)\right],`$ (37)
while the apparent horizon is located at
$$r=r_{AH}\xi ^1(t_0ϵt).$$
(38)
From the above equations, we can see that the solutions are singular on the hypersurface $`t=ϵt_0`$. When $`ϵ=+1`$, the singularity is hidden behind the apparent horizon, and the solutions represent the formation of black holes from the gravitational collapse of the fluid. When $`ϵ=1`$, the singularity is naked, the solutions can be considered as representing cosmological models or white holes. As in the type A case, the mass of such formed black holes also diverges at the limit $`t+\mathrm{}`$. Thus, to have finite masses of black holes, we may also make a “surgery” to the spacetimes. Since the fluid is comoving with the coordinates, too. Without loss of generality, we may also choose the boudary as $`r=r_0`$. Then, it can be shown that the contribution of the collapsing ball to the total mass of black hole now is given by
$$M_{BH}^Fm_{AH}(\tau _{AH},r_0)=\frac{B_N}{2^{N2}}\mathrm{sinh}^{(N3)(\xi +1)}(2r_0).$$
(39)
From the above expression we can see that for any given non-zero $`r_0,M_{BH}^F`$ is always finite and non-zero. Thus, in the present case black holes start to form with a mass gap.
Finally we note that, similar to the last case, the solution with $`\alpha =1`$ also corresponds to a massless scalar field with the scalar field being given by
$$\varphi =\pm \left[\frac{N1}{\kappa (N2)}\right]^{1/2}\mathrm{ln}\left\{\mathrm{tanh}[(N2)(t_0ϵt)]\right\}+\varphi _0,(\alpha =1).$$
(40)
## IV Concluding Remarks
The general form of metric for N-dimensional spherically symmetric and conformally flat spacetimes was found. As an application of it, all the Friedmann-Robertson-Walker-like solutions for a perfect fluid with an equation of state $`p=\alpha \rho `$ were given. These solutions were then used to model the gravitational collapse of a compact ball. It was shown that when the collapse has continuous self-similarity, the formation of black holes always starts with zero mass, and when the collapse has no such a symmetry, the formation of black holes always starts with a finite non-zero mass. This is the same as that obtained in the study of the problem in four-dimensional spacetimes . Thus, they provide further evidences to support the speculation that the formation of black holes always starts with zero-mass for the collapse with self-similarities.
###### Acknowledgements.
We would like to express our gratitude to S.K. Chatterjee for valuable discussions. The financial assistance from CAPES (JFVR), CNPq (AW) and FAPERJ (AW) is gratefully acknowledged.
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# (hep-ph/0007nnn, 31. July 2000) Implications of a Low 𝐬𝐢𝐧{𝟐𝜷}: A Strategy for Exploring New Flavor Physics
## 1 Introduction
One of the highlights of the ongoing ICHEP2000 Conference in Osaka has been the presentation of first results on CP violation in $`B_d`$$`\overline{B}_d`$ mixing by the BaBar and Belle Collaborations . The reported values obtained from the time-dependent CP asymmetry in $`BJ/\psi K_S`$ decays,
$$\mathrm{sin}2\beta _{\psi K}=\{\begin{array}{cc}0.12\pm 0.37\pm 0.09;\hfill & \text{BaBar,}\hfill \\ 0.45_{0.440.09}^{+0.43+0.07};\hfill & \text{Belle,}\hfill \end{array}$$
(1)
are smaller than the previous best measurement $`\mathrm{sin}2\beta _{\psi K}=0.79_{0.44}^{+0.41}`$ by the CDF Collaboration . They are also smaller than the value $`\mathrm{sin}2\beta =0.75\pm 0.06`$ obtained from a recent global analysis of the unitarity triangle . Although there is at present no statistically significant discrepancy, it is interesting to explore the implications of a measurement of $`\mathrm{sin}2\beta _{\psi K}`$ that would be inconsistent with the results of the global analysis. Under the assumption that there is no CP-violating New Physics in $`bc\overline{c}s`$ transitions (which is supported by the strong experimental bound on direct CP violation in $`B^\pm J/\psi K^\pm `$ decays reported by the CLEO Collaboration ) this would imply the existence of New Physics in $`B_d`$$`\overline{B}_d`$ mixing. (New Physics in $`K`$$`\overline{K}`$ mixing could not account for such a discrepancy, because of the minor impact of $`|ϵ_K|`$ on the global analysis.) The measured phase $`2\beta _{\psi K}`$ would then be the $`B_d`$$`\overline{B}_d`$ mixing phase $`2\varphi _d`$, which would differ from the CKM phase $`2\beta `$ because of New Physics. In such an event, it is likely that New Physics would also play a role in $`B_s`$$`\overline{B}_s`$ and $`K`$-$`\overline{K}`$ mixing.
The purpose of this Letter is to point out a strategy which provides a systematic exploration of the new flavor physics in this case. This strategy is different from the conventional route pursued at the $`B`$ factories, in which the main focus is on measurements that are sensitive to $`B`$$`\overline{B}`$ mixing, mainly because CP violation in the interference of mixing and decay can sometimes be interpreted without encountering large hadronic uncertainties. If there is New Physics in mixing, then the standard triangle obtained by combining information on $`|V_{ub}|`$ from semileptonic $`B`$ decay, $`\mathrm{\Delta }m_{d,s}`$ from $`B_{d,s}`$$`\overline{B}_{d,s}`$ mixing, and $`|ϵ_K|`$ from $`K`$$`\overline{K}`$ mixing does not agree with the true CKM triangle, and forcing it to close (as is done in the standard analysis) gives wrong results for the angles $`\gamma =\text{arg}(V_{ub}^{})`$ and $`\beta =\text{arg}(V_{td})`$. (We use the standard phase conventions; otherwise $`\gamma =\text{arg}[(V_{ub}^{}V_{ud})/(V_{cb}^{}V_{cd})]`$ and $`\beta =\text{arg}[(V_{tb}^{}V_{td})/(V_{cb}^{}V_{cd})]`$.)
In the absence of a reliable way to measure the magnitude and phase of $`V_{td}`$ in $`B`$ decays, it is important to base studies of the CKM matrix on a reference triangle obtained exclusively from measurements independent of particle–antiparticle mixing . In the $`B`$ system, this triangle is constructed from the measurement of the magnitude and phase of $`V_{ub}`$. (The use of $`\gamma =\text{arg}(V_{ub}^{})`$ replaces the use of $`|V_{td}|`$, determined from $`B`$$`\overline{B}`$ mixing, in the standard analysis.) Separate comparisons of particle–antiparticle mixing measurements in the $`B_d`$, $`B_s`$ and $`K`$ systems with information obtained from the reference triangle will allow extraction of the magnitude and phase of New Physics contributions to the mixing amplitudes. At a later stage, comparison of different reference triangle constructions can provide information about potential New Physics effects in the decay amplitudes, not related to mixing.
We stress that the reference triangle approach should be pursued regardless of whether or not the $`\mathrm{sin}2\beta _{\psi K}`$ measurements are consistent with the global analysis of the unitarity triangle. Agreement within errors could be accidental and would not exclude the possibility of large New Physics contributions in $`B_d`$$`\overline{B}_d`$ mixing. Proposals similar in spirit to ours have been discussed in the past. However, their feasibility is limited by their reliance on methods for extracting $`\gamma `$ that are extremely difficult, and are plagued by multiple discrete ambiguities . In the past two years, however, several strategies have been proposed that will allow a determination of $`\gamma `$ in the near future, without discrete ambiguities and with controlled theoretical uncertainties.
Our analysis is based on the following standard assumptions, which hold true for a vast class of extensions of the Standard Model (for a discussion, see e.g. Ref. ):
i) The determination of the CKM elements $`|V_{us}|`$, $`|V_{cb}|`$ and $`|V_{ub}|`$ from semileptonic decays is not affected by New Physics.
ii) The 3-generation CKM matrix is unitary.
iii) There are no (or negligibly small) New Physics effects in decays which in the Standard Model are dominated by tree topologies.
An experimental test for non-standard contributions in the semileptonic $`bul\nu `$ transition could be performed by comparing the values of $`|V_{ub}|`$ extracted from the exclusive $`B\pi l\nu `$ and $`B\rho l\nu `$ decays, and the inclusive $`BX_ul\nu `$ decays. (An analogous test for $`bcl\nu `$ decays has been proposed in Ref. .) Tests of the unitarity of the CKM matrix will be discussed in Section 4.
## 2 The reference triangle
Disregarding all information obtained from mixing measurements, not much is known about the Wolfenstein parameters $`(\overline{\rho },\overline{\eta })`$ determining the unitarity triangle. The magnitude of $`|V_{ub}|`$ measured in semileptonic $`B`$ decay fixes $`R_b=|(V_{ub}^{}V_{ud})/(V_{cb}^{}V_{cd})|=\sqrt{\overline{\rho }^2+\overline{\eta }^2}`$, corresponding to a circle centered at the origin in the $`(\overline{\rho },\overline{\eta })`$ plane. At present $`|V_{ub}|`$ is known with a precision of about 20%. A reduction of the uncertainty to the 10% level appears realistic within a few years. The phase $`\gamma `$ defining the orientation of the triangle ($`\mathrm{sin}\gamma =\overline{\eta }/R_b`$) is currently unknown.
### 2.1 Reference triangles from $`B`$ decays
In the very near future, ratios of CP-averaged $`B\pi K`$ and $`B\pi \pi `$ branching ratios can be used to extract $`\mathrm{cos}\gamma `$ using many different strategies. A method based on flavor symmetries, using little theory input, has been described in . It makes use of two experimentally determined rate ratios ($`R_{}`$ and $`\overline{\epsilon }_{3/2}`$) and the theoretical prediction that the relevant strong-interaction phase cannot be very large. Alternatively, it has been argued recently that in the heavy-quark limit most two-body hadronic $`B`$ decays admit a model-independent theoretical description based on a QCD factorization formula . Predictions for the $`B\pi K,\pi \pi `$ decay-rate ratios as a function of $`\mathrm{cos}\gamma `$ have been obtained (including the leading power corrections in $`1/m_b`$) . The combination of several independent determinations of $`\mathrm{cos}\gamma `$ from these modes will fix $`\gamma `$, up to a sign ambiguity $`\gamma \gamma `$, with reasonable precision. (We define all weak phases to lie between $`180^{}`$ and $`180^{}`$.) We believe that an uncertainty of $`\mathrm{\Delta }\gamma =25^{}`$ will be attainable in the near future.
Once the $`B^\pm (\pi K)^\pm `$ and $`B^\pm \pi ^\pm \pi ^0`$ decay rates are known with higher precision, $`\gamma `$ can be determined with minimal theory input (up to discrete ambiguities) using the method of Ref. . Here, in addition to CP-averaged decay rates, information about some direct CP asymmetries is added. Ultimately, this will reduce the theoretical uncertainty to a level of $`10^{}`$ or less. When supplemented with theoretical information on the strong-interaction phase this method can be used to completely eliminate the discrete ambiguities (for a detailed discussion, see Ref. ).
The QCD factorization approach can also be used to calculate the penguin-to-tree ratio in $`B\pi ^+\pi ^{}`$ decays, thereby turning a measurement of the mixing-induced CP violation into a determination of $`\gamma `$ without the need for an (impractical) isospin analysis . In Figure 1 we show the coefficient $`S`$ of $`\mathrm{sin}(\mathrm{\Delta }m_dt)`$ in the time-dependent CP asymmetry as a function of $`\gamma `$, assuming $`\mathrm{sin}2\varphi _d=0.3`$ for the $`B`$$`\overline{B}`$ mixing phase measured in $`BJ/\psi K_S`$. A measurement of $`\mathrm{sin}2\varphi _d`$ determines the phase $`2\varphi _d`$ up to a two-fold discrete ambiguity, $`2\varphi _d^{(1)}+2\varphi _d^{(2)}=\pi `$ mod $`2\pi `$. The width of the bands reflects the theoretical uncertainty. The eight-fold ambiguity can be reduced to a four-fold one, in principle, by measuring the direct CP asymmetry in this decay. Alternatively, using a Dalitz-plot analysis of the $`B\rho \pi `$ decay amplitudes one could determine $`\mathrm{sin}(2\varphi _d+2\gamma )`$ with small hadronic uncertainties .
On a much longer time-scale, it will be possible to obtain information on $`\gamma `$ (again up to discrete ambiguities) using only decays mediated by tree topologies in the Standard Model. Examples are the determination of $`\mathrm{sin}(2\varphi _d+\gamma )`$ from $`BD^{()\pm }\pi ^{}`$ decays , and the extraction of $`\gamma `$ from $`BDK`$ decays . This last method, in particular, will require very large data samples. Other methods make use of $`B_s`$-meson decays accessible at future $`B`$ factories at hadron colliders . In Figure 2 we show as an example the information obtainable from a determination of $`\mathrm{sin}(2\varphi _d+\gamma )`$. Combining this with the information derived from a measurement of the $`B\pi ^+\pi ^{}`$ CP asymmetry (see Figure 1), a unique pair of solutions $`(\gamma ,2\varphi _d^{(1)})`$ and $`(\gamma ,2\varphi _d^{(2)})`$ is obtained. We stress that combining any of the measurements sensitive to $`B_d`$$`\overline{B}_d`$ mixing described above with a determination of $`\gamma `$ (including its sign) from $`B\pi K`$ would remove the discrete ambiguity in the $`B_d`$ mixing phase $`2\varphi _d`$.
Up to now we have assumed that the various determinations of the reference triangle from $`B`$ decays are consistent with each other. This assumption will have to be tested as the data become increasingly precise. In Section 4 we will discuss how differences between these constructions would provide information about New Physics in $`B`$ decays rather than $`B_d`$$`\overline{B}_d`$ mixing.
### 2.2 Reference triangle from $`K`$ decays
An independent reference triangle can be constructed from measurements of very rare kaon decays. The branching ratios for the decays $`K^+\pi ^+\nu \overline{\nu }`$ and $`K_L\pi ^0\nu \overline{\nu }`$ measure $`|V_{ts}^{}V_{td}|`$ and $`|\text{Im}(V_{ts}^{}V_{td})|`$, respectively, and thereby determine $`R_t=\sqrt{(1\overline{\rho })^2+\overline{\eta }^2}`$ and $`|\eta |`$ independently of $`K`$$`\overline{K}`$ mixing . This provides a reference triangle up to a four-fold discrete ambiguity. Dedicated experiments would be necessary to measure $`R_t`$ and $`|\eta |`$ with useful precision. In Section 4 we will discuss what can be learned from the comparison of the kaon reference triangle with the $`B`$-meson triangle(s).
## 3 Exploring New Physics
Once the reference triangle is known, we can use it to explore New Physics contributions to $`B`$$`\overline{B}`$ and $`K`$$`\overline{K}`$ mixing. Measurement of $`R_b`$ and $`\gamma `$ fix the coordinates $`\overline{\rho }`$ and $`\overline{\eta }`$, which in turn determine the other side $`R_t`$ and the true angle $`\beta `$ of the reference triangle via
$`R_t`$ $`=`$ $`\sqrt{(1\overline{\rho })^2+\overline{\eta }^2},`$
$`\mathrm{sin}\beta `$ $`=`$ $`{\displaystyle \frac{\overline{\eta }}{R_t}},\mathrm{cos}\beta ={\displaystyle \frac{1\overline{\rho }}{R_t}}.`$ (2)
In Figure 3, we illustrate what the situation may look like both in the near and long-term future. We assume that in the near future $`\gamma `$ will be known only up to a sign ambiguity (from measurements of CP-averaged decay rates), which will be resolved after several more years of running at the $`B`$ factories (when certain CP asymmetries in rare decays will have been measured).
### 3.1 $`B_d`$$`\overline{B}_d`$ mixing
We now discuss how one can systematically study New Physics effects in the $`B_d`$$`\overline{B}_d`$ mixing amplitude $`M_{12}`$ by confronting measurements of the mass difference $`\mathrm{\Delta }m_d=2|M_{12}|`$ (or $`x_d=\mathrm{\Delta }m_d\tau _B`$) and of the mixing phase $`\mathrm{sin}2\varphi _d`$ with the reference triangle. Our approach is very similar to the one discussed in Ref. . Using (3) we construct the complex quantity $`R_t^2e^{2i\beta }`$ with the true CKM phase $`\beta `$. Up to a constant $`C_B`$, this quantity determines the Standard Model contribution to the mixing amplitude: $`M_{12}^{\mathrm{SM}}=C_BR_t^2e^{2i\beta }`$, where
$`C_B`$ $`=`$ $`{\displaystyle \frac{G_F^2}{12\pi ^2}}\eta _Bm_Bm_W^2S_0(x_t)B_Bf_B^2`$ (3)
$``$ $`0.24\text{ps}^1\times {\displaystyle \frac{B_Bf_B^2}{(0.2\text{GeV})^2}}.`$
The main uncertainty in this result comes from the hadronic matrix element parameterized by the product $`B_Bf_B^2`$. It is therefore convenient to focus on the ratio $`M_{12}/C_B`$, which in the Standard Model is given only in terms of CKM parameters: $`M_{12}^{\mathrm{SM}}/C_B=[(1\overline{\rho })^2\overline{\eta }^2)]2i\overline{\eta }(1\overline{\rho })`$. (If New Physics does not induce operators with non-standard Dirac structure, the ratio $`M_{12}/C_B`$ remains free of hadronic uncertainties.) The experimental value of the mixing amplitude is given by $`M_{12}^{\mathrm{exp}}/C_B=(\mathrm{\Delta }m_d/2C_B)e^{2i\varphi _d}`$. If the mixing phase is determined from the $`\mathrm{sin}2\varphi _d`$ measurement in $`BJ/\psi K_S`$ decays alone, then $`e^{2i\varphi _d}`$ has a two-fold discrete ambiguity. In the previous section we have discussed how this ambiguity may eventually be resolved by using data on CP violation in $`B`$ decays. The difference $`M_{12}^{\mathrm{NP}}=M_{12}^{\mathrm{exp}}M_{12}^{\mathrm{SM}}`$ is the New Physics contribution to the mixing amplitude.
In Figure 4 we illustrate this analysis using present-day values of the input parameters from Ref. ($`|V_{ub}/V_{cb}|=0.085\pm 0.018`$, $`\sqrt{B_B}f_B=(0.21\pm 0.04)\text{GeV}`$) and the average of the new BaBar and Belle results, $`\mathrm{sin}2\varphi _d=0.26\pm 0.29`$. The upper plot shows the allowed regions for the Standard Model contribution to the mixing amplitude as well as for its experimental value, taking into account the two-fold ambiguity in the mixing angle $`2\varphi _d`$. The difference between any point in the Standard Model regions with any point in the data regions defines an allowed vector in the complex $`M_{12}^{\mathrm{NP}}`$ plane. In the lower plot we show the resulting allowed regions for $`M_{12}^{\mathrm{NP}}`$. The origin in this plot corresponds to the Standard Model. We also show the results in the absence of any information on $`\gamma `$. An important message from this plot is that a potentially large New Physics contribution (of order the Standard Model contribution) to the mixing amplitude is allowed by the data in large portions of parameter space. In order to find out whether or not there is indeed such a large contribution it will be necessary to determine $`\gamma `$ and resolve the discrete ambiguity in $`2\varphi _d`$.
Figure 5 illustrates what the situation may look like several years from now. By then the uncertainties in the input parameters will most likely have been reduced significantly, and the mixing angle will have been measured with good precision. For the purpose of illustration we take $`|V_{ub}/V_{cb}|=0.085\pm 0.009`$, $`\sqrt{B_B}f_B=(0.21\pm 0.02)\text{GeV}`$, and $`\mathrm{sin}2\varphi _d=0.26\pm 0.10`$. Also, the phase $`\gamma `$ will have been measured more accurately, perhaps with $`\mathrm{\Delta }\gamma =10^{}`$. This would lead to the picture shown in Figure 5(a) and to the allowed regions for New Physics shown in (b). Most importantly, as we have explained, in the long term the discrete ambiguities in both $`\gamma `$ and $`2\varphi _d`$ will be removed, so that it would be possible to identify which one of the four regions in (b) is realized in nature. At this point, we will have achieved a precise determination of the New Physics contribution to $`B_d`$$`\overline{B}_d`$ mixing.
### 3.2 $`B_s`$$`\overline{B}_s`$ mixing
In the presence of New Physics in $`B_d`$$`\overline{B}_d`$ mixing, it is very likely that also the $`B_s`$$`\overline{B}_s`$ mixing amplitude is different from its value in the Standard Model. Therefore, measurements sensitive to this mixing amplitude should not be combined with measurements in the $`B_d`$ system. Rather, one should probe for New Physics in the $`B_s`$ system in an independent way. In the Standard Model the assumption of unitarity of the CKM matrix alone fixes the magnitude of $`|V_{tb}^{}V_{ts}|`$ (and hence the Standard Model contribution to $`\mathrm{\Delta }m_s`$), and in addition implies that the $`B_s`$ mixing phase is very small, $`\varphi _s^{\mathrm{SM}}=O(\lambda ^2)`$ (with $`\lambda 0.22`$ the Wolfenstein parameter). Even at the approved hadron collider experiments BTeV and LHCb it will not be possible to measure this small Standard Model phase. It follows that the complex amplitude $`M_{12}^{\mathrm{SM}}(B_s)`$ is determined by unitarity and is very nearly real. (In that sense the “$`B_s`$ reference triangle” is almost degenerate to a line.) Measuring the true values of the mass difference $`\mathrm{\Delta }m_s`$ and of the mixing phase $`2\varphi _s`$ (e.g., from the time-dependent CP asymmetry in $`B_sJ/\psi \varphi `$ decays), one can then construct the mixing amplitude from $`M_{12}(B_s)=(\mathrm{\Delta }m_s/2)e^{2i\varphi _s}`$. The difference $`M_{12}^{\mathrm{NP}}(B_s)=M_{12}(B_s)M_{12}^{\mathrm{SM}}(B_s)`$ determines directly the New Physics contribution to $`B_s`$$`\overline{B}_s`$ mixing.
### 3.3 $`K`$$`\overline{K}`$ mixing
The mass difference $`\mathrm{\Delta }m_K`$ between the neutral kaon mass eigenstates is dominated by long-distance physics and does not admit a clean theoretical interpretation. Therefore, constraints on the CKM matrix from $`K`$$`\overline{K}`$ mixing are derived only from the CP-violating quantity $`|ϵ_K|`$, which (to a very good approximation) is given by $`|ϵ_K||\text{Im}[M_{12}(K)]|/(\sqrt{2}\mathrm{\Delta }m_K)`$. Consequently, one can only derive information on the New Physics contribution to the imaginary part of the mixing amplitude in the kaon system.
The Standard Model contribution to $`|ϵ_K|`$ is the product of a function of the Wolfenstein parameters $`\overline{\rho }`$ and $`\overline{\eta }`$ with a hadronic matrix element parameterized by the quantity $`B_K`$ . Once we have determined $`\overline{\rho }`$ and $`\overline{\eta }`$ (as a function of $`\gamma `$ and $`R_b`$) from the reference triangle, we can compute the Standard Model contribution. Subtracting it from the measured value of $`|ϵ_K|`$ gives the New Physics contribution, $`|ϵ_K^{\mathrm{NP}}|=|\text{Im}[M_{12}^{\mathrm{NP}}(K)]|/(\sqrt{2}\mathrm{\Delta }m_K)`$.
In Figure 6 we show the New Physics contribution to $`|ϵ_K|`$ as a function of $`\gamma `$, for both present-day and more long-term uncertainties on $`B_K`$ and $`|V_{ub}/V_{cb}|`$. It is evident that a measurement of $`\gamma `$ is the key ingredient needed to answer the question of whether or not there is New Physics in $`K`$$`\overline{K}`$ mixing. Once $`\gamma `$ is known, $`|\text{Im}[M_{12}^{\mathrm{NP}}(K)]|`$ can be extracted with good precision.
## 4 New Physics in decays
In order for a reference triangle construction to give the true value of $`\gamma `$, and thus be useful for extracting potential New Physics contributions to the mixing amplitudes, it must not be contaminated by additional New Physics effects in the associated decay amplitudes. Up to now we have assumed consistency between the different constructions, which would imply that to a good approximation New Physics enters only in the mixing amplitudes. This is indeed the case in many extensions of the Standard Model, particularly if the scale of new flavor interactions is at a TeV or beyond. Of course, it would be extremely exciting if the different reference triangle constructions were not consistent with one another, implying that there is New Physics in $`B`$ decay amplitudes, or that the 3-generation CKM matrix is not unitary.
In Section 2 we have discussed various reference triangle constructions in roughly the chronological order in which it will be possible to carry them out. Interestingly, this order also corresponds to a progression from reliance on decays which in the Standard Model are penguin-dominated (and therefore more susceptible to New Physics) to those which are tree-dominated (and therefore less susceptible to New Physics), and finally to decays that are only based on tree topologies.
We briefly discuss tests for New Physics contributions to the penguin-dominated decays based on $`bs\overline{q}q`$ transitions in the Standard Model. These tests can be carried out at various stages of data collection. The measurement of several CP-averaged $`B\pi K`$ and $`B\pi \pi `$ decay rates itself provides for a series of internal consistency checks. For instance, there are upper and lower bounds on certain rate ratios, which are based on flavor symmetries and rely on minimal theoretical input. Violation of these bounds would be a signal for new isospin-violating New Physics contributions . In the longer term, as measurements of CP asymmetries for rare decays become available, additional tests will become possible. For example, one can then check whether the time-dependent CP asymmetries in $`BJ/\mathrm{\Psi }K_S`$ and $`B\varphi K_S`$ decays are in agreement. A discrepancy would imply new CP-violating contributions to the $`B\varphi K_S`$ decay amplitude, which to a good approximation is a pure $`bs\overline{s}s`$ penguin amplitude in the Standard Model . There are also upper bounds on the direct CP asymmetries for $`B^\pm \pi ^\pm K^0`$ and $`B^\pm \varphi K^\pm `$ in the Standard Model, which could turn out to be violated . Finally, a large direct CP asymmetry in $`B^\pm X_s\gamma `$ decays would imply significant New Physics contributions to $`bs`$ penguin transitions .
If all penguin-dominated determinations of $`\gamma `$ are consistent, then it is unlikely that there is New Physics in decays, and the program we have outlined above for extracting New Physics in mixing can already be carried out with confidence. If these determinations are not all consistent, however, the tree-dominated or pure tree reference triangle constructions will be required in order to reliably extract New Physics contributions to mixing. (At the same time, a clean determination of $`\gamma `$ would be be required in order to obtain a detailed picture of the New Physics in the penguin-dominated decays .) The determination of $`\gamma `$ from the mixing-induced CP asymmetry in $`B\pi ^+\pi ^{}`$ is less susceptible to New Physics effects, since these would have to compete with a Cabibbo-enhanced tree-level amplitude. Finally, the extractions of $`\gamma `$ from the pure tree processes $`BD^{()}\pi `$ and $`BDK`$ can only be affected by New Physics in rather exotic scenarios. Thus, checking for consistency between these two measurements and the $`B\pi ^+\pi ^{}`$ measurements essentially provides a test for New Physics effects in $`B\pi \pi `$.
If dedicated experiments to study the very rare $`K\pi \nu \overline{\nu }`$ decay modes will be performed, they will provide direct measurements of the magnitude and phase of $`V_{td}`$ independent of mixing. The comparison of the kaon reference triangle with the triangles obtained from $`B`$ decays primarily allows us to probe for New Physics in these rare kaon decays, which in the Standard Model are mediated by box and electroweak penguin diagrams. In addition, if the kaon and $`B`$-mesons triangles were to agree with one another, this would be a direct test of the assumption of 3-generation CKM unitarity, as it would check the relation $`V_{ub}^{}V_{ud}+V_{cb}^{}V_{cd}+V_{tb}^{}V_{td}=0`$ independent of any mixing measurements. Another test of CKM unitarity would be the direct measurement of the element $`|V_{tb}|`$ (and perhaps $`|V_{ts}|`$) in top decay.
## 5 Conclusions
Motivated by today’s announcement of the first $`\mathrm{sin}2\beta _{\psi K}`$ measurements from the dedicated $`B`$-factory experiments BaBar and Belle, we have reconsidered strategies for exploring New Physics in $`B`$$`\overline{B}`$ and $`K`$$`\overline{K}`$ mixing in a model-independent way. The low central values found by these experiments raise the possibility that there is New Physics in $`B_d`$$`\overline{B}_d`$ mixing. Therefore it becomes crucial to base studies of flavor physics on comparisons with a reference unitarity triangle whose construction is independent of mixing measurements. Ultimately, such a strategy is preferable whether or not the measured value of $`\mathrm{sin}2\beta _{\psi K}`$ agrees with the prediction from the global analysis of the unitarity triangle.
We have described in detail a program that in a few years could cleanly determine the New Physics contributions (in magnitude and phase) to the $`B_d`$$`\overline{B}_d`$ and $`K`$$`\overline{K}`$ mixing amplitudes. (A similar, more straightforward analysis for $`B_s`$$`\overline{B}_s`$ mixing can be performed at the BTeV and LHCb experiments.) Similar strategies have been proposed previously by several authors. Here we have stressed the relevance of the progress recently made in devising strategies for near-term measurements of the weak phase $`\gamma `$ based on charmless hadronic $`B`$ decays. Knowing $`\gamma `$ with reasonable accuracy, and without discrete ambiguities, is the key element that makes the program outlined in this Letter feasible and very powerful. We have also pointed out that the comparison of different constructions of the reference triangle provides several opportunities for probing New Physics in decay amplitudes, not related to mixing. Nothing would be more exciting than to follow the unfolding of New Physics at the $`B`$ factories in the next few years ahead of us.
Acknowledgements: We wish to thank the SLAC Theory Group for its hospitality while this work was carried out. A.K. is supported by the Department of Energy under Grant No. DE-FG02-84ER40153. M.N. is supported in part by the National Science Foundation.
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# RADEMACHER CHAOS IN SYMMETRIC SPACES, II
## 1 Introduction
This paper is a continuation of where we started the study of Rademacher chaos in functional symmetric spaces (s.s.) on the segment $`[0,1]`$. Let us first recall some definitions and notations from .
As usual,
$$r_k(t)=\mathrm{sign}\mathrm{sin}2^{k1}\pi t(k=1,2,\mathrm{})$$
denotes the system of Rademacher functions on $`I:=[0,1]`$. The set of all real-valued functions $`x(t)`$ that can be represented in the form
$$x(t)=\underset{1i<j<\mathrm{}}{}a_{i,j}r_i(t)r_j(t)(t[0,1])$$
is called a chaos of order $`2`$ with respect to the system $`\{r_k(t)\}`$ (Rademacher chaos of order 2 ). The same name is used, with no ambiguity, for the orthonormal system of functions $`\{r_ir_j\}_{1i<j<\mathrm{}}`$. In the sequel, as in , $`H`$ denotes the closure of $`L_{\mathrm{}}`$ in the Orlicz space $`L_M`$ where $`M(t)=e^t1`$.
In , we proved the following.
Theorem A. Let $`X`$ be a symmetric space. Then the following statements are equivalent:
$`1)`$ The system $`\{r_ir_j\}_{1i<j<\mathrm{}}`$ in $`X`$ is equivalent to the canonical basis in $`l_2`$;
$`2)`$ A continuous imbedding $`HX`$ takes place.
In this paper, we shall consider questions related to the unconditionality of Rademacher chaos. Our main result is:
The statements $`1)`$ and $`2)`$ in Theorem A are equivalent to the next one:
$`3)`$ The system $`\{r_ir_j\}_{1i<j<\mathrm{}}`$ is an unconditional basic sequence in s.s. $`X`$.
Let us recall the meaning of the central notions above.
Definition. A sequence $`\{x_n\}_{n=1}^{\mathrm{}}`$ of elements in Banach space $`X`$ is called a basic sequence if it is a basis in its closed linear span $`[x_n]_{n=1}^{\mathrm{}}`$.
As is well-known (see for example \[11, p.2\] ), the latter is equivalent to the following two conditions:
1) $`x_n\mathrm{\hspace{0.17em}0}`$ for all $`n`$;
2) The family of projectors
$$P_m\left(\underset{i=1}{\overset{\mathrm{}}{}}a_ix_i\right)=\underset{i=1}{\overset{m}{}}a_ix_i(m=1,2,\mathrm{}),$$
defined on $`[x_n]_{n=1}^{\mathrm{}}`$, is uniformly bounded. That is, a constant $`K>0`$ exists such that for all $`m,n`$, $`m<n`$, and $`a_i`$, the following inequality holds:
$$\underset{i=1}{\overset{m}{}}a_ix_iK\underset{i=1}{\overset{n}{}}a_ix_i.$$
$`(1)`$
One of the most important properties of a basic sequence is its unconditionality.
Definition. A basic sequence $`\{x_n\}_{n=1}^{\mathrm{}}`$ in a Banach space $`X`$ is said to be unconditional if, for any rearrangement $`\pi `$ of $``$, the sequence $`\{x_{\pi (n)}\}_{n=1}^{\mathrm{}}`$ is also a basic sequence in $`X`$.
This is equivalent, in particular, to the uniform boundedness of the family of operators
$$M_\theta \left(\underset{i=1}{\overset{\mathrm{}}{}}a_ix_i\right)=\underset{i=1}{\overset{\mathrm{}}{}}\theta _ia_ix_i(\theta _i=\pm 1)$$
which are defined on $`[x_n]_{n=1}^{\mathrm{}}`$ \[11, p.18\], and the last means that there is a constant $`K_0`$ such that for each $`n`$ and any couple of sequences of signs $`\{\theta _i\}`$ and real numbers $`\{a_i\}`$,
$$\underset{i=1}{\overset{n}{}}\theta _ia_ix_iK_0\underset{i=1}{\overset{n}{}}a_ix_i.$$
$`(2)`$
Finally, note that for sequences of real numbers $`(a_{i,j})_{1i<j<\mathrm{}}`$ we use the common notation
$$(a_{i,j})_2:=\left(\underset{1i<j<\mathrm{}}{}a_{i,j}^2\right)^{1/2}.$$
## 2 Rademacher chaos as a basic sequence
The system $`\{r_k\}_{k=1}^{\mathrm{}}`$ and Rademacher chaos $`\{r_ir_j\}_{1i<j<\mathrm{}}`$, both are special subsystems of Walsh system $`\{w_n\}_{n=0}^{\mathrm{}}`$. If the latter is considered with Paley indexing \[6, p.158\], then $`w_0=r_1,w_{2^k}=r_{k+2},k=0,1,\mathrm{}`$ We shall enumerate Rademacher chaos in correspondence to this indexing, namely,
$$\begin{array}{c}\phi _1=r_1r_2=r_2,\phi _2=r_1r_3=r_3,\phi _3=r_2r_3,\phi _4=r_1r_4=r_4,\mathrm{},\\ \\ \phi _{k(k1)/2+1}=r_1r_{k+1}=r_{k+1},\mathrm{},\phi _{k(k+1)/2}=r_kr_{k+1},\mathrm{}\end{array}$$
$`(3)`$
Before formulating our first theorem let us recall the definition of a fundamental notion in the interpolation theory of operators (for more details, see ).
Definition. A Banach space $`X`$ is said to be an interpolation space with respect to the Banach couple $`(X_0,X_1)`$ if $`X_0X_1XX_0+X_1`$ and, in addition, if the boundedness of a linear operator $`T`$ in both $`X_0`$ and $`X_1`$ implies its boundedness in $`X`$ as well.
Theorem 1. The Rademacher chaos $`\{r_ir_j\}_{1i<j<\mathrm{}}`$, ordered according to rule $`(3)`$, is a basic sequence in every interpolation with respect to the couple $`(L_1,L_{\mathrm{}})`$ s.s. $`X`$ on $`[0,1]`$ .
Proof. In the sequel we shall use the following property of Walsh systems (see \[4, p.45\]). Introduce the Fourier-Walsh partial sum operators
$$S_px(t):=\underset{i=0}{\overset{p}{}}_0^1x(s)w_i(s)𝑑sw_i(t)(p=0,1,\mathrm{})$$
and denote $`\sigma _k:=S_{2^k}`$. Set
$$\mathrm{\Delta }_j^{(k)}:=((j1)2^k,j2^k)(k=0,1,\mathrm{};j=1,2,\mathrm{},2^k).$$
Then
$$\sigma _kx(t)=\mathrm{\hspace{0.33em}2}^k_{\mathrm{\Delta }_j^{(k)}}x(u)𝑑u$$
for each $`t\mathrm{\Delta }_j^{(k)}`$. In other words, the operator $`\sigma _k`$ coincides with the averaging operator over the system of dyadic intervals $`\{\mathrm{\Delta }_j^{(k)}\}_{j=1}^{2^k}`$. It is easy to see that such an operator is bounded in $`L_1`$ and $`L_{\mathrm{}}`$, and more precisely, its norm is equal to $`1`$ in both spaces. Thus $`\sigma _k`$ is bounded in $`X`$ as well. Therefore, there exists a constant $`B=B(X)>0`$ such that
$$\sigma _kx_XBx_X$$
$`(4)`$
for all $`xX`$.
For given natural numbers $`m<n`$, and real numbers $`a_1,\mathrm{},a_n`$, set
$$y(t)=\underset{i=1}{\overset{n}{}}a_i\phi _i(t),z(t)=\underset{i=1}{\overset{m}{}}a_i\phi _i(t).$$
In the simplest case, when $`\phi _m=r_{k+2}`$ for some $`k=0,1,\mathrm{}`$, the orthonormality of Walsh system yields $`z=\sigma _ky`$. Therefore, taking into account (4), we get $`zBy`$. Thus, for $`x_i=\phi _i`$, inequality (1) holds with a constant $`K=B`$.
Consider now the general case: for some $`0kl`$, $`2p<k+2`$, $`2q<l+2`$,
$$z(t)=\sigma _ky(t)+\underset{j=2}{\overset{p}{}}b_jr_j(t)r_{k+2}(t)$$
and
$$y(t)=\sigma _ly(t)+\underset{j=2}{\overset{q}{}}c_jr_j(t)r_{l+2}(t).$$
There are two possibilities.
Case 1. $`k=l,pq`$.
Set
$$f(t)=z(t)\sigma _ky(t)=\underset{j=2}{\overset{p}{}}b_jr_j(t)r_{k+2}(t),$$
$$g(t)=y(t)z(t)=\underset{j=p+1}{\overset{q}{}}c_jr_j(t)r_{k+2}(t).$$
It follows from the definition of Rademacher functions that the absolute values of $`u(t)=f(t)+g(t)`$ and $`v(t)=f(t)g(t)`$ are equimeasurable. Then the symmetry of $`X`$ yields $`u_X=v_X`$. Since $`f=(u+v)/2`$ we have
$$f_Xu_X.$$
$`(5)`$
Taking into account that
$$u(t)=\underset{j=2}{\overset{q}{}}c_jr_j(t)r_{k+2}(t)=y(t)\sigma _ky(t),$$
the estimations (4) and (5) imply
$$f_Xy_X+\sigma _ky_X(B+1)y_X,$$
and consequently,
$$z_X\sigma _ky_X+f_X(2B+1)y_X.$$
$`(6)`$
Case 2. $`k<l`$.
Now set
$$g(t)=\sigma _{k+1}y(t)z(t)=\underset{j=p+1}{\overset{k+1}{}}d_jr_j(t)r_{k+2}(t)+d_{k+2}r_{k+3}(t).$$
Define $`f,u`$ and $`v`$ as above. Inequality (5) holds in this case also. Since
$$u=\sigma _{k+1}y\sigma _ky,$$
inequality (4) implies
$$f_X\sigma _{k+1}y_X+\sigma _ky_X\mathrm{\hspace{0.33em}2}By_X.$$
We have now
$$z_X\sigma _ky_X+f_X\mathrm{\hspace{0.33em}3}By_X.$$
$`(6^{})`$
The definitions of the functions $`z`$ and $`y`$, together with inequalities (6) and (6’), yield that relation (1) holds true for the Rademacher chaos which is ordered according to (3). The theorem is proved.
Remark 1. The requirement for the space $`X`$ to be an interpolation space with respect to the couple $`(L_1,L_{\mathrm{}})`$ is not very restrictive. The most important s.s. (Orlicz, Lorentz, Marcinkiewicz spaces and others) possess this property \[8, p.142\]. In addition, it is seen from the proof of the theorem that the above condition may be replaced by a weaker one: the boundedness in $`X`$ of the averaging operators corresponding to the dyadic partitionings of the interval $`[0,1]`$.
## 3 Rademacher chaos as unconditional basic sequence
We go now further to the study of the unconditionality of Rademacher chaos in s.s. We have already mentioned that the main result in this paper amplifies Theorem A proved in and formulated in Section 1.
Theorem 2. Let $`X`$ be s.s. on $`[0,1]`$. Then the following assertions are equivalent:
$`1)`$ The system $`\{r_ir_j\}_{1i<j<\mathrm{}}`$ in $`X`$ is equivalent to the canonical basis in the space $`l_2`$, that is, there is a constant $`C>0`$ that depends only on the space $`X;`$ such that for all real numbers $`a_{i,j}(1i<j<\mathrm{})`$,
$$C^1(a_{i,j})_2\underset{1i<j<\mathrm{}}{}a_{i,j}r_ir_j_XC(a_{i,j})_2.$$
$`(7)`$
$`2)`$ A continuous imbedding $`HX`$ takes place;
$`3)`$ The system $`\{r_ir_j\}_{1i<j<\mathrm{}}`$ is an unconditional basic sequence in $`X`$.
Remark 2. The implication $`1)3)`$ is evident and the equivalence $`1)2)`$ is proved in . Thus, it suffices to prove the implication $`3)1)`$.
First, we prove a weaker assertion. Let $`G`$ denote the closure of $`L_{\mathrm{}}`$ in the Orlicz space $`L_N`$ corresponding to the function $`N(t)=e^{t^2}1`$.
Proposition 1. Let the s.s. $`X`$ on $`[0,1]`$ be such that $`XG`$ and the system $`\{r_ir_j\}_{1i<j<\mathrm{}}`$ is an unconditional basic sequence in $`X`$. Then the assertion $`1)`$ in Theorem 2 holds true, that is, the system $`\{r_ir_j\}_{1i<j<\mathrm{}}`$ in $`X`$ is equivalent to the canonical basis in $`l_2`$.
For the proof we need a lemma that concerns spaces with a mixed norm. Let us recall the definition (for details, see \[5, p.400\] ).
Definition. Let $`X`$ and $`Y`$ be s.s. on $`[0,1]`$. The space with a mixed norm $`X[Y]`$ is the set of all measurable functions $`x(s,t)`$ on the square $`I\times I`$ satisfying the conditions:
1) $`x(,t)Y`$ for almost all $`tI`$;
2) $`\phi _x(t):=x(,t)_YX.`$
Define
$$x_{X[Y]}=\phi _x_X.$$
Let $`A=A(u)`$ be a $`N`$-function on $`[0,\mathrm{})`$. This means that $`A`$ is continuous, convex, and satisfies
$$\underset{u+0}{lim}\frac{A(u)}{u}=\underset{u+\mathrm{}}{lim}\frac{u}{A(u)}=\mathrm{\hspace{0.33em}0}.$$
As usual, denote by $`L_A`$ the Orlicz space of all functions $`x=x(t)`$ measurable on $`[0,1]`$ and having a finite norm,
$$x_{L_A}:=inf\{\lambda >0:_0^1A\left(\frac{|x(t)|}{\lambda }\right)𝑑t\mathrm{\hspace{0.17em}1}\}.$$
Finally, let $`A^{}`$ be the $`N`$-function conjugated to the $`N`$-function $`A`$, that is,
$$A^{}(u):=sup\{uvA(v):v0\}.$$
Lemma 1. The following imbeddings take place
$$L_{\mathrm{}}[L_A]L_A(I\times I),$$
$$L_A(I\times I)L_1[L_A],$$
where $`X(I\times I)`$ denotes a s.s. on the square $`I\times I`$.
Proof. If $`x=x(s,t)L_{\mathrm{}}[L_A]`$, then, according to the definition of the norm in an Orlicz space, for almost $`t[0,1]`$ we have
$$_0^1A\left(\frac{|x(s,t)|}{C}\right)𝑑s\mathrm{\hspace{0.33em}1}$$
where $`C=x_{L_{\mathrm{}}[L_A]}`$. After integrating this inequality and applying Fubini’s theorem we get
$$_0^1_0^1A\left(\frac{|x(s,t)|}{C}\right)𝑑s𝑑t\mathrm{\hspace{0.33em}1}.$$
Therefore $`xL_A(I\times I)`$ and $`x_{L_A(I\times I)}x_{L_{\mathrm{}}[L_A]}`$. The first imbedding is proved.
For the proof of the second imbedding we pass on to the dual spaces. Recall that the dual space $`X^{}`$ to the s.s. $`X`$ consists of all measurable functions $`y=y(t)`$ for which
$$y_X^{}:=sup\{_0^1x(t)y(t)𝑑t:x_X1\}<\mathrm{}.$$
We have already proved that
$$L_{\mathrm{}}[L_A^{}]L_A^{}(I\times I).$$
Therefore, for the dual spaces we have
$$\left(L_A^{}(I\times I)\right)^{}\left(L_{\mathrm{}}[L_A^{}]\right)^{}.$$
Finally, since $`(L_A)^{}=L_A^{}`$, $`(A^{})^{}=A`$ \[7, p.146, p.22\] and $`(X[Y])^{}=X^{}[Y^{}]`$ \[3, Th.3.12\], it follows that
$$L_A(I\times I)L_1[L_A].$$
Besides,
$$x_{L_1[L_A]}Cx_{L_A(I\times I)}$$
for a certain $`C>0`$.
Proof of Proposition 1. Let $`𝐫_{i,j}(u)`$ $`(1i<j<\mathrm{})`$ denote Rademacher functions, arbitrarily ordered by the couples $`(i,j)`$.
Since $`XG`$, by the lemma, for any given $`n`$ and real numbers $`a_{i,j}`$, we get
$$_0^1\underset{1i<jn}{}a_{i,j}𝐫_{i,j}(u)r_ir_j_X𝑑u=$$
$$=\underset{1i<jn}{}a_{i,j}𝐫_{i,j}(u)r_i(t)r_j(t)_{X(t)}_{L_1(u)}$$
$$C_1\underset{1i<jn}{}a_{i,j}𝐫_{i,j}(u)r_i(t)r_j(t)_{G(u)}_{L_{\mathrm{}}(t)}$$
(this follows from the fact that $`G`$ is a subspace of $`L_N`$ and thus, for the functions in $`L_{\mathrm{}}(I\times I)`$, the norms in the spaces $`L_1[G]`$ and $`L_1[L_N]`$ coincide and so do the norms in the spaces $`L_{\mathrm{}}[G]`$ and $`L_{\mathrm{}}[L_N]`$). By Khintchine’s inequality for the space $`G`$ (see, for example, ), this yields
$$_0^1\underset{1i<jn}{}a_{i,j}𝐫_{i,j}(u)r_ir_j_X𝑑uC_2(a_{i,j})_2.$$
$`(8)`$
On the other hand, as is known \[10, Ch.4\],
$$\underset{1i<j<\mathrm{}}{}a_{i,j}r_ir_j_{L_p}(a_{i,j})_2$$
for any $`p[1,\mathrm{})`$. (This means that a constant $`C>0`$ exists depending only on $`p`$ such that
$$C^1(a_{i,j})_2\underset{1i<j<\mathrm{}}{}a_{i,j}r_ir_j_{L_p}C(a_{i,j})_2).$$
Therefore, by the imbedding $`XL_1`$ which holds for each s.s. $`X`$ on $`[0,1]`$ \[8, p.124\], we get
$$\underset{1i<j<\mathrm{}}{}a_{i,j}r_ir_j_XC_3(a_{i,j})_2.$$
$`(9)`$
Thus, the inequality
$$_0^1\underset{1i<jn}{}a_{i,j}𝐫_{i,j}(u)r_ir_j_X𝑑uC_3(a_{i,j})_2,$$
$`(10),`$
which is opposite to (8), holds always true.
By the assumptions, with a constant depending only on the space $`X`$, we have
$$\underset{1i<jn}{}a_{i,j}r_ir_j_X_0^1\underset{1i<jn}{}a_{i,j}𝐫_{i,j}(u)r_ir_j_X𝑑u$$
for each $`n`$ and all real numbers $`a_{i,j}`$. In this way, the proposition follows from relations (8) and (10).
In (see also ) the notion of RUC (random unconditional convergence)-system was introduced. We shall give here an equivalent definition.
Definition. Let $`X`$ be a Banach space and let $`X^{}`$ be its dual space. The biorthogonal system $`(x_n,x_n^{})`$, $`x_nX`$, $`x_n^{}X^{}`$ $`(n=1,2,\mathrm{})`$ is said to be a RUC-system, if there exists a constant $`C>0`$ such that
$$_0^1\underset{i=1}{\overset{n}{}}r_i(s)x_i^{}(x)x_i_X𝑑sCx_X$$
for any $`n`$ and all $`x[x_n]_{n=1}^{\mathrm{}}`$ ($`r_i(s)`$ are Rademacher functions).
Inequalities (8) and (9) yield the following.
Corollary 1. If s.s. $`XG`$, then Rademacher chaos $`\{r_ir_j\}_{1i<j<\mathrm{}}`$ together with the basic coefficients is a RUC-system in $`X`$.
Corollary 2. For each s.s. $`X`$ the following assertions are equivalent:
$`1)XG;`$
$`2)_0^1_{1i<j<\mathrm{}}a_{i,j}𝐫_{i,j}(u)r_ir_j_Xdu(a_{i,j})_2`$.
Proof. The implication $`1)2)`$ follows from inequalities (8) and (10).
Suppose now that 2) takes place and let $`a_{i,j}=0`$ $`(i1)`$. From the definition of Rademacher functions and the assumptions we get
$$_0^1\underset{j=2}{\overset{\mathrm{}}{}}a_{1,j}r_{1,j}(u)r_1r_j_X𝑑u=\underset{j=2}{\overset{\mathrm{}}{}}a_{1,j}r_j_XC\left(\underset{j=2}{\overset{\mathrm{}}{}}a_{1,j}^2\right)^{1/2}.$$
Therefore (see \[12, p.134\] or ) $`XG`$.
Corollary 3. Suppose that s.s. $`XG`$. Then, for any set $`\{a_{i,j}\}_{1i<j<\mathrm{}}`$ of real numbers, there exists a set of signs $`\{\theta _{i,j}\}_{1i<j<\mathrm{}}`$, $`\theta _{i,j}=\pm 1`$, such that
$$\underset{1i<j<\mathrm{}}{}\theta _{i,j}a_{i,j}r_ir_j_X(a_{i,j})_2.$$
Proof. Corollary 2 yields
$$inf\{\underset{1i<j<\mathrm{}}{}\theta _{i,j}a_{i,j}r_ir_j_X:\theta _{i,j}=\pm 1\}$$
$$_0^1\underset{1i<j<\mathrm{}}{}a_{i,j}𝐫_{i,j}(u)r_ir_j_X𝑑uC(a_{i,j})_2.$$
Now the assertion follows from the fact that the opposite inequality holds always (see (10)).
In order to prove Theorem 2 we need some more auxiliary assertions.
Let $`\{n_k\}_{k=1}^{\mathrm{}}`$ be an increasing sequence of natural numbers. Set $`t_k=\mathrm{\hspace{0.17em}2}^{n_{k+1}}`$, $`m_k=(n_{k+1}n_k)(n_{k+1}n_k1)/2`$ and
$$y_k(t)=\underset{n_k<i<jn_{k+1}}{}r_i(t)r_j(t)(k=1,2,\mathrm{}).$$
Lemma 2. Let $`c_k>0`$ $`(k=1,2,\mathrm{})`$,
$$y(t)=\underset{k=1}{\overset{\mathrm{}}{}}c_ky_k(t)(t[0,1]).$$
If $`y^{}(t)`$ is a non-increasing rearrangement of the function $`|y(t)|`$ \[8, p.83\] , then
$$y^{}(t_k)\underset{l=1}{\overset{k}{}}m_lc_l.$$
$`(11)`$
Proof. By the definition of Rademacher functions $`y_l(t)=m_l`$ provided that $`0<t<2t_k`$ and $`1lk`$. Therefore
$$y(t)=\underset{l=1}{\overset{k}{}}m_lc_l+\underset{l=k+1}{\overset{\mathrm{}}{}}c_ly_l(t)(0<t<2t_k).$$
Besides, there exists a set $`E(0,2t_k)`$ of Lebesgue measure $`|E|=t_k`$ such that
$$\underset{l=k+1}{\overset{\mathrm{}}{}}c_ly_l(t)\mathrm{\hspace{0.33em}0}\text{for }tE.$$
Applying the previous equality we get
$$y(t)\underset{l=1}{\overset{k}{}}c_lm_l\text{ for }tE.$$
Inequality (11) follows now from the definition of the rearrangement and the fact that $`|E|=t_k`$.
The next assertion makes Theorem 8 from more precise. We shall use here the same notations as in .
If $`X`$ is a s.s. on $`[0,1]`$, then $`\overline{}(X)`$ denotes a subspace of $`X`$ consisting of all functions of the form
$$x(t)=\underset{1i<j<\mathrm{}}{}a_{i,j}r_i(t)r_j(t),(a_{i,j})_{1i<j<\mathrm{}}l_2.$$
For any arrangement of signs (that is, for any sequence $`\theta =\{\theta _{i,j}\}_{1i<j<\mathrm{}}`$, $`\theta _{i,j}=\pm 1`$), we define the operator
$$\overline{T}_\theta x(t)=\underset{1i<j<\mathrm{}}{}\theta _{i,j}a_{i,j}r_i(t)r_j(t)$$
on the subspace $`\overline{}(X)`$,
Proposition 2. There exists an arrangement of signs $`\theta =\{\theta _{i,j}\}_{1i<j<\mathrm{}}`$ such that for each $`\epsilon (0,1/2)`$ one can find a function $`x\overline{}(L_{\mathrm{}})`$ satisfying
$$(\overline{T}_\theta x)^{}(t)b\mathrm{log}_2^{1/2\epsilon }2/t$$
with a constant $`b>0`$ independent of $`t(0,1/16]`$.
Proof. By Theorem 6 in (see also Lemma 3 there), for each $`k=1,2,\mathrm{}`$ one can find $`\theta _{i,j}=\pm 1`$ $`(2^k<i<j2^{k+1})`$ such that the functions
$$z_k(t)=\underset{2^k<i<j2^{k+1}}{}\theta _{i,j}r_i(t)r_j(t)$$
satisfy
$$z_k_{\mathrm{}}\mathrm{\hspace{0.33em}2}^{3k/2}.$$
$`(12)`$
Set $`x_k(t)=\mathrm{\hspace{0.17em}2}^{(3+2\epsilon )k/2}z_k(t)`$ and
$$x(t)=\underset{k=1}{\overset{\mathrm{}}{}}x_k(t)=\underset{1i<j<\mathrm{}}{}a_{i,j}r_i(t)r_j(t),$$
where $`a_{i,j}=\mathrm{\hspace{0.17em}2}^{(3+2\epsilon )k/2}\theta _{i,j}`$, if $`2^k<i<j2^{k+1},k=1,2,\mathrm{}`$, and $`a_{i,j}=0`$, otherwise. It follows from (12) that
$$x_{\mathrm{}}C\underset{k=1}{\overset{\mathrm{}}{}}2^{\epsilon k}=C/(2^\epsilon 1),$$
that is, $`x\overline{}(L_{\mathrm{}})`$.
Let the arrangement of signs $`\theta `$ consist of the values $`\theta _{i,j}`$ just determined for $`2^k<i<j2^{k+1},k=1,2,\mathrm{}`$, and arbitrary $`\theta _{i,j}`$ for the other couples $`(i,j),i<j`$. We get
$$y(t)=\overline{T}_\theta x(t)=\underset{k=1}{\overset{\mathrm{}}{}}2^{(3+2\epsilon )k/2}y_k(t)$$
where
$$y_k(t)=\underset{2^k<i<j2^{k+1}}{}r_i(t)r_j(t)(k=1,2,\mathrm{}).$$
Next we apply Lemma 2 to the case when $`n_k=\mathrm{\hspace{0.17em}2}^k`$ and $`c_k=\mathrm{\hspace{0.17em}2}^{(3+2\epsilon )k/2}`$. Then clearly $`t_k=2^{2^{k+1}}`$ and $`m_k2^{2k2}`$. Therefore, for each $`k=1,2,\mathrm{}`$, we have
$$y^{}(t_k)\frac{1}{4}\underset{i=1}{\overset{k}{}}2^{2i}2^{(3+2\epsilon )i/2}\mathrm{\hspace{0.33em}2}^{(1/2\epsilon )k2}C_1\mathrm{log}_2^{1/2\epsilon }(2/t_k).$$
Given an arbitrary $`t(0,1/16]`$, one can find a $`k`$ so that $`t_{k+1}<tt_k`$. Taking into account the previous inequality, we get
$$y^{}(t)y^{}(t_k)C_1\mathrm{log}_2^{1/2\epsilon }(2/t_k)$$
$$\mathrm{\hspace{0.33em}4}^{\epsilon 1/2}C_1\mathrm{log}_2^{1/2\epsilon }(2/t_{k+1})b\mathrm{log}_2^{1/2\epsilon }2/t.$$
The proposition is proved.
Recall that Marcinkiewicz s.s. $`M(\phi )`$ ($`\phi (t)0`$ is a concave increasing function on $`[0,1]`$ ) consists of all measurable functions $`x(s)`$ having a finite norm
$$x_{M(\phi )}:=sup\{\frac{1}{\phi (t)}_o^tx^{}(s)𝑑s:\mathrm{\hspace{0.17em}0}<t1\}.$$
Corollary 4. Let $`X`$ be s.s. on $`[0,1]`$ such that, for any arrangement of signs $`\theta =\{\theta _{i,j}\}_{1i<j<\mathrm{}}`$, the operator $`\overline{T}_\theta `$ is bounded in $`\overline{}(X)`$. Then
$$X\underset{\epsilon (0,1/2)}{}M(\phi _\epsilon )$$
where $`M(\phi _\epsilon )`$ is Marcinkiewicz space determined by $`\phi _\epsilon (t)=t\mathrm{log}_2^{1/2\epsilon }2/t`$.
Proof. Making use of Proposition 2, we can find an arrangement of signs $`\theta `$ such that for a given number $`\epsilon (0,1/2)`$ and some function $`x\overline{}(L_{\mathrm{}})`$ ($`x`$ depends on $`\epsilon `$),
$$(\overline{T}_\theta x)^{}(t)b\mathrm{log}_2^{1/2\epsilon }2/t(0<t1/16).$$
By the fact that $`XL_{\mathrm{}}`$ for any s.s. $`X`$ on $`[0,1]`$ \[8, p.124\] and taking into account the assumption and the symmetry of $`X`$, we get
$$\overline{x}_\epsilon (t):=\mathrm{log}_2^{1/2\epsilon }2/tX.$$
$`(13)`$
Now it follows from the relation \[8, p.156\]
$$y_{M(\phi _\epsilon )}sup\{y^{}(t)\mathrm{log}_2^{1/2\epsilon }2/t:\mathrm{\hspace{0.17em}0}<t1\}.$$
that $`\overline{x}_\epsilon (t)`$ has maximal rearrangement in the space $`M(\phi _\epsilon )`$. Besides, by virtue of (13), $`XM(\phi _\epsilon )`$. The corollary is proved.
Proposition 3. An arrangement of signs $`\theta =\{\theta _{i,j}\}_{1i<j<\mathrm{}}`$ exists such that for each $`\epsilon (1/2,1/2)`$ and any $`\delta (0,1/4\epsilon /2)`$ one can find a function $`x\overline{}(M(\phi _\epsilon ))`$ satisfying
$$(\overline{T}_\theta x)^{}(t)d\mathrm{log}_2^{1/2+\delta }(2/t)$$
$`(14)`$
where the constant $`d>0`$ does not depend on $`t(0,1/16]`$.
Proof. Let $`\theta =\{\theta _{i,j}\}_{1i<j<\mathrm{}}`$, $`z_k`$ and $`y_k`$ $`(k=1,2,\mathrm{})`$ be defined in the same way as in the proof of Proposition 2. It is well-known that Marcinkiewicz space $`M(\phi _{1/2})=M(t\mathrm{log}_2(2/t))`$ coincides with the Orlicz space $`L_M`$, $`M(t)=e^t1`$. Therefore, by Theorem A from Section 1, we have
$$z_k_{M(\phi _{1/2})}\left(\underset{2^k<i<j2^{k+1}}{}|\theta _{i,j}|^2\right)^{1/2}\mathrm{\hspace{0.33em}2}^k(k=1,2,\mathrm{}).$$
$`(15)`$
For any $`u(0,1)`$ the space $`M(\phi _\epsilon )`$ with $`\epsilon =(12u)/2`$ is a space of the type $`u`$ with respect to the couple $`(L_{\mathrm{}},M(\phi _{1/2}))`$, that is, a constant $`C>0`$ exists such that
$$x_{M(\phi _\epsilon )}Cx_{\mathrm{}}^{1u}x_{M(\phi _{1/2})}^u$$
$`(16)`$
for all $`xL_{\mathrm{}}`$. In fact, as we have already mentioned in the proof of Corollary 4,
$$x_{M(\phi _\epsilon )}C^{}sup\{x^{}(t)\mathrm{log}_2^u(2/t):\mathrm{\hspace{0.17em}0}<t1\}$$
$$C^{}\left[sup\{x^{}(t):\mathrm{\hspace{0.17em}0}<t1\}\right]^{1u}\left[sup\{x^{}(t)\mathrm{log}_2^1(2/t):\mathrm{\hspace{0.17em}0}<t1\}\right]^u$$
$$Cx_{\mathrm{}}^{1u}x_{M(\phi _{1/2})}^u.$$
From (12), (15) and (16) we get
$$z_k_{M(\phi _\epsilon )}D\mathrm{\hspace{0.17em}2}^{(3u)k/2}(k=1,2,\mathrm{}).$$
$`(17)`$
For a given $`v>0`$ (to be determined later), set
$$x_k(t)=\mathrm{\hspace{0.33em}2}^{(u32v)k/2}z_k(t),x(t)=\underset{k=1}{\overset{\mathrm{}}{}}x_k(t).$$
By virtue of (17), $`x\overline{}(M(\phi _\epsilon ))`$. Applying Lemma 2 to the function
$$y(t)=\overline{T}_\theta x(t)=\underset{k=1}{\overset{\mathrm{}}{}}2^{(u32v)k/2}y_k(t)$$
for $`n_k=\mathrm{\hspace{0.17em}2}^k`$, $`c_k=\mathrm{\hspace{0.17em}2}^{(u32v)k/2}`$, we get
$$y^{}(t)\frac{1}{4}\underset{i=1}{\overset{k}{}}2^{2i}\mathrm{\hspace{0.17em}2}^{(u2v3)i/2}$$
$$C_1^{}\mathrm{\hspace{0.17em}2}^{(1+u2v)k/2}C_1\mathrm{log}_2^{(1+u2v)/2}(2/t_k)(k=1,2,\mathrm{}).$$
In the same way as in the proof of Proposition 2 we conclude that
$$y^{}(t)D_1\mathrm{log}_2^{(1+u2v)/2}(2/t)(0<t1/16).$$
Since $`\delta <1/4\epsilon /2`$ by assumption, we can choose $`v`$ so that
$$0<v<1/4\epsilon /2\delta .$$
Therefore $`(u2v)/2>\delta `$ and inequality (14) holds with some $`b>0`$.
Corollary 5. Let $`X`$ be s.s. on $`[0,1]`$ such that $`XM(\phi _\epsilon )`$ for some $`\epsilon (1/2,1/2)`$. If for any arrangement of signs $`\theta =\{\theta _{i,j}\}_{1i<j<\mathrm{}}`$ the operator $`\overline{T}_\theta `$ is bounded in $`\overline{}(X)`$, then
$$X\underset{0<\delta <1/4\epsilon /2}{}M(\phi _\delta ).$$
The proof is similar to that of Corollary 4.
Now we are ready to prove our main Theorem 2.
Proof of Theorem 2. As it has been already mentioned in Remark 2, it suffices to verify the implication $`3)1)`$.
If the system $`\{r_ir_j\}_{1i<j<\mathrm{}}`$ is unconditional in s.s. $`X`$, then, for each arrangement of signs $`\theta `$, the operator $`\overline{T}_\theta `$ is bounded in $`\overline{}(X)`$. In particular, we get by Corollary 4 that $`XM(\phi _{1/5})`$. Therefore, by Corollary 5, it follows that $`XM(\phi _{1/10})`$. Since $`M(\phi _{1/10})G`$, then all the more, $`XG`$. Finally, applying Proposition 1 we conclude that the system $`\{r_ir_j\}_{1i<j<\mathrm{}}`$ is equivalent to the canonical basis in $`l_2`$. The theorem is proved.
Remark 3. Assertions analogous to Theorems 1 and 2 are valid for the multiple Rademacher system $`\{r_i(s)r_j(t)\}_{i,j=1}^{\mathrm{}}`$ considered on the square $`I\times I`$, $`I=[0,1]`$, as well. This follows from the equivalence of the symmetric norms for the series with respect to the systems $`\{r_ir_j\}_{1i<j<\mathrm{}}`$ and $`\{r_i(s)r_j(t)\}_{i,j=1}^{\mathrm{}}`$ (see , and also ).
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# Modelling High-frequency Economic Time Series
## Abstract
The minute-by-minute move of the Hang Seng Index (HSI) data over a four-year period is analysed and shown to possess similar statistical features as those of other markets. Based on a mathematical theorem \[S. B. Pope and E. S. C. Ching, Phys. Fluids A 5, 1529 (1993)\], we derive an analytic form for the probability distribution function (PDF) of index moves from fitted functional forms of certain conditional averages of the time series. Furthermore, following a recent work by Stolovitzky and Ching, we show that the observed PDF can be reproduced by a Langevin process with a move-dependent noise amplitude. The form of the Langevin equation can be determined directly from the market data.
To appear in Proceedings of the Dynamics Days Asia Pacific Conference, 13-16 July, 1999, Hong Kong (Physica A, 2000).
The availability of high-frequency economic time series, with a sampling rate of every few seconds, has generated a great deal of theoretical interest in the econometrics and the econophysics community. Attempts have been made to devise models which produce time series with similar statistical characteristics as those of real markets. Many of these studies are based on variants of the Autogressive Conditional Heteroskedasticity (ARCH) process first introduced by Engle to analyze the quarterly consumer price index in the UK over the period 1958 to 1977, and the generalised ARCH (GARCH) process which offers a more flexible description of the volatility memory effect (i.e., lag structure). The nonlinearity in the regression models makes it possible to generate probability distribution functions (PDF) with fat tails, a characteristic of financial data first noted by Mandelbrot. However, since all these processes are discrete in time, an immediate question to ask is whether the quality of the modelling depends on the time unit chosen and if there is a time scale which is the most natural of all. Indeed, when the time step is not chosen properly, one has to either introduce many terms in the regression expression \[the GARCH($`p,q`$) model\] to take into account memory effects, or to miss some of the important short-time statistics.
An alternative approach, which partially circumvents the above difficulty, is to model the market price move as a continuous time process. Continuous time stochastic processes are quite familiar to physicists, ranging from simple Brownian motion to the fully-developed turbulence. In fact, the high-frequency market price movements have much in common with the velocity or temperature time series in turbulent flows, an analogy we exploit in this paper. To put this statement on more quantitative terms, let us first summarise two salient statistical features which seem to be universally true for all major stock indices.
(i) Short linear correlation of price moves — For a given stock index $`S(t)`$, one may define the price move over a fixed time interval $`\delta `$ (say one minute),
$$x(t)=S(t)S(t\delta ).$$
(1)
It has been shown that the “linear correlation”
$$C(\tau )=x(t+\tau )x(t)$$
(2)
decays to zero very rapidly, on the order of ten minutes. We have analysed the minute-by-minute Hang Seng Index (HSI) data collected over a four year period January 1994 — December, 1997. Figure 1 shows the linear correlation function $`C(\tau )`$ with $`\delta =1`$ min. It is seen that $`C(\tau )`$ becomes nearly zero after a period of ten minutes or so. The decay is however not completely monotone, suggesting that the market is slightly under-damped.
(ii) Nongaussian distribution of price moves with fat tails — The fat tails of the PDF $`P(x)`$ of stock price moves are well-known and have also been observed for the movement of foreign currency exchange rates. Mandelbrot has observed that $`P(x)`$ often decays as a power-law function of $`|x|`$, and hence, combining with (i), the stock index can be considered as a realisation of Lévy walk. From the analysis of the high-frequency S&P 500 data, Mantegna and Stanley showed that a truncated Lévy distribution offers a better description of the PDF. Figure 2 shows the PDF for the HSI minute-by-minute move data $`x(t)`$ (open circles) on a semi-log scale, collected over the same period as in Fig. 1. It is seen that the decay at large $`|x|`$ can be well-described by a simple exponential function, as observed previously in Ref. . For small $`|x|`$, a different behaviour is seen. A noticeable feature is the cusp-like singularity at $`x=0`$, which is so far unexplained.
The peculiar form of the PDF as seen in Fig. 2 has in fact been observed in a physical context. In analysing the temperature time series of thermal convection in the “hard-turbulence” regime, Pope and Ching considered the following conditional averages for a twice-differentiable time series $`x(t)`$,
$$r(x)=\frac{\ddot{x}|x}{\dot{x}^2},q(x)=\frac{\dot{x}^2|x}{\dot{x}^2}.$$
(3)
Here $`|x`$ denotes the average of a given quantity over those data points in the time series where $`x(t)=x`$. From the stationarity of the PDF, they proved that the PDF and the conditional averages $`r(x)`$ and $`q(x)`$ are related through the following equation,
$$P(x)=\frac{C}{q(x)}\mathrm{exp}\left[_0^x\frac{r(x^{})}{q(x^{})}𝑑x^{}\right].$$
(4)
Using the turbulent temperature time series data as input, they showed that $`r(x)`$ is generally linear in $`x`$ with a negative coefficient. In the soft turbulence regime, $`q(x)`$ is nearly constant. From Eq. (4), the resulting PDF is gaussian as observed. On the other hand, in the hard turbulence regime, $`q(x)`$ increases with increasing $`|x|`$, giving rise to fat tails in the PDF.
Figure 3 shows $`r(x)`$ and $`q(x)`$ computed using the HSI move time series. Indeed, the shape of these two functions are very much like the temperature data in Ref. , although there are small differences. The data can be fitted to the following functional forms,
$$r(x)=Rx,q(x)=Q(x^2+a^2)^{1/2},$$
(5)
where $`R=0.036`$ and $`Q=0.36`$. The round-off parameter $`a`$ can not be determined precisely because of the discreteness of the index data. Instead, we fix its value from the normalisation condition. Substituting Eq. (5) into Eq. (4), and carrying out the integration, we obtain,
$$P(x)=\frac{C^{}}{(x^2+a^2)^{1/2}}\mathrm{exp}\left[\alpha (x^2+a^2)^{1/2}\right],$$
(6)
where $`\alpha =R/Q=0.1`$ and $`C^1=2K_0(\alpha a)`$, with $`K_0(u)`$ being a modified Bessel function of the second kind. The solid line in Fig. 2 is produced by taking $`a^2=0.1`$. As can be seen, the agreement between the original data and the fitted form (6) is rather satisfactory.
It is worth noting that the functional form (6) we propose is quite different from those of earlier studies. The behaviour of the PDF around $`x=0`$ is controlled by the parameter $`a`$. A sharp peak is produced when $`a`$ becomes very small. In this respect, $`a`$ serves the purpose of a cut-off related to the discreteness of the underlying asset price. For large $`|x|`$, $`P(x)`$ crosses over to simple exponential decay. There is however no scale invariance as previously suggested.
The more challenging task is to devise a dynamic equation that generates a time series with the same conditional averages as those of the market data. This issue was considered recently by Stolovitzky and Ching. They studied a one-dimensional Langevin process defined by the following stochastic differential equation,
$$m\frac{d^2x}{dt^2}+\gamma \frac{dx}{dt}=F(x)+[2\gamma k_BT(x)]^{1/2}\xi (t),$$
(7)
where $`\xi (t)`$ is a gaussian white noise with $`\xi (t)=0`$ and $`\xi (t)\xi (t^{})=\delta (tt^{})`$. The main difference of (7) from the usual Brownian process is an $`x`$-dependent temperature (or noise amplitude) $`T(x)`$. In the over-damped limit $`\gamma \mathrm{}`$, they showed that the conditional averages are given by,
$$\ddot{x}|x=F(x)/m,\dot{x}^2|x=k_BT(x)/m.$$
(8)
Combining with Eq. (4), Stolovitzky and Ching showed that the PDF in this limit is given by a generalised Boltzmann form,
$$P(x)=\frac{C}{T(x)}\mathrm{exp}\left[_0^x\frac{F(x^{})}{k_BT(x^{})}𝑑x^{}\right].$$
(9)
The significance of the above result is as follows. Assuming that a given time series is generated by a Langevin process (7), one can then determine the effective force $`F(x)`$ and the effective temperature $`T(x)`$ uniquely by computing the conditional averages from the data, apart from an overall time constant. The PDF of the Langevin time series is identical to the PDF of the original data by construction.
The $`\gamma \mathrm{}`$ limit is quite suitable for the analysis of the HSI data, as we have seen from the two-point correlation function $`C(\tau )`$ that any memory effect about the direction of the move decays to zero rapidly. It is then suggestive to drop the inertia (i.e. mass) term in Eq. (7) altogether. Performing the scaling $`t\gamma t`$ and setting $`k_B=1`$, we can cast Eq. (7) in the form,
$$\frac{dx}{dt}=F(x)+[2T(x)]^{1/2}\xi (t).$$
(10)
The effective force and the effective temperature are related to the conditional averages through Eq. (8). The parameter $`m`$ can be chosen to include a short-time relaxation effect as seen in Fig. 1.
We have simulated the Langevin equation (10) using an Euler integration scheme with $`\mathrm{\Delta }t=0.1`$ min. The form of the functions $`F(x)`$ and $`T(x)`$ are determined from the conditional averages (5). The PDF of the simulated minute-by-minute moves, with the same number of events as the original data, is shown in Fig. 2 (crosses). In Fig. 4 we plot both the HSI data (daily close) and the simulated index data with an artificial annual yield of $`10\%`$. The gross features of the two data sets seem to be similar to the naked eye, though in the simulated data set there is no daily and weekly breaks.
One important feature which is missing in the Langevin equation (7) is the long-term volatility correlations which may last from a few days to several weeks or longer. The simulated time series has essentially the same relaxation time, on the order of a few minutes in our case, for the linear move correlation and for the volatility correlation. It is however possible to introduce volatility persistence by hand into the simulation. The effect of a nonstationary volatility on the PDF of the time series is a subject under current investigation.
Acknowledgements: We would like to thank Prof. Lam Kin at the Department of Finance and Decision Sciences, HK Baptist University for providing the HSI data, and Dr. Emily Ching for sending us Ref. prior to publication. The work is supported in part by the Hong Kong Baptist University under grant FRG/97-98/II-78.
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# HST Observations of the Interacting Galaxies NGC 2207 and IC 21631footnote 11footnote 1 Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract NAS5-26555
## 1 Introduction
The spiral galaxies IC 2163 and NGC 2207 are currently involved in a near-grazing encounter (Elmegreen, et al. 1995a,b; hereafter Papers III and IV). Tidal forces distorted IC 2163 in the inplane direction, forming a tidal arm with two velocity components on the anticompanion side, and forming an intrinsically oval disk with an eye–shaped (“ocular”) morphology. Streaming motions in the oval are in excess of 100 km s<sup>-1</sup>. Tidal forces also distorted NGC 2207 perpendicular to the plane, forming a strong, twisting warp. The H I in both galaxies has a large velocity dispersion of 30 - 50 km s<sup>-1</sup> and is concentrated in unusually large clouds ($`10^8`$ M; Elmegreen, Kaufman & Thomasson 1993; hereafter Paper II). Numerical simulations reproduced these peculiar structures and internal velocities with a prograde encounter affecting IC 2163 and a perpendicular encounter affecting NGC 2207 (Paper IV). The closest approach was around 40 My ago at a separation of only $`2`$ radii. This makes IC 2163 and NGC 2207 ideal for studying close encounters between galaxies and how they might trigger bars or starbursts.
Three systems of interacting galaxies (IC 2163, NGC 2535, and NGC 5394) have now been studied by our group to see how structures resulting from prograde, inplane, nearly-grazing encounters evolve in time (Papers III, IV; Kaufman et al. 1997; Kaufman et al. 1999). IC 2163 is the interaction which is least evolved. In all three cases, the prograde galaxy has two long tidal arms with a large arm/interarm contrast, misalignment between the kinematic and photometric axes indicative of an intrinsically oval disk, and widespread high velocity dispersions in the H I gas. NGC 2535 is an ocular galaxy but, unlike IC 2163, it has no enhanced HI or radio continuum emission on the rim of the oval; this indicates it is in a later phase than IC 2163 (Elmegreen et al. 1991; hereafter Paper I). There is also evidence for mass transfer from NGC 2535 to its small starburst companion. NGC 5394 has no ocular structure, but bright spiral arms in the inner disk and a nuclear starburst, making it more evolved than NGC 2535.
To study these grazing encounters in more detail, we observed the IC 2163/N2207 pair with the WFPC2 camera on the Hubble Space Telescope (HST) in May 1996 and November 1998. Based on these observations, this paper describes (a) an improved understanding of the formation mechanism of the tidal tail in IC 2163, with a detailed numerical simulation of the gas (Sect. 4.1 and 4.2); (b) a new view of the density-wave structure and star formation in a spiral arm, as seen in a backlit portion of NGC 2207 (Sect. 4.3; (c) bright star formation associated with a strong nonthermal radio source on an outer spiral arm (Sect. 5); (d) the possibility that interactions can induce transient nuclear gas inflow via bar-like hydrodynamics, even if a genuine bar is absent (Sect. 6); (e) peculiar jet–like or curved emission features (Sect. 7); and (f) the HI – star formation connection (Sect. 8). Other studies of the star–forming regions and dust opacities will be reported elsewhere (D. Elmegreen et al. 2000, hereafter Paper V; see also Berlind, et al. 1997). The HST observations also revealed peculiar dust spirals in the nucleus of NGC 2207; these were analyzed in terms of acoustic instabilities by Elmegreen et al. (1998).
## 2 Observations and Data Reduction
The galaxies IC 2163 and NGC 2207 were observed in 9 orbits with the Hubble Space Telescope WFPC2 camera using filters F336W (U band), F439W (B band), F555W (V band), and F814W (I band). Four exposures in each band concentrated on a field in the northwest, and four more exposures of each of two fields were made for the central and eastern parts. The images were dithered between each exposure, and the exposure times were 500 s in U, 500 s in B, 160 s in V, and 180 s in I. The average scale is $`0.0995^{\prime \prime }`$ per pixel for the wide field images and $`0.0455^{\prime \prime }`$ per pixel for the planetary camera images.
From these 48 individual pipeline-processed exposures (four filters, three pointings, two dither positions per pointing, and two CR-splits per position), cosmic ray split pairs were combined for each dither position chip by chip. Individual WFPC2 chips were then combined for each dither position into $`2907\times 1486`$-pixel mosaics. The dither positions were registered and combined, resulting in a single image for each filter and pointing (twelve images). A 3-color composite image was produced (using IDL) for each pointing and the three pointings were registered and mosaicked (using Adobe PhotoShop) into a single image. The initial step in creating each 3-color image involved clipping and logarithmically scaling the intensities in each filter. Then each filter was assigned, in chromatic order, to a color channel: red=log(I), green=log(V), blue=log(B+U). The 3 channels of the mosaicked image were combined using the “screen” algorithm in PhotoShop. The channels were adjusted in order to match more closely the visual color of the central wavelength of each filter. This produced an overall bright image with a neutral color background, and it increased the intensity and color dynamic ranges in order to emphasize individual features. The noisier PC chip was Gaussian smoothed to match the texture of the other chips. Finally, the remaining artifacts (hot pixels, chip seams, etc.) were removed using Photoshop.
For comparison with H I and radio continuum images, a B–band mosaic was transformed to right ascension and declination coordinates. The uncertainty in the registration to absolute coordinates is about $`\pm 0.8^{\prime \prime }`$, partly as a result of uncertainties in the positions of the Guide Stars and secondary standard stars and partly as a result of the standard deviations of the plate solution. The B-band mosaic in absolute coordinates was rotated 10.7 counterclockwise relative to the color mosaic displayed in Figure 1 below.
## 3 Basic Morphology of the Interacting Galaxies
A mosaic of three HST WFPC2 pointings of the galaxy pair IC 2163/NGC 2207 is shown in Figure 1, with North up relative to the figure caption. A digitized version of this figure may be found at the Hubble Heritage web site (http://heritage.stsci.edu). Here we comment on several unusual features. The relevant parameters for these galaxies are listed in Table 1. At the distance of 35 Mpc (H=75 km s<sup>-1</sup> Mpc<sup>-1</sup>), $`1^{\prime \prime }`$ corresponds to 170 pc.
The morphology, interaction dynamics, and orbital history of this system were discussed, along with numerical simulations of each galaxy, in Papers III and IV. These previous simulations considered stars only, without gas. This was enough to get the basic disk structure. A new simulation of the tidal arm discussed below includes gas. The orbit that we fit previously had the smaller galaxy, IC 2163 (presently in the east), pass from the near side of NGC 2207 to the far side several hundred million years ago, at a point in the west about two NGC 2207-radii away from the center of NGC 2207. IC 2163 then moved behind NGC 2207 in an easterly direction, passing perigalacticon at a distance of about two of its own radii some $`40`$ My ago. IC 2163 is currently moving southeasterly. These relative distances are uncertain because they depend on the sizes and masses of the halos around each galaxy. The orbit is constrained by the internal structures and velocities of the galaxies, by their similar line-of-sight velocities (which implies that they are presently moving nearly parallel to each other in the sky plane), and by the extended pool of optical and H I emission that lies south of NGC 2207. This pool seems to be the point of the previous disk crossing, where IC 2163 passed through the outer disk plane of NGC 2207, now rotated clockwise by about 1/4 revolution.
During this interaction, IC 2163 experienced the tidal force from NGC 2207 in a prograde, inplane direction, and subsequently developed an ocular, or oval-like, caustic structure midway out in the disk. Such ocular structure is the result of a rapid inward motion of both stars and gas in the disk of IC 2163, initiated by the stretching and compressing action of the tidal force and amplified by self-gravity in the disk of IC 2163 (Sundin 1989; Paper I; Donner, Engstrom & Sundelius 1991). The inward motion causes an azimuthal speed-up in the counter-clockwise direction, and the radial motion then bounces because of angular momentum conservation. The oval marks the inner extent of this radial excursion. The minor kinematic axis of IC 2163 lies $`117^{}`$ away from the minor photometric axis, indicating an intrinsically oval disk.
There is a peculiar streaming motion along the oval of about 65 km s<sup>-1</sup> counter-clockwise (Paper III). The line–of–sight component of the peculiar velocity increases to $`150`$ km s<sup>-1</sup> as the H I gas and stars enter the tidal arm from the northwest, and then it decreases suddenly to $`<50`$ km s<sup>-1</sup> toward the outer edge of the tidal arm (Fig. 10, Paper III). This motion was evident in the channel maps of Paper III (Fig. 4) and in a velocity-position map (Fig 9), but in the line-integrated velocity map the streaming motion was less obvious. Nevertheless, in the velocity field image in Paper III (Fig. 11), the highest velocity gas occurs at the eastern end of the northern eyelid, displaced $`63^{}`$ in position angle from the kinematic minor axis rather than $`90^{}`$ for a normal rotation curve.
In Paper III we referred to the eastern tip of the oval, where the non-circular motion peaks, as the launch point for gas leaving the caustic and entering the tidal arm. The gas streams outward here mainly because it is traveling too fast for a circular orbit, but tidal forces also produce some outward acceleration. By the time the gas reaches the southern edge of the tidal arm, its peculiar motion has slowed down because of the gravity of the galaxy and large scale shocks.
During the formation of the tidal tail in IC 2163, NGC 2207 experienced a perpendicular tidal forcing from the gravity of IC 2163. This perturbation seems to have created an inward-propagating tidal warp with an overall spiral shape, as determined by shear and disk self-gravity (Paper IV). The corresponding warp in the velocity field was clearly seen (Paper III), although the optical disk shows little evidence for it. The warp is such that the side of NGC 2207 closest to IC 2163 in projection is warped toward it by about 9 kpc compared to the average inner disk plane.
## 4 Peculiar Dust Structures
Dust features are sensitive indicators of shock fronts because even small compressions can change the gas from optically thin to optically thick. There are many interesting dust structures in these galaxies as a result of the complex and unusual dynamics of the interstellar media. The high resolution and the backlighting of the outer spiral arm of NGC 2207 by IC 2163 provide a unique opportunity to study these structures.
Two peculiar dusty regions will be highlighted here: the parallel striations in the tidal tail of IC 2163, and the foreground spiral arms in NGC 2207. A dense dark cloud in the far western arm of NGC 2207 will be considered in the next section.
### 4.1 Parallel Dust Filaments in the Tidal Tail of IC 2163
The tidal tail in IC 2163 begins at the eastern end of the inner oval with a broad and shallow distribution of starlight. Gas and stars flow outward and downward, ending abruptly in the southeast along a curving arc where there is a dust lane, and extending northwest along the length of the arm for $`15`$ kpc. A dwarf galaxy with unknown velocity lies just off the tip of the tidal arm (the presence of globular clusters around this nucleated dwarf elliptical suggests it is not related to the interaction – see Paper V). Along the tidal outflow, there are numerous, nearly-parallel dust filaments perpendicular to the flow direction $`0.5^{\prime \prime }=85`$ pc thick, spaced by about $`2^{\prime \prime }`$. The filaments are darker than the surrounding regions by 0.2–0.8 mag in V band, suggesting that the midplane densities are enhanced by a factor of $`5`$ compared to the inter-filament gas (Paper V). The filament widths increase toward the outer edge of the tidal arm, where they take the form of two long, dense dust lanes, still parallel to each other. There is a slight tendency for the highest opacity dust streamers to occur at smaller radii (Paper V).
In our previous simulations considering only stars (Paper IV), the broad plateau of light on the inside part of the eastern tidal arm is a region of outward streaming stars. The outward motion is 150 km s<sup>-1</sup> in places and is the result of angular momentum conservation correcting for the inward plunge that these same stars recently made in response to the tidal force and self-gravity (Papers I and IV). The stars stream outward and form the tidal arm, where they slow down to a lower streaming speed of 50 km s<sup>-1</sup>. These model velocities for stars were chosen to match the observed velocities for the gas. Ground-based observations could barely resolve the dense, double dust lane at the edge of the arm, and they could not see the fainter filaments in the broad plateau.
Paper III briefly considered the possibility that the split in the tidal arm of IC 2163 was the result of dust, but we favored a different interpretation in which it resulted from two separate stellar arms created for a short time during the interaction. One of these arms was a caustic, reproduced by the model, and the other was a normal tidal arm. Now we see with higher resolution that the split is from dust and conclude either that there was no caustic arm or that the present epoch is not the correct one for viewing such a short-lived feature. The two dust lanes are hydrodynamic features, probably shock fronts, and could not be modeled before with our pure-stellar code. Thus we performed a new model with gas.
This new model is for a nearly-inplane, prograde interaction (25 inclination) between a gaseous disk galaxy and a point source. The tidal arm part of the model is shown in Figure 2. The ocular structure (i.e., the caustic oval in the inner disk) is produced again because this is a general property of prograde nearly-inplane encounters. Here we highlight the parallel filaments in the outward-streaming regions of the tidal arm, which resemble the filamentary dust structures in Figure 1. A more detailed model of the whole two-galaxy system, including a fuller analysis of the gas motions between the galaxies and the warp in NGC 2207, will be deferred to a later paper. The filamentary arms in the tidal tail of IC 2163 are present in the detailed study too.
In the present model, the gravitational potential of IC 2163 is dominated by a dark matter halo that is effectively rigid over the relatively short duration of this interaction. Smoothed particle hydrodynamics (Struck 1997; Kaufman et al. 1999) with 18,000 particles calculates the dynamics in the disk, which is assumed to be pure gas, using an adiabatic equation of state with heating and cooling. The initial disk was in rotational equilibrium with small thermal motions. It had an exponential density distribution over four scale lengths, and a thickness perpendicular to the plane equal to one scale length. The thermal terms lead to the formation of a multiphase initial gas. The companion is assumed to be a pure halo with half the scale length of the IC 2163 halo and 1.25 times the total IC 2163 mass. The observed ratio of H and K band luminosities is 1.6 (Kaufman et al. 1997).
Figure 2 shows three face-on views of the model disk and a projected view: (a: top left) is the map of model points at the beginning of the interaction, using color to represent initial variations in the particle positions with galactocentric radius; (b: top right) is the point map at a time representative of the current NGC 2207/IC 2163 system, with the same color scheme; (c: bottom left) shows the velocity vectors at the same time as in (b), and (d: bottom right) shows the density at the same time again, but rendered from a 160x160 array with color proportional to the density of the model points inside each pixel (blue is low density and red is high density), and projected by an inclination of $`30^{}`$ to resemble IC 2163 in the sky. The center of the companion, NGC 2207, is off the figure in b, c, and d.
The tidal tail in the model has the same shape as the observed tidal tail, and an intricate network of parallel filaments is in the broad part of this tail (seen best in the lower right panel). These filaments are stretched flocculent spiral arms that were present in the disk before perigalacticon (cf. Figure 2a). The model generates only flocculent spirals because of the relatively high mass of the halo and the low temperature of the pure-gas disk (see Elmegreen & Thomasson 1993).
The velocity vectors in Figure 2c indicate there is a shock front along the outer edge of the tidal arm, where the velocities abruptly change. This front also corresponds to a region of enhanced density, and it is located at the same place as the dense dust lane in the tidal arm of IC 2163. This result indicates that the dust lane running along the edge of the IC 2163 tidal arm is a shock front where the velocities change from outward to inward streaming. Not all of the tidal arm material is unbound from the galaxy: the inner portion returns to orbit the main disk after it shocks in the dust lane.
In Figure 1, many of the filaments in the tidal arm of IC 2163 have bright star clusters adjacent to them. These clusters are typically on the inside edges, toward smaller galactocentric radii. This is true for the small and faint dust filaments, as well as the dense dust lanes near the edge. The systematic displacements of the star clusters relative to the tidal arm dust filaments is an important check on the dynamics of the gas. Presumably these clusters formed in the dust lanes and then became ballistic after they formed. The dust and gas are not ballistic, however, and respond to the pressure forces. Thus the relative displacement of the dust and the clusters indicates the direction of the pressure gradient. To get the dust filaments systematically outside the clusters, the gas pressure must be decreasing with increasing galactocentric radius, as expected for an exponential disk.
The displacement between stars and dust in this interpretation is unlike the situation in an idealized spiral arm shock. There the displacement inside corotation (where the stars are systematically outside the dust lanes) is the result of a spiral wave and associated shock front (the dust lanes) that move slower in the azimuthal direction than both the gas and the stars. The dust lane in a normal spiral is newly shocked material that did not yet form stars, while the stars and the gas clouds in which they formed are both ballistic particles that move “downstream.” The dust lane material has just been shocked, and so is moving at its most negative radial speed. It is also at its largest galactocentric radius, having just streamed outward through the interarm region. As the dust lane gas moves inward and along the arm, it presumably forms clouds and stars that feel a Coriolis force deflecting them outward. In this way, their inward motion is converted into an azimuthal motion, and the young stars begin to emerge from the region of the shock toward the interarm region. For this normal spiral arm, the stars and the dust lane are displaced from each other because these two features are at different parts of their epicycles: the dust lanes are where the radial motion is greatest inward and the azimuthal motion relatively small, while the stars and associated clouds just downstream from the dust lanes are entering the parts of their epicycles where the azimuthal motion is greatest in the prograde sense. This juxtaposition creates the illusion that the stars formed out of gas that is located in the dust lane right next to them, but this is not the case. Instead, the young stars formed in gas clouds that are moving downstream with them, and the dust lane next to these stars is newly shocked material that has just arrived in the wave. The process of star formation in spiral arms will be discussed further in Section 4.3.
For the dust lanes in the tidal tail of IC 2163, the situation should be different. This is a stretched and distorted part of the former galaxy, and the motions through the arm are largely radial. In this case, the radial disk gradient of the interstellar pressure contributes to the force on the gas, pushing it faster than the stars which recently formed.
### 4.2 Comparison Between the Dust and Atomic Hydrogen in the Tidal Arm of IC 2163
Papers III and IV found H I streaming motions in the tidal arm of IC 2163 that are consistent with the outward flow of gas expected from the arm formation models. Figure 3 shows the H I overlaid on the B band HST image. The triangles correspond to the ridge of maximum streaming speed, where the line-of-sight velocity is fairly constant at 2970–2990 km s<sup>-1</sup> (the galaxy systemic velocity is $`2765\pm 20`$ km s<sup>-1</sup>). The pentagons correspond to a ridge of lower speed, ranging from 2780 to 2930 km s<sup>-1</sup> (Paper III, Sect. 4.2). The plus symbols are along the dividing line between the high velocity and the low velocity H I ridges, where the H I intensity is low.
In Paper III, we expected the dividing line to lie at a dust lane, where the high speed stream suddenly shocked into the low speed gas and the H I got converted into H<sub>2</sub>. This does not appear to be the case. Instead, the streaming ridge corresponds to a diffuse, optically faint region on the northern part of the tidal arm, and the low-speed gas in the south corresponds to the unresolved combination of the two main dust filaments. The dividing line (plus symbols in Fig. 3) is unrelated to the main dust lanes because it has a different curvature. The H I arm also has an S-shaped wiggle at $`\alpha _{1950}=6^h14^m25^s`$$`6^h14^m27^s`$, where the high and low velocity components of the arm appear to merge. The dust lanes do not show this structure.
The new model in Figure 2 explains the origin of these moving streams. Most of the high-velocity stream in the northern part of the tidal arm is from the original outer disk of IC 2163, which was H I-dominated and had a low stellar density before the encounter. The optical faintness of the northern part of the arm is consistent with an outer disk origin. The H I gas in this region moves quickly on its way to the tip of the arm where it may form a large gas pool (cf. Paper II). The intermediate zone of low H I emission (plus symbols) is mostly stellar, and it comes from the intermediate- to outer-radii of the optical disk of the former IC 2163. This part of the tidal tail moves in a southeasterly direction through the arm at a steady projected speed until the gas shocks along the curved arm edge. Then it forms the prominent parallel dust lanes at a relatively low line-of-sight velocity, corresponding to some of the gas moving inward.
The asymmetry in the fall-off of the H I intensity on the northern and southern sides of the tidal arm is consistent with this picture: the northern H I is extended gas that may have come all the way from the companion side of the IC 2163 outer disk, whereas the southern H I is mostly shocked gas at the leading front of the tidal arm. Figure 3b of the model in Paper II illustrates this morphology.
### 4.3 Foreground Dust in the NGC 2207 Spiral Arms
The background lighting by IC 2163 provides a unique view of the dust lane and star formation in an outer spiral arm of NGC 2207. What appears in a ground-based image to be a single dust lane in a spiral arm of NGC 2207 is seen at higher resolution to consist of 4 to 7 nearly parallel dust streamers that span the full width of the arm. The V-band magnitude decrements in many of the dust features of these arms are in the range from 0.4 to 1.5 mag (Paper V).
The spiral arm dust lanes in the HST image of NGC 2207 show an intricate internal structure reminiscent of Galactic cirrus clouds (Low et al. 1984) or other diffuse interstellar structures in the Milky Way, but on a much larger scale. There are also blue stars or clusters adjacent to many of the dust clouds, as if they just formed there. These spiral arms are density waves in NGC 2207, and their dust structures presumably delineate the shocks in these arms, as in the standard theory (Roberts 1969). However, the shocks are not smooth, and they are not just clumpy in the usual sense either (e.g. Combes & Gerin 1985; Roberts & Steward 1987; Elmegreen 1988). They are mostly composed of long, knotty filaments, which in some places run side by side in parallel streaks. Such structure suggests that the density wave shock occurred in several separate places along the width of the arm, or that the interarm gas is not in the form of spherical clouds, but filamentary like in M51 (Block et al. 1997).
We do not see the old stars in these arms because they are too faint compared to the background disk of IC 2163. We can see them in other parts of the foreground arms, however, such as the region north of IC 2163 before the arms cross in front of the disk. There are many faint red stars in these parts too, alongside young clusters and other dust lanes.
The blue stars in the foreground spiral arms are interesting because of what they tell us about the processes by which density waves trigger or organize star formation. This is the first example in an external galaxy where we can see such triggered star formation in projection against a background source. We see at this resolution (1 pixel $`=`$ 17 pc) blue stellar-like images mixed with the dust. These blue objects are probably young clusters (Paper V).
The star formation process looks normal at this perspective. The filaments contain clumps, and the clumps form stars. This is the same morphology as in local regions that are much smaller, such as Taurus. We cannot see any dust structure that might be indicative of cloud disruption, such as comet tails pointing away from the blue clusters, but the dust opacity may be too low for such detail on these scales.
The tidal forces acting on NGC 2207 are primarily perpendicular to its disk, and this galaxy is not flocculent, so it is unlikely that the parallel dust streamers in the spiral arm are from tidally stretched flocculent arms, as in IC 2163. The parallel dust filaments in NGC 2207 look similar to parallel dust filaments in the arms of M81. These M81 filaments are often at the upstream and downstream edges of the strings of giant HII regions in the arms, and may be parts of dense shells (Kaufman, Elmegreen, & Bash 1989), but it is also possible that M81 and other grand-design spirals have the same multi-stranded shock structures in their arms as NGC 2207.
Not all of the spiral arms in NGC 2207 have parallel dust structure. In the west, midway out in the optical disk, the dust lanes cut through the arm like spurs, but farther to the west, in the outer arm, some of the dust lanes are parallel again. Perhaps this change to spurs marks an orbital resonance.
## 5 The Dense Dark Cloud and Star-Forming Region in the Far West of NGC 2207
There is a peculiar dust region that appears to be associated with star formation in the far western part of the outer spiral arm of NGC 2207. It is the site of an intense radio continuum source, as shown in Figure 4, and the most luminous H$`\alpha `$ source in the system (see map in Paper V). An enlargement of the HST image is labeled feature i in Figure 6. In Paper III, we conjectured that this radio continuum source was a background radio galaxy, not associated with NGC 2207, because if it were at the distance of 35 Mpc its luminosity would be unlike anything else in either galaxy. Now, its coincidence with a dense dust cloud and bright young star cluster, and the similarity between the velocities of the associated H$`\alpha `$ emission and the H I emission from NGC 2207 in this vicinity, make the association of this radio continuum source with NGC 2207 more likely.
With a resolution of $`1.0^{\prime \prime }\times 1.9^{\prime \prime }`$ FWHM, Vila et al. (1990) detected a compact, radio continuum core with a deconvolved Gaussian size of $`1^{\prime \prime }`$ (marked with a plus sign in Fig. 4) at the location of the star cluster, surrounded by a more extended radio component. The core has flux densities of $`S(20)`$ =3.4 mJy at $`\lambda `$ 20 cm and $`S(6)=1.4`$ mJy at $`\lambda `$ 6 cm and a spectral index $`\alpha `$ of –0.7. Integrating over a $`7.5^{\prime \prime }`$ = 1.3 kpc region, they measured $`S(20)=10.3`$ mJy, $`S(6)=3.4`$ mJy and $`\alpha =0.9`$. The lower resolution radio continuum observations in Paper III found even more extended emission here, with $`S(20)`$ = 22 mJy and a deconvolved size of $`2.5\times 2.2`$ kpc (FWHM). Thus the core has a radio continuum luminosity 300 times that of Cas A; the core plus extended component forms a big, nonthermal radio source with a radio continuum luminosity $`1500`$ times that of Cas A if the spectral index of –0.9 applies throughout (taking the $`\lambda `$ 6 cm luminosity of Cas A as $`7\times 10^{17}`$ W Hz<sup>-1</sup> from Weiler, et al. 1989). Although comparable in radio luminosity to the most luminous radio supernovae (e.g. SN 1986J, SN 1988Z; Van Dyk, et al. 1993), the large linear size of the source in NGC 2207 suggests it is something else.
The H$`\alpha `$ source in this region has a luminosity (from Paper V) of $`1.4\times 10^{40}`$ erg s<sup>-1</sup> and a diameter of $`6^{\prime \prime }`$. Uncorrected for reddening, the H$`\alpha `$ flux is equivalent to S(6) = 0.13 mJy of optically thin free–free emission if $`T_e=10^4`$ K. This is much smaller than the measured $`\lambda `$ 6 cm flux density from approximately the same region. Its H$`\alpha `$ luminosity uncorrected for reddening is similar to that of 30 Dor (Kennicutt & Hodge 1986) and the brightest HII region in M51 (Van der Hulst et al. 1988). The strong nonthermal emission from the NGC 2207 source suggests that some of the H$`\alpha `$ flux may be from shocks or synchrotron sources rather than photoionization. A significant amount of H$`\alpha `$ could be occulted by the dark cloud, but not all of it because we still see star clusters, and the radio continuum and H$`\alpha `$ sources are larger than the cloud.
The structure of the region in the HST images is also peculiar because of a V-shaped feature, possibly a conical outflow, that has a bright star cluster at the apex and opens up to the north. This V-shape seems to extend on its western side to a bright star-like object in the NNW (which may not be related to it) at a distance of $`500`$ pc from the central cluster. A dense dust cloud trails off in the other direction. If this V-shape comes from a conical outflow, then there may be a peculiar and energetic star or collapsed object in this region, possibly with an accretion disk and jet. An x-ray survey might reveal a compact source.
To the southeast of the dense cloud and central star cluster, there are a number of smaller star clusters (Fig. 6i) that appear to be distributed in concentric arcs $`400`$ pc long. A similar pattern was found in an equally intense region of star formation in the galaxy NGC 6946 (Elmegreen, Efremov, & Larsen 2000), where several cluster arcs are inside and at the edge of a cleared $`600`$-pc cavity surrounding a $`15`$ My old globular cluster. If these arcs are not optical illusions (e.g. Bhavsar & Ling 1988), then the clusters could have been triggered in expanding partial shells. The large scale of this process, the unusual morphology, and the strong nonthermal radio emission from the NGC 2207 source suggest peculiar conditions. Considering the arc sizes, the explosions that triggered the clusters would have been extremely energetic, more than just single supernovae, and they were possibly one-sided since the apparent cluster arcs are not complete circles. It may be that a few extremely massive stars and their hypernovae (Paczyński 1998) are involved, in which case there could have been a gamma-ray burst from these regions several tens of millions of years ago (see also Efremov 2000). The one-sided nature of the cluster distribution could also be the result of an initial gas concentration on that side.
There is no significant excess of H I gas at the location of this dust cloud and star-forming region (see Figure 8 below). An H I cloud exists slightly to the south, along the spiral arm in a large darkish region, and another H I cloud is to the northeast, in the dark interarm region. These two clouds could be pieces of the envelope of a former giant molecular cloud centered on the dust feature. To the south along the same arm, the next big star-forming region does have a prominent H I cloud associated with it. Thus the more intense star-forming region associated with the dense dust cloud and the strong radio continuum source could have blown apart the low density parts of its cloud. We will discuss the other H I structure in these galaxies in more detail in Section 8.
The mass of the dust cloud can be estimated from its size ($`300`$ pc $`\times 140`$ pc) and apparent darkness. Paper V measured brightness deficits of $`\mathrm{\Delta }m_U=1.78\pm 0.3`$, $`\mathrm{\Delta }m_B=1.73\pm 0.3,`$ $`\mathrm{\Delta }m_V=1.24\pm 0.3,`$ and $`\mathrm{\Delta }m_I=1.38\pm 0.3`$ mag in the dark cloud. Since these values are all about equal, the dark cloud is optically thick in even the I band. If we assume an average visual extinction of 3 mag (twice $`\mathrm{\Delta }m_V`$ because of foreground stars), then the mass would be $`10^6`$ M.
## 6 Inner Disk Spirals in IC 2163
Paper III speculated that the nuclear region of IC 2163 might have a weak bar, or possibly be forming a bar now as the result of the interaction, as predicted by Noguchi (1987), and shown again by Gerin et al. (1990) and Paper I. The HST image shows no evidence for such a bar but there is a spiral arm system inside the oval that is elongated to make the overall structure resemble the inner Lindblad resonance (ILR) region of a bar.
The elongation of the two inner spirals arms is actually much more pronounced in the image than it is in reality because the galaxy is inclined with a line of nodes nearly parallel to the minor axis. This is the opposite of what most galaxies have: usually the projection line of nodes is along the major axis. For IC 2163, the kinematic minor axis is approximately the same as the morphological major axis in the oval region (Paper III). The kinematic minor axis from H I velocities has a position angle of $`155^{}`$, and the morphological major axis from optical and H I isophotes has a position angle of $`128^{}`$, which is only $`27^{}`$ different. The inclination is $`30^{}`$$`40^{}`$ from the model fits. Thus the bright oval in IC 2163 is actually more elongated than the image shows (cf. Fig. 2), by a factor of $`1.21.3`$, and the spirals inside the oval are more circular.
Figure 5 displays the radial behavior of the arm/interarm contrast of the inner–disk spirals in V and I-bands. The arm contrast is nearly constant with radius, which is unusual for the inner parts of density-wave arms (Elmegreen & Elmegreen 1995).
There is a curious dynamical property of this interaction that may point to the origin of the inner–disk spiral and even suggest a common mechanism for close interactions to trigger nuclear gas accretion and nuclear starbursts. This property follows from the relative position of the companion, NGC 2207, and the rotation curve inside IC 2163.
Figure 13 in Paper I showed the rotation curve of IC 2163, from the H I data. It is steeply rising as a solid body between the center and $`7^{\prime \prime }`$, and then it is more slowly rising in a steady fashion after that. This inner steep rise means that there will be an inner ILR somewhere inside $`7^{\prime \prime }`$ and an outer ILR somewhere outside $`7^{\prime \prime }`$ for any bi-symmetric spiral pattern that has a sufficiently low pattern speed. The semi-minor axis of the oval is about $`15^{\prime \prime }`$ and the inner radius of the inner spiral is about $`5^{\prime \prime }`$. Thus the inner spiral could extend between the inner and outer ILRs.
For a conventional bar potential with an ILR, there is a shock, or dust lane, parallel to the bar on the leading side and then a twist to a ring or inner spiral at around the ILR (Athanassoula 1992). If there is both an outer ILR and an inner ILR, then there will often be a 90 turn of this shock between the two resonances (Sanders & Huntley 1976). The spiral inside the oval of IC 2163 turns about 90, and may have its outer radius at the outer ILR, which would then be at the inner extent of the oval. The spiral could have its inner radius at the inner ILR. The whole oval would then be showing the response of the galaxy IC 2163 to a bar-like or $`\mathrm{cos}(2\theta )`$ potential: the long axis of the oval, which has an axial ratio of about 3:1 corrected for inclination, is along the “bar”, where the inward radial force and $`\mathrm{cos}(2\theta )`$ are maxima. The inner near-circular spiral, which extends from the inner minor axis of the “bar” to a smaller radius, twisting 90 along the way, would be the ILR shock feature. This bar-like structure presumably gave IC 2163 the de Vaucouleurs type of SB(rs)c pec, even though there is no bar in the classical sense.
An interesting thing about this system is that the tidal potential from NGC 2207 has the correct orientation to contribute to the “bar” potential in IC 2163, along with the oval itself. The tidal force from NGC 2207 pinches inward on IC 2163 in a direction that is perpendicular to the line connecting the two galaxies, and it pulls outward on IC 2163 in a direction that is parallel to this intergalactic line. This makes a $`\mathrm{cos}2\theta `$ forcing, much like in a bar (Combes 1988). For a real bar, the perturbation in the inward direction, relative to a circular potential, is along the bar, and the perturbation in the outward direction, relative to a circular potential, is on the minor axis of the bar. Thus the inward tidal pinch and outward tidal pull from NGC 2207 makes IC 2163 feel like it has a bar oriented in a direction perpendicular to the line connecting the two galaxies. This is in fact the mean orientation of the oval during the last $`80`$ My, which is the total time from the current position $`40`$ My after perigalacticon back to the symmetric position prior to perigalacticon. Recall that IC 2163 is quickly moving along the backside of NGC 2207 in a direction roughly from west to east and parallel to the disk of NGC 2207. The transient deformation of the disk into the oval accentuates this tidal force, increasing the bar-like potential, and this transient bar can drive ILR spirals.
What makes the tidally-induced bar potential so prominent in IC 2163 is the fortunate circumstance of catching the galaxy pair close to perigalacticon, and the extreme proximity of the encounter (perigalacticon was estimated to be $`2.3`$ times the radius of IC 2163 in Paper IV). Presumably the bar-like response in IC 2163 will last only during the closest approach, unless IC 2163 actually forms a permanent bar during this time, which is possible (Paper I).
The corotation position of the bar-like potential in this interpretation is the companion galaxy, NGC 2207. This is a large distance from IC 2163, and it gives a low pattern speed (approximately $`5.6`$ km s<sup>-1</sup> kpc<sup>-1</sup>). For such a pattern speed, the inner and outer ILRs should both exist for the rotation curve of IC 2163, and they should be sufficiently well separated to account for the inner-disk spiral structures. Unfortunately, the H I rotation curve is not sufficiently accurate, considering the streaming motions, to determine these inner and outer ILR positions any better.
A transient, tidally-induced, bar-like potential in an interacting galaxy should be able to drive mass inflow to the nucleus in much the same way as a permanent bar structure. The bar produces torques on the gas, and the gas moves to the center, primarily along the dust lanes and ILR spirals, with the resulting loss of angular momentum. If the encounter is too weak and fast to make a bar (the criteria for bar-making were discussed in Paper I), then the system will end up with a starburst in the nucleus of one or both galaxies (if they both felt an inplane tidal force), and there will be no evident bar structures at late times to indicate how the gas got to the center so quickly. In this way, tidal interactions may trigger gas accretion and nuclear starbursts in a variety of seemingly mild encounters.
The inner-disk spiral arms of IC 2163 may be an earlier form of the very bright inner-disk spiral arms of NGC 5394, discussed in Section 1. The model in Kaufman et al. (1999) finds that the latter developed from an ocular structure present at a slightly earlier time. NGC 5394 has a central starburst and its inner-disk spiral arms are unusual in that two of the three bright arms show no evidence for star formation.
## 7 Peculiar Emission Structures
We examined the HST image in detail for all peculiar emission features that might be conical or jet-like. Several interesting candidates were found. Most have a size in the range from 100 to 1000 pc.
### 7.1 Features on 100 to 1000 pc Scales
A collection of peculiar emission features on 100 to 1000 pc scales is found in Figure 6 and a finding chart is given in Figure 7. North is up for all the images in Figure 6 and the physical scale is on the right, assuming a galaxy distance of 35 Mpc. No feature appears in only one filter, which would make it an image flaw.
There are several linear or arc-shaped structures of blue knotty emission, such as c, h, s, t, and o. They could be composed of young stars except for their peculiar linear shapes. Feature f is similar to these but brighter; u is a curved line of star formation. Features a, g and p are other multiple-point structures, apparently made from young stars.
Jet-like features are seen in several places: e contains two jet-like structures that point toward each other; q has two jet-like features with a star-like object in the middle; b has a small jet-like object pointing away from a bright diffuse patch and other linear structures perpendicular to this; g (mentioned above) also contains a streak of faint emission between the central pointy structure and a star-like object in the north; n contains three faint linear emission streaks in what seems to be a region of star formation.
V-shaped features that could be conical emission regions include: i, the strong radio continuum source in the far west of NGC 2207, containing intense star formation, a dense dust cloud and a possible conical emission region in the north (see Sect. 5); l, a sideways V-shape, possibly a conical dust and reflection feature in the NGC 2207 spiral arm with blue star formation knots in it; and possibly r, a peculiar dark feature just to the upper right of the center, with a bright rim around it. Feature j contains two nearly parallel streaks running vertically in the figure, one streak is bright and the other is dark, giving it a three-dimensional quality.
There are two dust shells: k, a dust shell with three star formation knots in it, most likely in the NGC 2207 spiral because of the blue color of the star formation; and m, a dust shell with a red star in the middle (possibly a foreground star, see Paper V).
To the northwest of the dust shell k, there is a bright kpc-long arc (feature d) that looks like reflection off the inner edge of the dust lane that is inside the inner-disk arm.
The reddish features (k, d, l, m, n) juxtaposed on top of IC 2163 could be inside IC 2163 or in the foreground spiral arms of NGC 2207. The other features are mostly white or blue in color, and, except for a and possibly r, are positioned to be likely members of NGC 2207.
Individual jet-like objects are typically $`1^{\prime \prime }2^{\prime \prime }`$ or 150-300 pc long; the two streaks in feature e span 1000 pc. Analogous objects are not known in our Galaxy; much smaller versions of jets could be those associated with SS 443 (Margon 1984), the Crab Nebula jet (van den Bergh 1970; Gull & Fesen 1983), or Cas A (Fesen & Gunderson 1996). The jet-like regions in NGC 2207 might also be radio sources that are too weak and small to have been seen in our VLA survey. The origin of these features is not clear. Some may be related to mass transfer events in which gas from IC 2163 recently impacted the disk of NGC 2207, some may be optical illusions, and some may be normal galactic features not previously recognized in other galaxies.
### 7.2 Larger-Scale Linear Features
The top panel of Figure 6 shows four unusual elongated features that lie on approximately the same line. Each cuts across a spiral arm at a large angle. The most prominent is the straight ridge of star formation along the innermost northern spiral arm of NGC 2207 (near the middle of this panel at $`\alpha _{1950}=6^h14^m15.7^s`$, $`\delta _{1950}=21^{}21^{}5^{\prime \prime }`$). This ridge alone is not so peculiar because it has a dust lane on the inner edge and normal–looking young stars, but it is unusually straight for a density wave feature, and it does not follow the curvature of the rest of the arm. Just to the west of this bright feature (and along the same line) is a faint, thinner emission streak composed of several diffuse spots with no obvious dust lane. It cuts across the same spiral arm. In the opposite direction, in the next outer arm of NGC 2207 to the east of the brightest linear streak, there is another linear feature composed of dust in the center of two streaks of star formation; this is also not aligned with its local arm. The fourth feature is in the middle western arm of NGC 2207, on the right side of the top panel at $`\alpha _{1950}=6^h14^m10.7^s`$, $`\delta _{1950}=21^{}20^{}47^{\prime \prime }`$ (enlarged in panel u). The linear part crossing the spiral arm is made of dust and star formation, but just to the west of the arm, in the interarm region, is a large, diffuse “bubble-like” object. Several blue emission spots and a dust filament curve upward from the end of the bubble, and a faint, sideways V-shaped object that looks like a bow shock is further to the west, along the same line (cf. feature u in Fig. 6). The whole line of 4 features points back to within 1<sup>′′</sup> of the nucleus of IC 2163. This alignment resembles a jet, but there is no prominent radio continuum source in the center of IC 2163 nor along the line. Paper III found a strong radio continuum ridge $`15^{\prime \prime }`$ north of the line connecting these four linear features, and aligned at an angle of $`30^{}`$ relative to it.
The original WFPC2 chip seam on the mosaic runs underneath and almost parallel to this linear structure, beginning at the bright emission knot in the arm to the east of the main ridge and ending at the edge of the image on the west. However the individual features are visible in the separate fields before the mosaic was made, so these features are not artifacts of the seam.
Detailed numerical models of both galaxies in the interaction (Struck et al. 2000) occasionally find linear features from debris trails with material pulled out of IC 2163 and brushing the backside of NGC 2207. The lifetimes of any linear features composed of stars and dust are constrained by galactic rotation to be fairly short, less than several million years.
## 8 H I Emission and Star Formation
Maps of the H I emission from each galaxy are shown superposed on the HST images in Figures 3 and 8. Two peculiarities of the H I emission are its large velocity dispersion ($`50`$ km s<sup>-1</sup>) and large cloud masses ($`10^8`$ M; Paper III), which are characteristic of interacting systems (Paper II; Kaufman et al. 1997; 1999; Irwin 1994; Hibbard 1995).
The largest clouds in these galaxies are not clearly associated with star formation. Figure 14 in Paper III was a finding chart for $`10^8`$ M clouds that were also outlined on the contour diagrams in Figures 8 and 16 of that paper. The same contours are shown here in Figures 3 and 8, but without the outlines for giant clouds. Here the clouds in NGC 2207 are identified by the notation N1, N2, and so on.
In IC 2163, there are five giant concentrations of H I, two in the eastern tidal arm, one in the streaming region of the eastern tidal arm to the north, one in the western tidal arm near the center of NGC 2207, and one in the northern eyelid region (cf. Fig. 3). None are associated with specific star formation regions, but the one on the northern eyelid has star formation all along its inside border. The one in the western tidal arm is heavily obscured by NGC 2207. There is a less massive H I cloud in the southern eyelid region which is centered on a region with bright patches of star formation.
In NGC 2207 there is a giant H I cloud in the far northwest (N1) with only faint emission in the B band and three more in the western spiral arm below this (N2, N3, N4) with bright star formation only in the lowest one. The other two giant clouds in NGC 2207 (N5, N6) are in the east in Figure 8, in the region that is foreground to IC 2163. One (N5) is at the end of the outer eastern arm of NGC 2207 and is associated with a dust feature and faint blue cluster (cf. Fig. 1). The other (N6) is in a region crossed by dust lanes between the tidal bridge of IC 2163 and the middle arm on the eastern side of NGC 2207. The cloud is not centered on a bright star-forming region. It coincides with part of a 10 kpc long, radio continuum ridge (see Paper III).
The bright star-forming regions are mostly associated with smaller H I clouds, having masses of $`10^7`$ M or less. At this level of H I emission, many star-forming regions have some H I nearby, as is often the case in non-interacting galaxies. For example, in NGC 2207, the small patches of star formation in the southern arm are all associated with these less massive H I clouds, as is the small patch south of feature i (Fig. 6) in the western arm. The middle arm on the western side of NGC 2207 winds around to the north and then crosses in front of the central part of IC 2163. In the north, this arm has several bright regions of star formation along the H I ridge, some of which coincide with H I clumps. On the eastern side, the arm, seen as the dust streamers discussed in Section 4.3, also coincides with the H I ridge.
The giant star-forming region (feature i) that is associated with the intense radio continuum source on the outer western arm has no H I cloud, although there are smaller clouds to the south and west of it (N3, N2). The dense dust in this region suggests the hydrogen is present in molecular form.
There is no direct optical evidence for the large H I velocity dispersions that are present in these galaxies, which is typically 30–50 km s<sup>-1</sup> (Paper II). Such motions would be supersonic for H I. Large random motions also imply that the gas disks are thicker than normal, by perhaps a factor of $`5`$. The models in Paper IV predict a warp in NGC 2207 that is pointed away from us on the IC 2163-side and toward us on the western side. Paper III found that on the western side of NGC 2207, the optical radial surface brightness profile has a plateau from $`40^{\prime \prime }90^{\prime \prime }`$, which corresponds to the broad ring of H I emission that peaks at $`70^{\prime \prime }`$. The increase in the line-of-sight thickness produced by the warp may be responsible for the prominence of the H I ring on the eastern and western sides and for the unusual behavior of the radial surface brightness profile in the optical on the western side. Since most of the massive clouds in NGC 2207 are in these directions, their surface mass densities would be slightly smaller than computed by neglecting the warp.
The high H I velocity dispersion could affect the optical observations in other ways too. The dispersion could broaden the gas response to the spiral arms, for example, or remove the tendency for the gas to shock once in a thin dust lane. Instead, it could be shocking several times in the parallel dust lanes that are observed in front of IC 2163 (see Sect. 4.3). The H I line profiles are typically not Gaussian, so they could be composed of a small number of streams with different velocities. In that case, the high velocity dispersion could result from variations in the systematic motions within the H I synthesized beam. These systematic motions could lead to the multiple shocks seen as dust lanes. Because the H I observations have a large beam ($`13.5^{\prime \prime }\times 12^{\prime \prime }`$), any streaming motions on a kpc scale would be unresolved and appear only as turbulence.
Strong turbulence or high speed streaming motions could also make the arms in NGC 2207 thicker than those in other spiral galaxies, such as M100 or M81. The NGC 2207 arms are similar to those in Luminosity Class III galaxies, which are usually much smaller than NGC 2207 (Iye & Kodaria 1976). This makes sense if the relative thickness of the arms scales with the ratio of the velocity dispersion to the orbit speed.
The origin of the high HI velocity dispersion is not known. Other strongly interacting galaxies have these motions too, so the turbulence could come from internal adjustments to tidal forces. We show in Paper V that the regions of highest dispersion contain superstar clusters, so there could be a connection with star formation too.
## 9 Conclusions
The galaxies IC 2163 and NGC 2207 are involved in a near encounter that has compressed IC 2163 in the plane and warped NGC 2207 out of the plane. Optical observations with HST show many interesting and peculiar features. Two types of extinction patterns were discussed in Section 4: long parallel filaments in the tidal tail of IC 2163, and clumpy filaments in the spiral arms of NGC 2207 that are seen in projection against IC 2163.
The tidal-tail filaments seem to be normal flocculent spiral arms that were in the disk of IC 2163 before the encounter and then got stretched out into the tidal tail after the encounter by the overall disk deformation. A numerical simulation reproduced these filaments well. The extreme youth of the interaction ($`40`$ My) explains why the spiral arms were not individually affected much, except for the overall distortion that followed the tidal flow. Many of the filaments have star formation that trails systematically behind them, suggesting a large-scale acceleration of the gas relative to the stars that form inside it.
The foreground spiral arm in NGC 2207 shows multiple parallel filaments as well, but probably for a different reason. There could be multiple shocks in the density wave, or independent filaments that came in from the interarm region. Such structure could be normal for galaxies but not commonly observed because of the rarity of background lighting. In any case, star formation occurs in the clumps of these filaments in a way that resembles local star formation in small filamentary dark clouds. The process of star formation in the density waves of NGC 2207 is therefore one in which filaments, either made by the wave or pre-existing, collapse gravitationally into globules which then continue to collapse into clusters and individual stars. The star formation does not look like it is occurring at the interfaces between colliding clouds or colliding filaments.
Another interesting region is a dense dark cloud in a region of star formation on the outer western arm of NGC 2207. This is the site of the most luminous H$`\alpha `$ source in the galaxies, and a large, nonthermal radio continuum source 1500 times more luminous that Cas A. There is a massive cluster in the center, a conical feature like an outflow to the north, and other clusters in the southeast that could have been triggered by energetic explosions.
The spirals inside the oval of IC 2163 were discussed in Section 6. The dynamical time of this inner region is much shorter than the interaction time, so these spirals could be part of the response. We suggested, on the basis of the rotation curve, that the central spirals could extend from the outer ILR to the inner ILR for a pattern speed that places corotation at the companion. In that case, the whole oval may be viewed as a bar-like stellar+gaseous flow pattern in the transient $`\mathrm{cos}(2\theta )`$ potential of the tidal field. The inner spirals are then a normal resonance response for such a bar. This circumstance highlights the interesting possibility that transient tidal forces may trigger nuclear gas inflow via bar-like hydrodynamics even when there is no bar. It also suggests that post-encounter starburst systems without obvious bars today may have had their major accretion events close to perigalacticon when the orbits were highly distorted.
Numerous peculiar emission features were found after a close examination of the HST image. Many are linear with either smooth or clumpy internal structures, and some have star-like objects at one or both ends. They are typically several hundred parsecs long, although smaller features would be difficult to see. Some could be jets or conical outflows, but they are much larger than protostellar jets found in the Milky Way. They may be coincidental alignments of stars or clusters, but the smoothest ones do not look like this. A long line of four co-linear features that are at odd angles to their local spiral arms was also noted.
The H I is weakly associated with star formation in these galaxies, but the largest H I clouds, with masses of around $`10^8`$ M, are not generally producing rich clusters. The H I velocity dispersion is $`5\times `$ higher than normal, which may explain why the spiral arms in NGC 2207 look thick and feathery.
The high-resolution images obtained by the Hubble Space Telescope have presented us with an unprecedented opportunity to study the morphology of a pair of interacting galaxies in the early stages of interaction. It has shown us detailed features not previously observed elsewhere, and given us important insights into the early stages of galaxy interactions.
Acknowledgements:
This work was supported by NASA through grant number GO-06483-95A from the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5-26555. Support for the Hubble Heritage Team’s contribution also came from STScI.
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# Permutations avoiding a pattern from 𝑆_𝑘 and at least two patterns from 𝑆₃
## 1 Introduction
Let $`[k]=\{1,\mathrm{},k\}`$ be a (totally ordered) alphabet on $`k`$ letters, and let $`\alpha [k]^m`$, $`\beta [l]^m`$ with $`lk`$. We say that $`\alpha `$ is order-isomorphic to $`\beta `$ if the following condition holds for all $`1i<jn`$: $`\alpha _i<\alpha _j`$ if and only if $`\beta _i<\beta _j`$.
We say that $`\tau S_n`$ contains $`\alpha S_k`$ if there exist $`1i_1<\mathrm{}<i_kn`$ such that $`(\tau _{i_1},\mathrm{},\tau _{i_k})`$ is order-isomorphic to $`\alpha =(\alpha _1,\mathrm{},\alpha _k)`$, and we say that $`\tau `$ avoids $`\alpha `$ if $`\tau `$ does not contain $`\alpha `$. The set of all permutations in $`S_n`$ avoiding $`\alpha `$ is denoted $`S_n(\alpha )`$. More generally, for any finite set of permutations $`T`$ we write $`S_n(T)`$ to denote the set of permutations in $`S_n`$ avoiding all the permutations in $`T`$. Two sets, $`T_1,T_2`$, are said to be Wilf equivalent (or to belong to the same Wilf class) if and only if $`|S_n(T_1)|=|S_n(T_2)|`$ for any $`n0`$; the Wilf class of $`T`$ we denote by $`\overline{T}`$.
The study of the sets $`S_n(\alpha )`$ was initiated by Knuth , who proved that $`|S_n(\alpha )|=\frac{1}{n+1}\left(\genfrac{}{}{0pt}{}{2n}{n}\right)`$ for any $`\alpha S_3`$. Knuth’s results where further extended in two directions. West and Stankova analyzed $`S_n(\alpha )`$ for $`\alpha S_4`$ and obtained the complete classification, which contains $`3`$ distinct Wilf classes. This classification, however, does not give exact values of $`S_n(\alpha )`$. On the other hand, Simion and Schmidt studied $`S_n(T)`$ for arbitrary subsets $`TS_3`$ and discovered $`7`$ Wilf classes. The study of $`S_n(\alpha ,\tau )`$ for all $`\alpha S_3`$, $`\tau S_4(\alpha )`$, was completed by West , Billey, Jockusch and Stanley , and Guibert .
In the present paper, we calculate the cardinalities of the sets $`S_n(T,\tau )`$ for all $`TS_3`$, $`|T|2`$, and all permutations $`\tau S_k`$ such that $`k3`$.
###### Remark 1
West observed that if $`\tau `$ contains a pattern in $`T`$, then $`|S_n(T,\tau )|=|S_n(T)|`$. Therefore, in what follows we assume that $`\tau S_k(T)`$.
Throughout the paper, we often make use of the following simple statement.
###### Lemma 1
Let $`\{s_i(x)\}_{i=1}^r`$, $`\{A_i(x)\}_{i=1}^r`$ and $`\{B_i(x)\}_{i=1}^r`$ be sequences of functions such that
$$s_i(x)=A_i(x)s_{i+1}(x)+B_i(x),$$
where $`1ir1`$, and $`s_r(x)=h(x)`$. Then
$$s_1(x)=\left|\begin{array}{ccccc}B_1(x)& A_1(x)& 0& \mathrm{}& 0\\ B_2(x)& 1& A_2(x)& \mathrm{}& 0\\ B_3(x)& 0& 1& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& 0& \mathrm{}& 0\\ B_{r1}(x)& 0& 0& \mathrm{}& A_{r1}(x)\\ h(x)& 0& 0& \mathrm{}& 1\end{array}\right|.$$
###### Proof
Immediately, by definitions and induction on $`r`$.
Our calculation is divided into three sections corresponding to the cases $`|T|=2,3`$ and $`|T|4`$. In the last section, as an example, we will obtain the complete classification (Table 1) of all cardinalities of the sets $`S_n(T,\tau )`$ where $`TS_3`$, $`\tau S_4`$.
## 2 Avoiding a pair from $`S_3`$ and a pattern from $`S_k`$
In this section, we calculate the cardinality of the sets $`S_n(\beta ^1,\beta ^2,\tau )`$ where $`\beta ^1,\beta ^2S_3`$, $`\tau S_k`$, $`k3`$. By Remark 1 and by the three natural operations, the complementation, the reversal and the inverse (see Simion and Schmidt , Lemma $`1`$), we have to consider the following four possibilities:
$`\begin{array}{cccc}\hfill 1)& S_n(123,132,\tau ),\hfill & \text{where}\hfill & \tau S_k(123,132),\hfill \\ \hfill 2)& S_n(123,231,\tau ),\hfill & \text{where}\hfill & \tau S_k(123,231),\hfill \\ \hfill 3)& S_n(132,213,\tau ),\hfill & \text{where}\hfill & \tau S_k(132,213),\hfill \\ \hfill 4)& S_n(213,231,\tau ),\hfill & \text{where}\hfill & \tau S_k(213,231).\hfill \end{array}`$
The main body of this section is divided into four subsections corresponding to the above four cases.
### 2.1 $`T=\{123,132\}`$.
Let $`a_\tau (n)=|S_n(123,132,\tau )|`$, and let $`a_\tau (x)`$ be the generating function of the sequence $`a_\tau (n)`$, that is,
$$a_\tau (x)=\underset{n0}{}a_\tau (n)x^n.$$
We find an explicit expression for the generating function $`a_\tau (x)`$.
###### Theorem 2.1
Let $`\tau S_k(123,132)`$. Then:
1. there exist $`r_1,r_2,\mathrm{},r_m1`$ with $`r_1+r_2+\mathrm{}+r_m=k`$ such that
$$\tau =(\beta _1,\beta _2,\mathrm{},\beta _m),$$
where $`\beta _i=(t_i1,t_i2,\mathrm{},t_ir_i+1,t_i)`$, and $`t_i=k(r_1+\mathrm{}+r_{i1})`$ for $`i=1,2,\mathrm{},m`$;
2. $$a_\tau (x)=\left|\begin{array}{ccccc}f_{r_1}(x)& g_{r_1}(x)& 0& \mathrm{}& 0\\ f_{r_2}(x)& 1& g_{r_2}(x)& \mathrm{}& 0\\ f_{r_3}(x)& 0& 1& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& g_{r_{m1}}(x)\\ f_{r_m}(x)& 0& 0& 0& 1\end{array}\right|,$$
where $`f_d(x)=\frac{1x}{12x+x^d}`$ and $`g_d(x)=\frac{x^d}{12x+x^d}`$.
###### Proof
$`(i)`$ Let $`\tau S_k(123,132)`$; choose $`r_1`$ such that $`\tau _{r_1}=k`$. Since $`\tau `$ avoids $`132`$, we see that $`\tau _jkr_1+1`$ for all $`jr_1`$, and since $`\tau `$ avoids $`123`$, we get that $`\tau =(\beta _1,\tau ^{})`$, where $`\tau ^{}S_{kr_1}(123,132)`$, and so on.
$`(ii)`$ Let $`\tau =(\beta _1,\mathrm{},\beta _m)`$, and let $`\alpha S_n(123,132,\tau )`$; choose $`t`$ such that $`\alpha _t=n`$. Similarly to $`(i)`$, $`\alpha =(n1,\mathrm{},nt+1,n,\alpha _{t+1},\mathrm{},\alpha _n)`$, therefore
$$a_\tau (n)=\underset{t=1}{\overset{r_11}{}}a_\tau (nt)+\underset{t=r_1}{\overset{n}{}}a_\tau ^{}(nt),$$
which means
$$a_\tau (n)=2a_\tau (n1)a_\tau (nr_1)+a_\tau ^{}(nr_1)$$
for all $`nk+1`$, where $`\tau ^{}=(\beta _2,\mathrm{},\beta _m)`$. Hence
$$a_{(\beta _i,\mathrm{},\beta _m)}(n)=2a_{(\beta _i,\mathrm{},\beta _m)}(n1)a_{(\beta _i,\mathrm{},\beta _m)}(nr_i)+a_{(\beta _{i+1},\mathrm{},\beta _m)}(nr_i)$$
for all $`nt_i+1`$, or equivalently,
$$\begin{array}{cc}\underset{nt_i+1}{}a_{(\beta _i,\mathrm{},\beta _m)}(n)x^n\hfill & =2x\underset{nt_i}{}a_{(\beta _i,\mathrm{},\beta _m)}(n)x^n\hfill \\ & x^{r_i}\underset{nt_i+1}{}a_{(\beta _i,\mathrm{},\beta _m)}(n)x^n\hfill \\ & +x^{r_i}\underset{nt_i+1}{}a_{(\beta _{i+1},\mathrm{},\beta _m)}(n)x^n.\hfill \end{array}$$
Since $`a_{(\beta _i,\mathrm{},\beta _m)}(n)=2^{n1}`$ for all $`nt_i1`$, $`a_{(\beta +i,\mathrm{},\beta _m)}(t_i)=2^{t_i1}1`$, and $`a_{(\beta _i,\mathrm{},\beta _m)}(0)=1`$ (Simion and Schmidt , Proposition $`7`$), we obtain that
$$(12x+x^{r_i})a_{(\beta _i,\mathrm{},\beta _m)}(x)=1x+x^{r_i}a_{(\beta _{i+1},\mathrm{},\beta _m)}(x)$$
for all $`im1`$, and
$$(12x+x^{r_m})a_{(\beta _m)}(x)=1x.$$
Hence, by Lemma 1, the theorem holds.
###### Example 1
Let $`T=\{123,132\}`$. By Theorem 2.1,
1. $`|S_n(T,3214)|=t_n`$, where $`t_n`$ is the $`n`$-th Tribonacci number , and $`|S_n(T,3241)|=f_{n+2}1`$, where $`f_n`$ is the $`n`$-th Fibonacci number.
2. $`|S_n(T,3412)|=|S_n(T,4231)|=\left(\genfrac{}{}{0pt}{}{n}{2}\right)+1`$.
3. $`|S_n(T,3421)|=3n5`$.
### 2.2 $`T=\{123,231\}`$.
In this subsection, we calculate the cardinality of the set $`S_n(123,231,\tau )`$, where $`\tau S_k(123,231)`$. This cardinality we denote by $`b_\tau (n)`$.
###### Lemma 2
Let $`\tau S_k(123,231)`$. Then, either there exists $`r`$, $`1rk1`$, such that $`\tau =(r,\mathrm{},2,1,k,k1,\mathrm{},r+1)`$, and hence
$$b_\tau (n)=(k2)n\frac{k(k3)}{2}\text{for all}nk,$$
or $`\tau =(k,\tau ^{})(k,\mathrm{},2,1)`$ such that $`\tau ^{}S_{k1}(123,231)`$, and hence
$$b_\tau (n)=b_\tau ^{}(n1)+n1\text{for all}nk.$$
###### Proof
Let $`\tau S_k(123,231)`$; put $`r=\tau _1`$. Since $`\tau `$ avoids $`123`$, we see that $`\tau `$ contains $`(r,k,\mathrm{},r+1)`$, since $`\tau `$ avoids $`231`$, we get that $`\tau =(r,\tau ^{},k,k1,\mathrm{},r+1)`$, and since $`\tau `$ avoids $`123`$, we have two cases: either $`\tau =(r,\mathrm{},1,k,\mathrm{},r+1)`$ for $`1rk1`$, or $`\tau =(k,\tau ^{})`$ such that $`\tau ^{}S_{k1}(123,231)`$. Now let us consider the two cases:
1. Let $`\alpha S_n(123,231,\tau )`$, where $`\tau =(r,\mathrm{},1,k,\mathrm{},r+1)`$, $`1rk1`$. Similarly to the above, we have two cases for $`\alpha `$: in the first case $`\alpha =(t,\mathrm{},1,n,n1,\mathrm{},t+1)`$ for $`1tn1`$, so there are $`k2`$ permutations like $`\alpha `$. In the second case $`\alpha =(n,\alpha _2,\mathrm{},\alpha _n)`$, so there are $`b_\tau (n1)`$ permutations, which means $`b_\tau (n)=b_\tau (n1)+k2`$. Besides, $`a_\tau (k)=k(k1)/2`$ (see Simion and Schmidt , Proposition $`11`$), hence $`b_\tau (n)=(k2)n\frac{k(k3)}{2}`$.
2. Let $`\alpha S_n(123,231,\tau )`$, and let $`\tau =(k,\tau ^{})(k,\mathrm{},1)`$ such that $`\tau ^{}S_{k1}(123,231)`$. Similarly to the above, we have two cases for $`\alpha `$: in the first case $`\alpha =(t,\mathrm{},1,n,n1,\mathrm{},t+1)`$ for $`1tn1`$, so there are $`n1`$ permutations like $`\alpha `$. In the second case $`\alpha =(n,\alpha _2,\mathrm{},\alpha _n)`$, so there are $`b_\tau ^{}(n1)`$ permutations. Hence $`b_\tau (n)=b_\tau ^{}(n1)+n1`$.
###### Theorem 2.2
Let $`\tau S_k(123,231)`$. Then :
1. there exist $`m`$, $`2mk+1`$, and $`r`$, $`1rm2`$, such that
$$\tau =(k,\mathrm{},m,r,\mathrm{},1,m1,\mathrm{},r+1);$$
2. for all $`nk`$
$$b_\tau (n)=(k2)n\frac{k(k3)}{2}.$$
###### Proof
$`(i)`$ Immediately, by Lemma 2,
$$\tau =(k,\mathrm{},m,r,\mathrm{},1,m1,\mathrm{},r+1),$$
where $`2mk+1`$, $`1rm2`$.
$`(ii)`$ Again, by Lemma 2, for all $`nk`$
$$b_\tau (n)=\underset{j=1}{\overset{km+1}{}}(nj)+(m3)(n(km+1))+\frac{(m1)(4m)}{2},$$
hence, this theorem holds.
###### Example 2
Let $`T=\{123,231\}`$; by the Theorem 2.2
$$|S_n(T,4312)|=|S_n(T,1432)|=|S_n(T,2143)|=|S_n(T,3214)|=2n2.$$
### 2.3 $`T=\{132,213\}`$.
Let $`c_\tau (n)=|S_n(132,213,\tau )|`$, and let $`c_\tau (x)`$ be the generating function of the sequence $`c_\tau (n)`$. We find an explicit expression for the generating function $`c_\tau (x)`$.
###### Theorem 2.3
Let $`\tau S_k(132,213)`$. Then:
1. there exist $`k+1=r_0>r_1>\mathrm{}>r_m1`$ such that
$$\tau =(r_1,r_1+1,\mathrm{},k,r_2,r_2+1,\mathrm{},r_11,\mathrm{},r_m,r_m+1,\mathrm{},r_{m1}1);$$
2. $$c_\tau (x)=\left|\begin{array}{ccccc}f_{r_0r_1}(x)& g_{r_0r_1}(x)& 0& \mathrm{}& 0\\ f_{r_1r_2}(x)& 1& g_{r_1r_2}(x)& \mathrm{}& 0\\ f_{r_2r_3}(x)& 0& 1& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& g_{r_{m2}r_{m1}}(x)\\ f_{r_{m1}r_m}(x)& 0& 0& 0& 1\end{array}\right|,$$
where $`f_d(x)`$ and $`g_d(x)`$ have the same meaning as in Theorem 2.1.
###### Proof
$`(i)`$ Let $`\tau S_k(132,213)`$, and let $`r_1=\tau _1`$. Since $`\tau `$ avoids $`132`$, we see that $`\tau `$ contains $`(r_1,r_1+1,\mathrm{},k)`$, and since $`\tau `$ avoids $`213`$, we get that $`\tau =(r_1,r_1+1,\mathrm{},k,\tau ^{})`$, where $`\tau ^{}S_{r_11}(132,213)`$, and so on.
$`(ii)`$ Let $`\alpha S_n(132,213,\tau )`$, and let $`t=\alpha _1`$; similarly to $`(i)`$, $`\alpha =(t,t+1,\mathrm{},n,\alpha _{nt+2},\mathrm{},\alpha _n)`$. Therefore, for $`tnk+r_1`$ we have $`\alpha S_n(132,213,\tau )`$ if and only if $`(\alpha _{nt+2},\mathrm{},\alpha _n)S_{t1}(132,213,\tau ^{})`$, and for $`tnk+r_1+1`$ we have $`\alpha S_n(132,213,\tau )`$ if and only if $`(\alpha _{nt+2},\mathrm{},\alpha _n)S_{t1}(132,213,\tau )`$. Hence
$$c_\tau (n)=\underset{t=1}{\overset{nk+r_1}{}}c_\tau ^{}(t1)+\underset{t=nk+r_1+1}{\overset{n}{}}c_\tau (t1),$$
which means that
$$c_\tau (n)=2c_\tau (n1)c_\tau (nk+r_11)+c_\tau ^{}(nk+r_11).$$
Let us define $`t_i=r_{i1}r_i`$ for $`i=1,2,\mathrm{},m`$, so
$$c_\tau (n)=2c_\tau (n1)c_\tau (nt_1)+c_\tau ^{}(nt_1).$$
If $`\tau =(1,2,\mathrm{},k)`$, then immediately $`c_\tau (x)=\frac{1x}{12x+x^k}`$, hence, by Lemma 1, this theorem holds.
###### Corollary 1
Let $`k2`$. For all $`n0`$,
$$c_{(k,\mathrm{},2,1)}(n)=\underset{j=0}{\overset{k2}{}}\left(\genfrac{}{}{0pt}{}{n1}{j}\right).$$
###### Proof
By the proof of Theorem 2.3,
$$c_{(k,\mathrm{},2,1)}(n)=2^k+\underset{t=k1}{\overset{n1}{}}c_{(k1,\mathrm{},2,1)}(t),$$
which means that
$$c_{(k,\mathrm{},2,1)}(n)=c_{(k,\mathrm{},2,1)}(n1)+c_{(k1,\mathrm{},2,1)}(n1).$$
Besides, $`c_{(k,\mathrm{},2,1)}(k)=2^{k1}1`$, and $`c_{(k,\mathrm{},2,1)}(n)=2^{n1}`$ for $`1nk1`$ (see Simion and Schmidt , Proposition $`8`$). Hence, the corollary is true.
###### Example 3
Let $`T=\{132,213\}`$.
1. By Corollary 1, $`|S_n(T,4321)|=\left(\genfrac{}{}{0pt}{}{n}{2}\right)+1`$.
2. By Theorem 2.3, $`|S_n(T,1234)|=t_n`$ where $`t_n`$ is the $`n`$-th Tribonacci number.
3. By Theorem 2.3, $`|S_n(T,2341)|=f_{n+2}1`$ where $`f_n`$ is the $`n`$-th Fibonacci number.
4. By Theorem 2.3, $`|S_n(T,3412)|=|S_n(T,3421)|=|S_n(T,4231)|=\left(\genfrac{}{}{0pt}{}{n}{2}\right)+1`$.
### 2.4 $`T=\{213,231\}`$.
Let $`d_\tau (n)=|S_n(213,231,\tau )|`$, and let $`d_\tau (x)`$ be the generating function of the sequence $`d_\tau (n)`$. We find an explicit expression for the generating function $`d_\tau (x)`$.
###### Theorem 2.4
Let $`\tau S_k(213,231)`$. Then:
1. $`\tau _i`$ is either the right maximum, or the righr minimum, for all $`1ik1`$;
2. $$d_\tau (x)=\left|\begin{array}{ccccc}1& g(x)& 0& \mathrm{}& 0\\ 1& 1& g(x)& \mathrm{}& 0\\ 1& 0& 1& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& 0\\ 1& 0& 0& 0& g(x)\\ 1& 0& 0& 0& 1\end{array}\right|,$$
where $`g(x)=\frac{x}{1x}`$.
###### Proof
$`(i)`$ Let $`\tau S_k(213,231)`$; if $`2\tau _1k1`$, then $`\tau `$ contains either $`(\tau _1,1,k)`$ or $`(\tau _1,k,1)`$, which means $`\tau `$ contains either $`213`$ or $`231`$, hence $`\tau _1=1`$ or $`\tau _1=k`$, and so on.
$`(ii)`$ Let $`\alpha S_n(213,231,\tau )`$; similarly to $`(i)`$, $`\alpha _1=1`$ or $`\alpha _1=n`$. Let $`\tau =(\tau _1,\tau ^{})`$, hence in the above two cases ($`\tau _1=1`$ or $`\tau _1=k`$) we obtain
$$d_\tau (n)=d_\tau (n1)+d_\tau ^{}(n1)$$
for all $`nk`$. Besides, $`d_\tau (0)=1`$, $`d_\tau (k)=2^{k1}1`$, and $`d_\tau (n)=2^{n1}`$ for all $`1nk1`$ (see Simion and Schmidt , Proposition $`10`$). Hence, similarly to Theorem 2.1, the theorem holds.
Immediately by Theorem 2.4,
$$|S_n(123,132,231)|=|S_n(132,231,321)|=n,$$
for all $`n0`$, which means, we have a generalization Proposition $`16`$ and Lemma $`6(b)`$ of Simion and Schmidt .
###### Example 4
Let $`T=\{213,231\}`$. By Theorem 2.4,
$$|S_n(T,1234)|=|S_n(T,1243)|=|S_n(T,1423)|=|S_n(T,1432)|=\left(\genfrac{}{}{0pt}{}{n}{2}\right)+1.$$
## 3 Three patterns from $`S_3`$ and a pattern from $`S_k`$
In this section, we calculate the cardinality of the sets $`S_n(T,\tau )`$ such that $`TS_3`$, $`|T|=3`$ and $`\tau S_k(T)`$ for $`k3`$. By Remark 1 and by three natural operations the complementation, the reversal and the inverse (see Simion and Schmidt , Lemma $`1`$), we have to consider the following five possibilities:
$`\begin{array}{cccc}\hfill 1)& S_n(123,132,213,\tau ),\hfill & \text{where}\hfill & \tau S_k(123,132,213),\hfill \\ \hfill 2)& S_n(123,132,231,\tau ),\hfill & \text{where}\hfill & \tau S_k(123,132,231),\hfill \\ \hfill 3)& S_n(123,213,231,\tau ),\hfill & \text{where}\hfill & \tau S_k(123,213,231),\hfill \\ \hfill 4)& S_n(123,231,312,\tau ),\hfill & \text{where}\hfill & \tau S_k(123,231,312),\hfill \\ \hfill 5)& S_n(132,213,231,\tau ),\hfill & \text{where}\hfill & \tau S_k(132,213,231).\hfill \end{array}`$
###### Remark 2
By Erdös and Szekeres , $`|S_n((1,2,\mathrm{},a),(b,b1,\mathrm{},1))|=0`$ for all $`n(a1)(b1)+1`$, where $`a,b1`$. Therefore, in what follows we assume that $`\tau S_k(T)`$ and $`\tau (k,k1,\mathrm{},1)`$, since $`123T`$.
The main body of this section is divided into five subsections corresponding to the above five cases.
### 3.1 $`T=\{123,132,213\}`$.
Let $`e_\tau (x)`$ be the generating function of the sequence $`|S_n(T,\tau )|`$. We find an explicit expression for the generating function $`e_\tau (x)`$.
###### Lemma 3
Let $`T=\{123,132,213\}`$ and $`\tau S_k(T)`$. Then, either there exists $`\tau ^{}S_{k1}(T)`$ such that $`\tau =(k,\tau ^{})(k,k1,\mathrm{},1)`$, and hence
$$|S_n(T,\tau )|=|S_{n1}(T,\tau ^{})|+|S_{n2}(T,\tau ^{})|\text{for any}nk,$$
or there exists $`\tau ^{\prime \prime }S_{k2}(T)`$ such that $`\tau =(k1,k,\tau ^{})`$, and hence
$$|S_n(T,\tau )|=|S_{n1}(T,\tau )|+|S_{n2}(T,\tau ^{\prime \prime })|\text{for any}nk.$$
###### Proof
Let $`\tau S_k(T)`$; since $`\tau `$ avoids $`123`$ and $`132`$ we have either $`\tau _1=k`$ or $`\tau _1=k1`$. If $`\tau _1=k1`$, then, since $`\tau `$ avoids $`213`$, we see that $`\tau =(k1,k,\tau ^{\prime \prime })`$. Now we consider the two cases:
1. Let $`\tau =(k,\tau ^{})`$, $`\alpha S_n(T,\tau )`$. Similarly to the above, either $`\alpha =(n,\alpha _2,\mathrm{},\alpha _n)`$, or $`\alpha =(n1,n,\alpha _3,\mathrm{},\alpha _n)`$, so evidently
$$|S_n(T,\tau )|=|S_{n1}(T,\tau ^{})|+|S_{n2}(T,\tau ^{})|.$$
2. Let $`\tau =(k1,k,\tau ^{\prime \prime })`$, $`\alpha S_n(T,\tau )`$. Similarly to the above, either $`\alpha =(n,\alpha _2,\mathrm{},\alpha _n)`$, or $`\alpha =(n1,n,\alpha _3,\mathrm{},\alpha _n)`$, so evidently
$$|S_n(T,\tau )|=|S_{n1}(T,\tau )|+|S_{n2}(T,\tau ^{\prime \prime })|.$$
For any permutation $`\tau S_k`$ such that $`\tau _1=k`$ we define $`p(\tau )=1`$ and $`q(\tau )=(\tau _2,\mathrm{},\tau _k)`$, and for any permutation $`\tau S_k`$ such that $`\tau _1=k1`$, and $`\tau _2=k`$ we define $`p(\tau )=2`$ and $`q(\tau )=(\tau _3,\mathrm{},\tau _k)`$. Also, let $`m(\tau )=(m_1,m_2,\mathrm{},m_r)`$ where $`m_i=p(q^{i1}(\tau )))`$ for $`1ir`$, and $`q^0(\tau )=\tau `$, $`q^{i1}(\tau )=q(q^{i2}(\tau )`$ for $`i2`$.
###### Theorem 3.1
Let $`T=\{123,132,213\}`$, $`\tau S_k(T)`$, and $`m(\tau )=(m_1,\mathrm{},m_r)`$. Then
$$e_\tau (x)=\left|\begin{array}{ccccc}u_{m_1}(x)& v_{m_1}(x)& 0& \mathrm{}& 0\\ u_{m_2}(x)& 1& v_{m_2}(x)& \mathrm{}& 0\\ u_{m_3}(x)& 0& 1& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& v_{m_{r1}}(x)\\ u_{m_r}(x)& 0& 0& 0& 1\end{array}\right|,$$
where $`u_1(x)=1`$, $`u_2(x)=\frac{1}{1x}`$, $`v_1(x)=x(1+x)`$, and $`v_2(x)=\frac{x^2}{1x}`$.
###### Proof
Let $`\tau S_k(T)`$; by Lemma 3, there are two cases:
1. $`\tau _1=k`$. So $`p(\tau )=1`$, $`q(\tau )=(\tau _2,\mathrm{},\tau _k)`$, and for all $`nk`$
$$|S_n(T,\tau )|=|S_{n1}(T,\tau ^2)|+|S_{n2}(T,\tau ^2)|.$$
Besides, $`|S_n(T,\tau )|=f_n`$ for all $`nk1`$ (see Simion and Schmidt , Proposition $`15`$), where $`f_n`$ is the $`n`$-th Fibonacci number. Hence
$$e_\tau (x)=x(1+x)e_{q(\tau )}(x)+1.$$
2. $`\tau _1=k1`$. So $`p(\tau )=2`$, $`q(\tau )=(\tau _3,\mathrm{},\tau _k)`$, and for all $`nk`$,
$$|S_n(T,\tau )|=|S_{n1}(T,\tau )|+|S_{n2}(T,\tau ^{\prime \prime })|.$$
Besides, $`|S_n(T,\tau )|=f_n`$ for all $`nk1`$ (see Simion and Schmidt , Proposition $`15`$), where $`f_n`$ is the $`n`$-th Fibonacci number. Hence
$$e_\tau (x)=\frac{x^2}{1x}e_{q(\tau )}(x)+\frac{1}{1x}.$$
Hence, by the definitions and Lemma 1, the theorem holds.
###### Example 5
Let $`T=\{123,132,213\}`$. By Theorem 3.1,
1. $`|S_n(T,3412)|=n`$.
2. $`|S_n(T,4231)|=|S_n(T,3421)|=4`$.
### 3.2 T={123, 132, 231}.
###### Theorem 3.2
Let $`T=\{123,132,231\}`$ and $`\tau S_k(T)`$. Then:
1. there exists $`r`$, $`1rk`$, such that $`\tau =(k,\mathrm{},r+1,r1,\mathrm{},1,r)`$;
2. for all $`nk`$
$$|S_n(T,(k,\mathrm{},r+1,r1,\mathrm{},1,r)|=k1,$$
where $`2rk`$.
###### Proof
$`(i)`$ Let $`\tau S_k(T)`$; put $`r=\tau _n`$. Since $`\tau `$ avoids $`123`$, we see that $`\tau `$ contains $`(r1,\mathrm{},1,r)`$, since $`\tau `$ avoids $`132`$, we see that $`\tau =(\tau _1,\mathrm{},\tau _{kr},r1,\mathrm{},1,r)`$, and since $`\tau `$ avoids $`231`$, we get that $`\tau =(k,\mathrm{},r+1,r1,\mathrm{},1,r)`$.
$`(ii)`$ Let $`\alpha S_n(T,\tau )`$; similarly to $`(i)`$, $`\alpha =(n,\mathrm{},t+1,t1,\mathrm{},1,t))`$ for $`1tn`$, hence $`|S_n(T,\tau )|=k1.`$
###### Example 6
Let $`T=\{123,132,231\}`$; by Theorem 3.2,
$$|S_n(T,4312)|=|S_n(T,4213)|=|S_n(T,3214)|=3.$$
### 3.3 T={123, 213, 231}.
###### Theorem 3.3
Let $`T=\{123,213,231\}`$ and $`\tau S_k(T)`$. Then:
1. there exists $`r`$, $`1rk`$, such that $`\tau =(k,\mathrm{},r+1,1,r,\mathrm{},2)`$;
2. for all $`nk`$
$$|S_n(T,(k,\mathrm{},r+1,1,r,\mathrm{},2))|=k1,$$
where $`2rk`$.
###### Proof
$`(i)`$ Let $`\tau S_k(T)`$ and choose $`r`$ such that $`\alpha _{kr+1}=1`$. Since $`\tau `$ avoids $`213`$, we get that $`\tau _i>\tau _j`$ for all $`i<kr+1<j`$, since $`\tau `$ avoids, $`123`$ we see that $`\tau `$ contains $`(1,r,\mathrm{},2)`$, and since $`\alpha `$ avoids $`231`$, we get that $`\tau =(k,\mathrm{},r+1,1,r,\mathrm{},2)`$.
$`(ii)`$ Let $`\alpha S_n(T,\tau )`$; similarly to $`(i)`$, $`\alpha =(n,\mathrm{},t+1,1,t,\mathrm{},2)`$ for $`1tn`$, hence $`|S_n(T,\tau )|=k1`$.
###### Example 7
Let $`T=\{123,213,231\}`$; by Theorem 3.3,
$$|S_n(T,4312)|=|S_n(T,4132)|=|S_n(T,1432)|=3.$$
### 3.4 T={123, 231, 312}.
###### Theorem 3.4
Let $`T=\{123,231,312\}`$ and $`\tau S_k(T)`$. Then :
1. there exists $`r`$, $`1rk`$, such that $`\tau =(r,\mathrm{},2,1,k,\mathrm{},r+1)`$;
2. for all $`nk`$
$$|S_n(T,(r,\mathrm{},2,1,k,\mathrm{},r+1))|=k1,$$
where $`1rk1`$.
###### Proof
$`(i)`$ Let $`\tau S_k(T)`$; put $`r=\tau _1`$. Since $`\tau `$ avoids $`123`$, we get that $`\tau `$ contains $`(r,k,\mathrm{},r+1)`$, and since $`\tau `$ avoids $`231`$, we see that $`\tau =(r,\mathrm{},k,\mathrm{},r+2,r+1)`$, and since $`\tau `$ avoids $`312`$, we get that $`\tau =(r,\mathrm{},2,1,k,\mathrm{},r+2,r+1)`$ for $`r=1,\mathrm{},k`$.
$`(ii)`$ Let $`rk1`$ and $`\alpha S_n(T,\tau )`$; similarly to $`(i)`$, $`\alpha =(t,\mathrm{},2,1,n,\mathrm{},t+2,t+1)`$ for $`1tn`$. Hence $`|S_n(T,\tau )|=k1`$.
###### Example 8
Let $`T=\{123,231,312\}`$; by Theorem 3.4,
$$|S_n(T,1432)|=|S_n(T,2143)|=|S_n(T,3214)|=3.$$
### 3.5 T={132, 213, 231}.
###### Theorem 3.5
Let $`T=\{132,213,231\}`$ and $`\tau S_k(T)`$. Then :
1. there exists $`r`$, $`1rk`$ , such that $`\tau =(k,\mathrm{},r+1,1,2,\mathrm{},r)`$;
2. for all $`nk`$
$$|S_n(T,(k,\mathrm{},r+1,1,2,\mathrm{},r))|=k1.$$
###### Proof
$`(i)`$ Let $`\tau S_k(T)`$; put $`r=\tau _n`$. Since $`\tau `$ avoids $`231`$, we get that $`\tau `$ contains $`(k,\mathrm{},r+1,r)`$, since $`\tau `$ avoids $`132`$, we see that $`\tau =(k,\mathrm{},r+1,\tau _{kr+1},\mathrm{},\tau _{k1},r)`$, and since $`\tau `$ avoids $`213`$, we get that $`\tau =(k,k1,\mathrm{},r+1,1,2,\mathrm{},r)`$ for $`r=1,\mathrm{},k`$.
$`(ii)`$ Let $`\alpha S_n(T,\tau )`$; similarly to (i), $`\alpha =(n,\mathrm{},t+1,1,2,\mathrm{},t)`$ for $`1tn`$. Hence $`|S_n(T,\tau )|=k1`$.
###### Example 9
Let $`T=\{132,213,231\}`$; by Theorem 3.5,
$$|S_n(T,4321)|=|S_n(T,4312)|=|S_n(T,4123)|=|S_n(T,1234)|=3.$$
## 4 At least four patterns from $`S_3`$ and a pattern from $`S_k`$
By Simion and Schmidt , Proposition $`17`$,
$$|S_n(T)|=2,$$
where $`\{123,321\}TS_3`$, $`|T|=4,5`$, and
$$|S_n(T)|=0,$$
where $`\{123,321\}T`$. Hence, we obtain the following theorem.
###### Theorem 4.1
Let $`TS_3`$, $`|T|4`$ and $`\tau S_k(T)`$. For all $`nk`$
$$|S_n(T,\tau )|=\{\begin{array}{cc}2\delta _{\tau ,(1,2,\mathrm{},k)}\delta _{\tau ,(k,\mathrm{},2,1)}\hfill & 123,321T\hfill \\ 2\delta _{\tau ,(k,\mathrm{},2,1)}\hfill & 123T,321T\hfill \\ 2\delta _{\tau ,(1,2,\mathrm{},k)}\hfill & 123T,321T\hfill \\ 0\hfill & 123,321T.\hfill \end{array},$$
where $`\delta _{x,y}`$ the Kronecker symbol.
## 5 Wilf classes of $`\{T,\tau \}`$, where $`TS_3`$, $`\tau S_4`$
By all the examples in all the sections we obtain Table 1. This table describes all the Wilf classes of sets of permutations avoiding a pattern from $`S_4`$ and a set of patterns from $`S_3`$. It contains $`22`$ Wilf classes for sets $`\{T,\tau \}`$ where $`TS_3`$, $`\tau S_4`$.
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# Charge Ordering and Spin Dynamics in NaV2O5
\[
## Abstract
We report high-resolution neutron inelastic scattering experiments on the spin excitations of NaV<sub>2</sub>O<sub>5</sub>. Below $`T_c`$, two branches with distinct energy gaps are identified. From the dispersion and intensity of the spin excitation modes, we deduce the precise zig-zag charge distribution on the ladder rungs and the corresponding charge order: $`\mathrm{\Delta }_c0.6`$. We argue that the spin gaps observed in the low-T phase of this compound are primarily due to the charge transfer.
\]
The low dimensional inorganic compound NaV<sub>2</sub>O<sub>5</sub> undergoes a phase transition at $`T_c=34`$ K associated with both a lattice distortion and the opening of an energy gap to the lowest triplet spin excitations. While the nature of the low-$`T`$ phase in NaV<sub>2</sub>O<sub>5</sub> is not fully understood it is clear that, unlike CuGeO<sub>3</sub>, the spin-Peierls model does not apply simply to this compound . The spin gap may result from charge-order (CO) rather than the lattice distortion. Indeed, NMR measurements indicate two inequivalent vanadium sites below $`T_c`$, while there exists only one site above. There has been no direct evidence for the connection between CO and a spin gap, nor to distinguish various conjectured spatial distributions of charge . In this letter, we present new results of neutron inelastic scattering (NIS) on the spin excitations in the low-$`T`$ phase that can now resolve these issues.
In NaV<sub>2</sub>O<sub>5</sub>, the vanadium ions have a formal valence of 4.5+. Initially, this was proposed to correspond to an alternation of V<sup>4+</sup> ions, with a spin value $`S=1/2`$, and V<sup>5+</sup> ions with $`S=0`$. At room temperature, NaV<sub>2</sub>O<sub>5</sub> is well described by a quarter-filled two-leg ladder system, with only one type of vanadium site V<sup>4.5+</sup>. From calculations of electronic structure, the strongest orbital overlaps are on the ladder rungs. One expects that the $`S=1/2`$ spins are carried by the V-O-V molecular bonding orbitals, with charge fully delocalized on two sites. As the energy of the anti-bonding orbital is much higher, it can be projected out, and above $`T_c`$, these spins, as they interact in the leg direction ($``$ $`𝐛`$ axis), form an effective uniform quantum Heisenberg spin chain with interactions between chains that are both weaker and frustrated.
At low temperatures, NMR shows this can no longer be so. On each rung, a charge transfer $`\mathrm{\Delta }_c`$ may occur. Taking the average charge on vanadium sites to be 1/2, the charges on the two vanadium sites on a rung are defined through $`n_\pm =(1\pm \mathrm{\Delta }_c)/2`$. Two forms of CO can be considered, the in-line, with the same charge transfer on each rung, and the zig-zag with alternation in the charge along the ladders as shown in fig. 1b. Recent X-ray diffraction measurements established that the lattice structure below $`T_c`$ consists of a succession of distorted and non-distorted ladders of vanadium ions (see fig. 1a). Neglecting inter-ladder diagonal couplings $`J_{}`$, the ladders would behave magnetically as independent spin chains. For one ladder (chain 2 in fig. 1a) distortions in the exchange paths both within the ladder and via neighboring ladders result in an alternation of the effective exchange coupling in the b direction, $`J_{b1}`$ and $`J_{b2}`$. The ladders in which the rungs are distorted (chains 1 and 3), however, remain magnetically uniform as a mirror plane passes through each rung. A minimum magnetic model without CO would be a succession of alternating and uniform chains. An energy gap (expected to be small as it results primarily from the alternation in $`J_{}`$) would characterize the excitation branch of the alternating chains, and there would be no gap for the uniform chains. The initial NIS found, however, two excitation branches, with the same gap at the antiferromagnetic point $`E_g^+=E_g^{}10`$ meV. To analyse these results, a recent spin model used an explicit relation between the spin excitations and the CO. This model assumed a single gap and, as it implies zero intensity of one excitation branch, must be extended to explain the NIS data. A more precise determination of excitations in the low-T phase of NaV<sub>2</sub>O<sub>5</sub> is therefore crucial. In the present work, using high-resolution, the dispersion of the excitations is re-explored in a wider part of the reciprocal space. Moreover, the evaluation of the structure factor, i.e., the (energy-integrated) intensity of each excitation mode, allows us to determine the charge transfer $`\mathrm{\Delta }_c`$.
The single crystal ($`8`$x$`5`$x$`2`$ mm<sup>3</sup>) was grown by a flux method. The NIS measurements were performed at $`T4.2`$ K, on two thermal neutron three-axes spectrometers - IN8 and CRG/CEA-IN22 - at the Institut Laue-Langevin (ILL). On IN8, vertically focusing monochromator PG($`002`$) and Cu($`111`$) were used in conjunction with a vertically focusing analyzer PG($`002`$) and horizontal collimations $`60^{}`$-$`40^{}`$-$`60^{}`$. The final wave vector was kept fixed at $`k_f=4.1`$ Å<sup>-1</sup>. IN22 was operated at $`k_f=2.662`$ Å<sup>-1</sup>, with a PG($`002`$) monochromator and a PG($`002`$) analyzer used in horizontal monochromatic focusing condition (resulting wave vector resolution: $`\delta q0.03`$ r.l.u.) with no collimation. The sample was installed in an “orange” ILL cryostat, with the scattering wave-vector $`𝐐`$ lying in the reciprocal ($`a^{},b^{}`$) ladder plane.
The two branches characterizing the low-energy excitations in NaV<sub>2</sub>O<sub>5</sub> have distinct energy gaps: $`E_g^+E_g^{}`$. This important result will be established when we consider the dispersions in the transverse a direction. First, however, we determine the dispersion in the leg direction (b axis). Examples of constant-energy scans obtained on IN8 as a function of $`Q_b`$ are displayed in fig. 2a (wave-vector components are expressed in reciprocal lattice units, r.l.u.). Increasing the energy, one resolves the single peak seen at 10 meV into propagating modes (dashed lines), whose peak position is shown in fig. 2b. They describe the dispersion of the elementary excitations in the b direction, near the AF chain wave-vector component $`Q_b^{AF}`$ ($`Q_b=0.5`$). They are compared to a dispersion law characteristic of a gapped spin chain: $`E(q_b)=\sqrt{E_g^2+(E_m^2E_g^2)\mathrm{sin}^2(2\pi Q_b)}`$ where $`E_g`$ and $`E_m`$ are the gap and maximum energies of the dispersion, respectively. For both $`E_g=E_g^+`$, $`E_g^{}`$ (the solid and dashed lines, respectively), we evaluate $`E_m=93\pm 6`$ meV. Compared to the prediction for a uniform Heisenberg chain, $`E_m=\pi J_b/2`$, where $`J_b`$ is the exchange in the chain, one obtains for the low-$`T`$ phase, $`J_b60`$ meV.
Second, we consider the dispersions in the a* direction (i.e., along the rungs). A few examples of energy scans performed on IN22 at constant $`𝐐`$ are reported in fig. 3. In general, two peaks are observed. At a few $`Q_a`$ values, however, an extinction occurs. This extinction may concern one of the two modes, or the two modes simultaneously. The left and right panels report data obtained for $`Q_b^{AF}=0.5`$ and $`Q_b=1`$ (equivalent to the zone center of the AF chains $`Q_b^{ZC}`$). As examples, we show for $`Q_b^{AF}`$ that the smallest energy difference between the two observed peaks is obtained for integer $`Q_a`$ values (here, $`Q_a=3)`$, the largest one for half-integer values ($`Q_a=2.5)`$ and an extinction of the two modes occurs at $`Q_a1.75`$ (identical results have been obtained for $`Q_b=1.5`$). Surprisingly, at the chain zone-center $`Q_b^{ZC}`$ we found a small but non-zero intensity for the two excitation branches. The smallest energy difference between the two peaks is now obtained for half-integer values (here, $`Q_a=1.5`$). Extinctions of one of the two modes is observed at $`Q_a=1`$ (on the high-energy mode) and at $`Q_a=2`$ (on the low-energy mode). For all the spectra recorded on IN22, the background (the dotted lines in fig. 3) was carefully determined. Several procedures have been used involving data recorded at low and high temperature (i.e., above $`T_c`$). In fig. 3a, for instance, the open dots used to define the background are obtained from measurements performed at 40 K while in fig. 1b, it is determined from $`Q_b`$ scans performed at low temperature for different energies.
For the analysis, we assume the two observed peaks to belong to two distinct contributions. Their unsymmetrical lineshape is characteristic of gapped excitations undergoing a rapid energy dispersion (as established in fig. 2b). In such a case, a dynamical response function (shown by the dashed lines in fig. 3) is well-suited to fitting, conveniently defined with only 3 parameters: the peak energy $`E_+`$ ($`E_{}`$), an intensity factor $`A_+`$ ($`A_{}`$) and an energy damping $`\mathrm{\Gamma }`$. As $`\mathrm{\Gamma }`$ is mainly fixed by the resolution conditions, it is assumed to be the same for the two contributions ($`\mathrm{\Gamma }0.40.8`$ meV). Together with the background, the agreement with the experiments, shown by the solid lines, is good. The values obtained for $`E_\pm `$ as a function of $`Q_a`$ are shown in fig. 4. We establish several new features. At both $`Q_b^{AF}`$ (solid symbols) and $`Q_b^{ZC}`$ (open symbols), the transverse dispersion consists of two distinct excitation branches which never cross. This justifies our previous statement, namely that there are two distinct energy gaps, $`E_g^+`$ and $`E_g^{}`$. In each branch, the periodicity is $`2\pi Q_a`$: this is twice that previously determined. The two dispersions have the same amplitude, $`\delta J1`$ meV but, remarkably, the upper and lower branches are out of phase, and there is phase inversion between branches at $`Q_b^{ZC}`$ and $`Q_b^{AF}`$. For each excitation branch, the corresponding structure factors, $`S_{b\pm }^{AF}(Q_a)`$ and $`S_{b\pm }^{ZC}(Q_a)`$ are evaluated by integrating the fitted dynamical response functions over a wide energy range (from $`0`$ up to $`E300\mathrm{\Gamma })`$. We estimate systematic error for varying the upper cut off alters the results by at most 5%. The resulting values (dots and squares) are reported in figs. 5a and b. The sums $`S_b^{AF}(Q_a)=S_{b+}^{AF}(Q_a)+S_b^{AF}(Q_a)`$ and $`S_b^{ZC}(Q_a)=S_{b+}^{ZC}(Q_a)+S_b^{ZC}(Q_a)`$ are shown as the stars.
The interpretation of these results is developed in three steps. First, the charge order: each spin is associated with an electronic wave function on the two sites of a rung that depends on $`n_\pm `$. The structure factors for the in-line and zig-zag models are $`S_{Q_b}(Q_a,\omega )=\mathrm{cos}^2\left(\pi Q_a\rho \right)\stackrel{~}{S}(Q_a,Q_b,\omega )+\mathrm{\Delta }_c^2\mathrm{sin}^2\left(\pi Q_a\rho \right)\stackrel{~}{S}(Q_a+\frac{1}{2},Q_b,\omega )`$ and $`S_{Q_b}(Q_a,\omega )=`$ $`\mathrm{cos}^2\left(\pi Q_a\rho \right)\stackrel{~}{S}(Q_a,Q_b,\omega )`$ \+ $`\mathrm{\Delta }_c^2\mathrm{sin}^2\left(\pi Q_a\rho \right)\stackrel{~}{S}`$($`Q_a+\frac{1}{2}`$, $`Q_b+\frac{1}{2},\omega `$), respectively, with $`\omega =E/\mathrm{}`$ and where $`\rho =l/a0.304`$ ($`l`$, rung length and $`a`$, lattice parameter) and $`\stackrel{~}{S}(Q_a,Q_b,\omega )`$ is the structure factor for spins localized on the center of each rung. From the ratios of intensities for different values of momentum, one can extract the charge transfer $`\mathrm{\Delta }_c`$ independent of the form of $`\stackrel{~}{S}`$. In particular, we verify that the order cannot be in-line but predictions agree with a zig-zag order with $`\mathrm{\Delta }_c^20.35`$, i.e., $`\mathrm{\Delta }_c0.6`$. The agreement is particularly good for the sums $`S_b^{AF}(Q_a)`$ and $`S_b^{ZC}(Q_a)`$. For the absolute intensities, we need an explicit form for $`\stackrel{~}{S}`$ which we take from the strongly dimerized limit (SDL), in which the wave function is simply a product of singlets on the stronger bonds. For the in-line model, for example, the SDL would give zero intensity at $`Q_b^{ZC}`$ in contradiction with the observation (the data in fig. 5b). The in-line model can be ruled out. In figs. 5a and b, the predictions provided by the zig-zag model (solid lines) are compared with the experimental total structure factors $`S_b^{AF}(Q_a)`$ and $`S_b^{ZC}(Q_a)`$ (solid and open stars). In fig. 5a, the agreement is obtained with no adjustable parameter except for an overall amplitude factor. Once this factor is determined, the results in fig. 5b depends only on $`\mathrm{\Delta }_c^2`$. As can be seen, a good agreement is obtained for $`\mathrm{\Delta }_c^20.35`$. The low-T phase of NaV<sub>2</sub>O<sub>5</sub> is very well described by the zig-zag model, with a rather large charge transfer.
Second, we consider the transverse dispersions. Later, we explain that the gap is induced by the CO. In fact, there are two distinct gaps because of the structural distortions (implying distinct couplings $`J_{b1,2}^A`$ and $`J_{b1,2}^B`$ as shown in fig. 1b). Due to the CO, 4 different interchain exchange integrals must be considered, $`J_1^{}`$, $`J_2^{}`$, $`J_3^{}`$ and $`J_4^{}`$. The two branches (associated with the gaps $`E_g^+`$ and $`E_g^{}`$) acquire a transverse dispersion, described at $`Q_b^{AF}`$ and $`Q_b^{ZC}`$ by $`E_{b\pm }^{AF}=(E_g^++E_g^{})/2\pm \sqrt{\left(E_g^+E_g^{}\right)^2/4+\delta J^2\mathrm{sin}^2\left(\pi Q_a\right)}`$ and $`E_{b\pm }^{ZC}=(E_g^++E_g^{})/2\pm \sqrt{\left(E_g^+E_g^{}\right)^2/4+\delta J^2\mathrm{cos}^2\left(\pi Q_a\right)}`$, respectively, with $`\delta J=J_1^{}J_2^{}+J_3^{}J_4^{}`$. In fig. 4, these predictions (solid and dashed lines) are compared to the data. Again, a very good agreement is obtained yielding the following evaluation $`E_g^+=10.1\pm 0.1`$, $`E_g^{}=9.1\pm 0.1`$ and $`\delta J=1.2\pm 0.1`$ meV. The dispersion gives directly the alternation in the inter-ladder diagonal bonds $`\delta J`$. If we assume it is dominated by the CO, we can also estimate the average from $`\delta JJ_{}\mathrm{\Delta }_c^2`$ giving $`J_{}2.4`$ meV.
The structure factors can be calculated for each branch. Within the SDL approach, one obtains the contributions $`S_{b\pm }^{AF}(Q_a)`$ and $`S_{b\pm }^{ZC}(Q_a)`$ shown by the dotted and dashed lines in fig. 5. In fig. 5b, the agreement with the data is rather good. In particular, the extinction phenomenon observed for each branch is well reproduced. In fig. 5a, while the sum $`S_b^{AF}(Q_a)`$ is well described, we note a discrepancy between the SDL predictions and the individual structure factors $`S_{b+}^{AF}(Q_a)`$ and $`S_b^{AF}(Q_a)`$. This difficulty could be explained as follows. As shown in fig. 1b, two successive chains are not identical. A charge transfer giving different average valence on the chains will mix intensities at $`Q_b^{AF}`$ without affecting the fluctuations at $`Q_b^{ZC}`$ in qualitative agreement with the observation. Experimental supports for such a charge transfer would be useful. Our proposition for the charge ordering, i.e., the zig-zag model sketched in fig. 1a, and our estimate $`\mathrm{\Delta }_c0.6`$ are based on the total intensities.
Finally, we consider the origin of the gap. As discussed in the introduction, the lattice distortion alone, in isolated ladders, cannot explain the presence of two energy gaps and their size. To analyze the effects of the diagonal couplings $`J_{}`$, we refer to fig. 1b. Each ladder is seen to be a succession of two distinct clusters (shown by ovals in the figure). By exact diagonalization of each cluster, using the effective parameters of a t-J model and adding a potential imposing a charge transfer, we evaluated $`J_{b1}`$ and $`J_{b2}`$ as a function of $`\mathrm{\Delta }_c`$ and the bond alternation $`d=\left|J_{b1}J_{b2}\right|/\left(J_{b1}+J_{b2}\right)`$. For $`\mathrm{\Delta }_c0.6`$, the couplings underestimate the experimental value by a factor of about 2, but as the parameters calculated on the high temperature structure and as such cluster calculations are rather crude, the agreement is satisfactory. The value obtained for the bond alternation can be considered as reasonable: $`d0.0250.030`$. Then, using the experimental value $`J_b60`$ meV, one finds an energy gap $`E_g68`$ meV. This is in a fairly good agreement with the experimental value $`E_g10`$ meV. This simple analysis supports the view that, in NaV<sub>2</sub>O<sub>5</sub>, the gaps are primarily due to the CO. The lattice distortion plays a secondary role, explaining why two distinct branches are observed experimentally, and their separation. In our picture, the magnetic anisotropies are unnecessary.
In conclusion, the CO in NaV<sub>2</sub>O<sub>5</sub> is quantitatively determined by the present NIS measurements . It explains also the energy gaps observed in the low-T phase of this compound.
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# Multiwavelength optical observations of chromospherically active binary systems
## 1 Introduction
This paper is a continuation of our ongoing project of multiwavelength optical observations aimed at studying the chromosphere of active binary systems using the information provided for several optical spectroscopic features that are formed at different heights in the chromosphere. In Paper I (Montes et al. M97 (1997)) we focussed our study on the analysis of the extensively used H$`\alpha `$ chromospheric activity indicator together with simultaneous observations of the less studied He i D<sub>3</sub> and Na i D<sub>1</sub>, D<sub>2</sub> spectral features in a sample of 18 northern active binary systems. In Paper II (Montes et al. 1998a ) the H$`\alpha `$, H$`\beta `$, Na i D<sub>1</sub>, D<sub>2</sub>, He i D<sub>3</sub>, Mg i b triplet, Ca ii H & K and Ca ii infrared triplet lines (Ca ii IRT) of the RS CVn system EZ Pegasi were studied at different orbital phases, also including a high resolution echelle spectrum.
As shown in Paper I and II, with the simultaneous analysis of the different optical chromospheric activity indicators and using the spectral subtraction technique, it is possible to study in detail the chromosphere, discriminating between the different structures: plages, prominences, flares and microflares (see recent studies by Gunn & Doyle G&D97 (1997); Gunn et al. G97 (1997); Lázaro & Arévalo L&A97 (1997); Arévalo & Lázaro A&L99 (1999); Hall & Wolovitz H&W98 (1998); Montes et al. 1998b ; Montes & Ramsey MR98 (1998), MR99 (1999); Montes et al. M99 (1999); Eibe et al. Eibe99 (1999); Oliveira & Foing 1999; Berdyugina et al. Ber99 (1999)).
In this paper we present high resolution echelle spectra of 16 systems selected from ”A Catalog of Chromospherically Active Binary Stars (second edition)” (Strassmeier et al. 1993, hereafter CABS). The spectra were taken during four observing runs (from 1996 to 1999) and include all the optical chromospheric activity indicators from the Ca ii H & K to Ca ii IRT lines. Preliminary results of some of the systems included here can be found in Sanz-Forcada et al. (1998, 1999); Montes et al. (1998c , 2000b ); Latorre et al. (1999, 2000).
In Sect. 2 we give the details of our observations and data reduction. In Sect. 3 we describe the spectroscopic features analysed in this paper: the different chromospheric activity indicators and the Li i $`\lambda `$6707.8 line. Individual results about stellar parameters and the behaviour of the chromospheric excess emission in each system is reported in Sect. 4. Finally, in Sect. 5 the discussion and conclusions are given.
## 2 Observations and Data Reduction
The spectroscopic echelle observations of the chromospherically active binaries analysed in this paper were obtained during four observing runs. Two of them were carried out in 1-2 March 1996 (NOT96 hereafter) and in 5-10 April 1998 (NOT98 hereafter), using the 2.56 m Nordic Optical Telescope (NOT) located at the Observatorio del Roque de Los Muchachos (La Palma, Spain). The Soviet Finnish High Resolution Echelle Spectrograph (SOFIN) was used with an echelle grating (79 grooves/mm), camera Astromed-3200 and a 1152$`\times `$770 pixel EEV P88200 CCD detector. The wavelength range covers from 3765 to 9865 Å, for the first run, and from 3632 to 10800 Å, for the second run. The reciprocal dispersion ranges from 0.07 to 0.18 Å/pixel and the spectral resolution, determined as the full width at half maximum (FWHM) of the arc comparison lines, ranges from 0.15 to 0.60 Å, in both runs.
We also analyse echelle spectra obtained during a 10-night run in 12-21 January 1998 (McD98 hereafter), using the 2.1 m Otto Struve telescope at McDonald Observatory (Texas, USA) and the Sandiford Cassegrain Echelle Spectrograph (McCarthy et al. 1993). This instrument is a prism cross-dispersed echelle mounted at the Cassegrain focus and it is used with a 1200$`\times `$400 Reticon CCD detector. The spectrograph setup was chosen to cover the H$`\alpha `$ (6563 Å) and Ca ii IRT (8498, 8542, 8662 Å) lines. The wavelength coverage is about 6400-8800Å and the reciprocal dispersion ranges from 0.06 to 0.08 Å/pixel. The spectral resolution, determined as the FWHM of the arc comparison lines, ranges from 0.13 to 0.20 Å. In one of the nights, we changed the spectrograph setup to include the He i D<sub>3</sub> (5876 Å) and Na i D<sub>1</sub> and D<sub>2</sub> (5896, 5890 Å) lines, with wavelength coverage of 5600-7000 Å.
In addition, we dispose of echelle spectra taken in 7-8 January 1999 (INT99 hereafter), with the 2.5 m Isaac Newton Telescope (INT) at the Observatorio del Roque de Los Muchachos (La Palma, Spain) using the ESA-MUSICOS spectrograph. This is a fibre-fed cross-dispersed echelle spectrograph, built at the ESA Space Science Department in ESTEC as a replica of the first MUSICOS spectrograph built at Meudon-Paris Observatory (Baudrand & Böhm B&B92 (1992)) and developed as part of MUlti-SIte COntinuous Spectroscopy (MUSICOS <sup>1</sup><sup>1</sup>1http://www.ucm.es/info/Astrof/MUSICOS.html) project. During this observing run, a 2148$`\times `$2148 pixel SITe1 CCD detector was used obtaining wavelength coverage from 3950 Å to 9890 Å. The reciprocal dispersion ranges from 0.06 to 0.12 Å and the spectral resolution (FWHM of the of the arc comparison lines) from 0.15 to 0.4 Å.
In Table 1 we give the observing log. For each observation we list date, UT, orbital phase ($`\phi `$) and the signal to noise ratio (S/N) obtained in the H$`\alpha `$ line region. Table 2 shows the HD number, name and stellar parameters of the active binary systems and the non active stars used as reference stars in the spectral subtraction. The B–V and V–R color indexes and the radius are obtained from the relation with spectral type given by Landolt-Börnstein (Schmidt-Kaler 1982) when individual values are not given in the literature. Other parameters are given by CABS Catalog (Strassmeier et al. 1993) or taken from the references given in the individual results of each star.
The spectra have been extracted using the standard reduction procedures in the IRAF<sup>2</sup><sup>2</sup>2IRAF is distributed by the National Optical Observatory, which is operated by the Association of Universities for Research in Astronomy, Inc., under contract with the National Science Foundation. package (bias subtraction, flat-field division and optimal extraction of the spectra). The wavelength calibration was obtained by taking spectra of a Th-Ar lamp, for NOT and McDonald runs, and a Cu-Ar lamp, for INT run. Finally, the spectra have been normalized by a polynomial fit to the observed continuum.
## 3 Spectroscopic features analysed
### 3.1 Chromospheric activity indicators
The echelle spectra analysed in this paper allow us to study the behaviour of the different optical chromospheric activity indicators formed at different atmospheric heights: Na i D<sub>1</sub>, D<sub>2</sub> and Mg i b triplet (upper photosphere and lower chromosphere), Ca ii IRT lines (lower chromosphere), H$`\alpha `$, H$`\beta `$, Ca ii H & K (middle chromosphere) and He i D<sub>3</sub> (upper chromosphere). The chromospheric contribution in these features has been determined using the spectral subtraction technique described in detail by Montes et al. (1995a, b, c), (see Paper I and II). The synthesized spectrum was constructed using the program STARMOD developed at Penn State (Barden B85 (1985)). The inactive stars used as reference stars in the spectral subtraction were observed during the same observing run as the active stars or were taken from our libraries of late-type stars (see Montes 1998). We have determined the excess emission equivalent width (EW) (measured in the subtracted spectra) and converted to absolute chromospheric flux at the stellar surface. We have estimated the errors in the measured EW taking into account the typical internal precisions of STARMOD (0.5 - 2 km s<sup>-1</sup> in velocity shifts, $`\pm `$5 km s<sup>-1</sup> in V$`\mathrm{sin}i`$, and 5% in intensity weights), the rms obtained in the fit between observed and synthesized spectra in the regions outside the chromospheric features (typically in the range 0.01-0.03) and the standard deviations resulting in the EW measurements. The estimated errors are in the range 10-20%. For low active stars errors are larger and we have considered as a clear detection of excess emission or absorption in the chromospheric lines only when these features in the difference spectrum are larger than 3 $`\sigma `$.
Table 3 gives the H$`\alpha `$ line parameters, measured in the observed and subtracted spectra of the sample. The column (3) gives the orbital phase ($`\phi `$) for each spectrum. In the column (4), H and C mean emission from the hot and cool components, respectively, and T means that at these phases the spectral features cannot be deblended. The column (5) gives the relative contribution of the hot and cool components to the total continuum (S<sub>H</sub> and S<sub>C</sub>), respectively. The column (6) describes the observed H$`\alpha `$ profile, i.e. whether the line is in absorption (A), in emission (E) or totally filled-in by emission (F). The columns (7), (8) and (9) give the following parameters measured in the observed spectrum: the full width at half maximum (W<sub>obs</sub>), the residual intensity (R<sub>c</sub> = $`\frac{\mathrm{F}_\mathrm{c}}{\mathrm{F}_{\mathrm{cont}}}`$) and the equivalent width (EW). The last four columns give the following parameters measured in the subtracted spectrum: the full width at half maximum (W<sub>sub</sub>), the peak emission intensity (I), the excess H$`\alpha `$ emission equivalent width (EW(H$`\alpha `$)), corrected for the contribution of the components to the total continuum and the logarithm of the absolute flux at the stellar surface (logF<sub>S</sub>(H$`\alpha `$)) obtained with the calibration of Hall (Hall96 (1996)) as a function of (V–R).
In Table 4 we list the parameters (I, FWHM, EW) of the broad and narrow components used in the two Gaussian-component fit to the H$`\alpha `$ subtracted emission profile. We have performed this fit in the stars that show broad wings. See the comments for each individual star in Sect. 4 and the interpretation of these components given in Sect. 5.
Table 5 gives the H$`\beta `$ line parameters, measured in the observed and subtracted spectra, as in the case of the H$`\alpha `$ line. In this table we also give the ratio of excess emission EW in the H$`\alpha `$ and H$`\beta `$ lines, $`\frac{\mathrm{EW}(\mathrm{H}\alpha )}{\mathrm{EW}(\mathrm{H}\beta )}`$, and the ratio of excess emission $`\frac{\mathrm{E}_{\mathrm{H}\alpha }}{\mathrm{E}_{\mathrm{H}\beta }}`$ with the correction:
$$\frac{\mathrm{E}_{\mathrm{H}\alpha }}{\mathrm{E}_{\mathrm{H}\beta }}=\frac{\mathrm{EW}(\mathrm{H}\alpha )}{\mathrm{EW}(\mathrm{H}\beta )}0.24442.512^{(BR)}$$
given by Hall & Ramsey (H&R92 (1992)) that takes into account the absolute flux density in these lines and the color difference in the components. We have used this ratio as a diagnostic for discriminating between the presence of plage-like and prominence-like material at the stellar surface, following the theoretical modelling by Buzasi (Buz89 (1989)) who found that low E/E ($``$ 1-2) can be achieved both in plages and prominences viewed against the disk, but that high ratios ($``$ 3-15) can only be achieved in extended regions viewed off the limb. The study of chromospherically active binaries by Hall & Ramsey (H&R92 (1992)) has demonstrated the presence of large amounts of extended, prominence-like material in these stars.
We also analyse the possible filling-in of the core of the Na i D<sub>1</sub> and D<sub>2</sub> lines as other chromospheric activity indicator as well as the behaviour of the He i D<sub>3</sub> line, which can be in absorption, filled-in due to frequent low-level flaring or in emission due to flares (see Paper I and II; Saar et al. 1997; Montes et al. 1996b , M99 (1999)).
Table 6 gives the Ca ii H & K and H$`ϵ`$ lines parameters, measured in the observed and subtracted spectra. In columns (5) and (6) we list the EW for the K and H lines, obtained by reconstruction of the absorption line profile (described by Fernández-Figueroa et al. FFMCC94 (1994), hereafter FFMCC). In columns (7), (8) and (9) we give the EW for the K, H and H$`ϵ`$ lines using the spectral subtraction technique (explained by Montes et al. 1995c, 1996a) and corrected for the contribution of the components to the total continuum. In columns (10), (11) and (12) we list the corresponding logarithm of the surface flux obtained by means of the linear relationship between the absolute surface flux at 3950 Å (in erg cm<sup>-2</sup> s<sup>-1</sup> Å<sup>-1</sup>) and the colour index (V–R) by Pasquini et al. (1988).
Table 7 gives the Ca ii IRT lines parameters, measured in the observed spectra by reconstruction of the absorption line profile and using the spectral subtraction. The columns of this table have the same meaning as in Table 6 for the Ca ii H & K lines. The absolute fluxes at the stellar surface have been obtained using the calibration of Hall (Hall96 (1996)) as a function of (V–R). For the observing runs in which the $`\lambda `$8542 and $`\lambda `$8498 lines are included we also give the ratio of excess emission EW, E<sub>8542</sub>/E<sub>8498</sub>, which is also an indicator of the type of chromospheric structure that produces the observed emission. In solar plages, values of E<sub>8542</sub>/E<sub>8498</sub> $``$ 1.5-3 are measured, while in solar prominences the values are $``$ 9, the limit of an optically thin emitting plasma (Chester Ch91 (1991)). However, the observations of active stars (Chester et al. Ch94 (1994); Arévalo & Lázaro A&L99 (1999)) indicate that these lines exhibit markedly different behaviour. The E<sub>8542</sub>/E<sub>8498</sub> ratios found in these stars are smaller (closer to the optically thick value of one) than solar plages. These values indicate that the Ca ii IRT emission arises predominantly in chromospheric plages.
### 3.2 The Li i $`\lambda `$6707.8 line
The resonance doublet of Li i at $`\lambda `$6708 Å is an important diagnostic of age in late-type stars since it is destroyed easily by thermonuclear reactions in the stellar interior. It is well-known that a large number of chromospherically active binaries shows Li i abundances higher than other stars of the same mass and evolutionary stage (Barrado et al. Ba97 (1997), Ba98 (1998); Paper I, Montes & Ramsey MR98 (1998)). This line is only included in our echelle spectra in the McD98 and INT99 observing runs. In Fig. 25 we have plotted representative spectra of OU Gem, BF Lyn and HU Vir in this spectral region. A K1III reference star with some photospheric lines identified has been also plotted in order to better identify the expected position of the Li i line. It was only possible to measure the equivalent width of the Li i absorption line in the SB1 system HU Vir. In the case of the SB2 systems OU Gem and BF Lyn the possible small absorption Li i of one or both components is blended with photospheric lines of the other component.
## 4 Individual Results
In this section we describe the behaviour of the above mentioned chromospheric activity indicators for each star of the sample. The profiles of H$`\alpha `$, H$`\beta `$, Ca ii H & K and Ca ii IRT are displayed from Fig. 1 to 23. For each system we have plotted the observed spectrum (solid-line) and the synthesized spectrum (dashed-line) in the left panel and the subtracted spectrum (dotted line) in the right panel. The name of the star, the observing run (NOT96, NOT98, McD98, INT99) and the orbital phase ($`\phi `$) of each spectrum are given in every figure. The He i D<sub>3</sub> line, for selected stars of the sample, is displayed in Fig. 24.
### 4.1 UX Ari (HD 21242)
This double-lined spectroscopic binary (G5V/K0IV) is a well known RS CVn system and extensively studied in the literature (Carlos & Popper C&P71 (1971); Bopp & Talcott BoT78 (1978); Huenemoerder et al. 1989; Raveendran & Mohin 1995). Recently, Duemmler & Aarum (D&A00 (2000)) have given a new orbit determination, which we have adopted in Table 2. Our previous H$`\alpha `$ observations (Montes et al. 1995a, b) showed clear H$`\alpha `$ emission above the continuum from the cool component. This emission was superimposed to the weak absorption of the hot component. Our spectrum in the Ca ii H & K region (Montes et al. 1995c) showed strong emission from the cool component and a weak H$`ϵ`$ emission line. We also detected a flare in this system through simultaneous H$`\alpha `$, Na i D<sub>1</sub>, D<sub>2</sub> and He i D<sub>3</sub> observations (Montes et al. 1996b; Paper I).
In the new observation (NOT96) we observe intense emission in the Ca ii H & K, H$`ϵ`$, H$`\alpha `$ and Ca ii $`\lambda `$8662 lines and a filled-in absorption H$`\beta `$ line from the cool component (see Fig. 1). The excess emission measured in all these activity indicators is larger than in our previous observations of this system in quiescent state in 1992 and 1995. The small emission we observe in He i D<sub>3</sub> (Fig. 24) confirms the high level of activity of UX Ari in this observation. The He i D<sub>3</sub> line has been observed as clear emission in this system in other occasions associated to flare-like events (Montes et al. 1996b and references therein).
### 4.2 12 Cam (BM Cam, HD 32357)
This single-lined spectroscopy binary was classified by Bidelman (Bi64 (1964)) as a K0 giant. He also noticed Ca ii H & K emission. Later, Hall et al. (Hall95 (1995)) revised the system spectroscopically and photometrically and obtained new values of orbital parameters, given in Table 2. Eker et al. (Ek95 (1995)) observed that H$`\alpha `$ profiles showed an asymmetric shape with a round shoulder in the red wing and a steeper blue wing. They also confirmed the variable H$`\alpha `$ filling, which was suspected by Strassmeier et al. (1990) too. In our previous observations of this system (FFMCC) we found the H$`\alpha `$ line filled-in and strong Ca ii H & K emission.
In the new spectrum (Fig. 2) we find strong emission in the Ca ii H & K lines. Thanks to the higher resolution of this spectrum it is possible to see now that both lines exhibit self-absorption with blue asymmetry. Small H$`ϵ`$ emission is observed, but it was impossible to deblend it from the Ca ii H line. Both H$`\alpha `$ and H$`\beta `$ appear in absorption with a slight filling-in, the H$`\alpha `$ filling-in is smaller than the one corresponding to other epochs. The H$`\alpha `$ line shows excess absorption in the red wing similar to the asymmetric shape of H$`\alpha `$ observed by Eker et al. (Ek95 (1995)). The application of the spectral subtraction technique reveals that the He i D<sub>3</sub> line appears as an absorption feature (Fig. 24). This fact is more frequent in giants than in dwarfs (Paper I). Finally, the Ca ii IRT ($`\lambda `$$`\lambda `$ 8542, 8662 Å) absorption lines are clearly filled-in.
### 4.3 V1149 Ori (HD 37824)
This single-lined spectroscopic binary, classified as K1III + F by Bidelman & MacConnell (BiMcC73 (1973)), was listed as a G5IV star by Hirshfeld & Sinnott (1982). Our previous observations revealed clear excess H$`\alpha `$ emission (Montes et al. 1995a, b; Paper I), strong Ca ii H & K and H$`ϵ`$ emission (Montes et al. 1995c) and clear absorption in the He i D<sub>3</sub> line in the subtracted spectrum (Paper I).
In the new observation (NOT98, Fig. 3), the H$`\alpha `$ and H$`\beta `$ lines show a filled-in absorption profile, and clear absorption is observed in the blue wing of both lines. This excess absorption in the blue wing, not observed in our previous observations of this system, could be indicative of variable mass motion. To confirm this behaviour we took a new spectrum with the ESA-MUSICOS spectrograph in January 2000 (forthcoming paper) in which the blue wing of the H$`\alpha `$ line is in emission, confirming the high variability of the H$`\alpha `$ line profile in this system. Strong emission is observed in the Ca ii H & K lines but the H$`ϵ`$ line is not detected. The Ca ii IRT lines ($`\lambda `$$`\lambda `$ 8542, 8662 Å) show a strong filling-in. The He i D<sub>3</sub> line appears in absorption (Fig. 24).
### 4.4 OU Gem (HD 45088)
OU Gem is a bright (V= 6.79, Strassmeier et al. 1990) and nearby (d= 14.7 pc, ESA ESA97 (1997)) BY Dra-type SB2 system (K3V/ K5V) with an orbital period of 6.99 days and a noticeable eccentricity (Griffin & Emerson Gr&E75 (1975)). Both components show Ca ii H & K emission, though the primary shows slightly stronger emission than the secondary. The H$`\alpha `$ line is in absorption for the primary and filled-in for the secondary (Bopp Bo80 (1980); Bopp et al. 1981a , b; Strassmeier et al. 1990; Montes et al. 1995a, b, 1996). Dempsey et al. (1993a ) observed that the Ca ii IRT lines were filled-in. This binary was detected by the WFC on board the ROSAT satellite during the all-sky survey (Pounds et al. 1993; Pye et al. 1995). OU Gem has 1.7$`\times `$10<sup>29</sup> ergs<sup>-1</sup> X-ray luminosity, typical value of the BY Dra systems (Dempsey et al. 1993b , De97 (1997)). The photometric variability was discovered by Bopp et al. (1981a ) and they also computed a 7.36-day photometric period. It is interesting that the orbital and rotational periods differ in 5% due to the appreciable orbital eccentricity (e= 0.15), according to Bopp (Bo80 (1980)). Although BY Dra systems are main-sequence stars, their evolutionary stage is not clear. OU Gem has been listed by Soderblom et al. (1990) and Montes et al. (2000a , 2000c ) as a possible member of the UMa moving group (300 Myr), indicating that it may be a young star.
The H$`\alpha `$ line: In the observed spectra, we see an absorption line for the primary star and a nearly complete filling-in for the secondary star. After applying the spectral subtraction technique, clear excess H$`\alpha `$ emission is obtained for the two components, being stronger for the hot one (see Fig. 4 upper panel). The excess H$`\alpha `$ emission EW is measured in the subtracted spectrum and corrected for the contribution of the components to the total continuum. We took one spectrum in this region in Dec-92 (Montes et al. 1995b). At the orbital phase of this observation ($`\phi `$= 0.48) we could not separate the emission from both components and we measured the total excess H$`\alpha `$ emission EW relative to the combined continuum. We obtained a similar value to Mar-96, Apr-98 and Jan-99 values obtained adding up the excess emission EW from the two components.
The H$`\beta `$ line: Looking at the observed spectra, we only see the H$`\beta `$ line for the primary, in absorption. After applying the spectral subtraction technique small excess H$`\beta `$ emission is obtained for the two components (see Fig. 4 lower panel). We have obtained, in general, $`\frac{E_{H\alpha }}{E_{H\beta }}`$ values larger than three for the two components, so the emission can come from prominences.
The Ca ii H $`\&`$ K and H$`ϵ`$ lines: We observe that both components of this binary have the Ca ii H & K and H$`ϵ`$ lines in emission. We can also see that the excess Ca ii H & K emission of the hot star is larger than the one of the cool star (Fig. 5 upper panel). The measured excess Ca ii H & K emission of both components is larger in the two spectra of the NOT98 observing run than in the NOT96 spectrum. Overlapping between the H$`ϵ`$ line of one star and the Ca ii H line of the other only allows to see the H$`ϵ`$ line of the cool star at orbital phase 0.19, and H$`ϵ`$ of the hot star otherwise.
The Ca ii IRT lines: In the observed spectra, we can see that both components of OU Gem show the Ca ii IRT lines in emission superimposed to the corresponding absorption. After applying the spectral subtraction technique, clear excess emission appears for the two components, being clearly stronger for the hot one (see Fig. 5 lower panel).
### 4.5 $`\sigma `$ Gem (HD 62044)
This single-lined spectroscopic binary belongs to the long-period group of RS CVn binary systems. It was classified as K1III, but the radius obtained by Duemmler et al. (Du97 (1997)) seems to be too small for a giant star. In Table 2 we have adopted the orbital and physical parameters updated by Duemmler et al. (Du97 (1997)). Strong and variable Ca ii H & K emission always centered at the absorption line has been reported by Bopp (Bo83 (1983)), Strassmeier et al. (1990), FFMCC, Montes M95 (1995), Montes et al. 1996b . Our previous observation in the H$`\alpha `$ line region (Montes et al. 1995a, b) revealed small excess emission, similar to that reported by Strassmeier et al. (1990) and Frasca & Catalano (F&C94 (1994)). Variable excess H$`\alpha `$ emission anti-correlated with spot regions has been found by Zhang & Zhang (1999).
The new observations (NOT96, NOT98) show strong Ca ii H & K emission lines and small emission in the H$`ϵ`$ line. After applying the spectral subtraction, a small filling-in is observed in the H$`\alpha `$, H$`\beta `$ and Ca ii IRT lines (Fig. 6). We observe in this system notable He i D<sub>3</sub> absorption (Fig. 24). All the activity indicators show an increase in the 1998 observation ($`\phi `$= 0.88) with respect to 1996 one ($`\phi `$= 0.81). The emission in the Ca ii H & K lines in these two observations is noticeably larger than in our previous observations at different epochs and orbital phases.
### 4.6 BF Lyn (HD 80715)
This double-lined spectroscopic binary with spectral types K2V/\[dK\] shows variable H$`\alpha `$ emission and strong Ca ii H & K, H$`ϵ`$ and Ca ii IRT emission from both components (Barden & Nations BN85 (1985); Strassmeier et al. 1989b). In our previous observation (Montes et al. 1995c) we found strong emission in the Ca ii H & K lines from both components with very similar intensity and the H$`ϵ`$ line in emission. The orbital period is 3.80406 days (Barden & Nations BN85 (1985)), and Strassmeier et al. (1989b), from photometric observations, found that BF Lyn is a synchronized binary with a circular orbit.
In the four runs analysed in this paper we have obtained 11 spectra of BF Lyn at different orbital phases. We have used the original heliocentric Julian date on conjunction (T<sub>conj</sub>) given by Barden & Nations (BN85 (1985)) to calculate the orbital phases since we discovered a mistake in the Strassmeier et al. (1993) catalog where the orbital data from Barden & Nations (BN85 (1985)) are compiled. In the original paper the epoch was determined using Modified Julian Date (MJD), that is why 0.5 days must be added to the Strassmeier et al. (1993) date, who used the 2440000.0 Julian day as a reference date.
The H$`\alpha `$ line: We took several spectra of BF Lyn in the H$`\alpha `$ line region in four different epochs and at different orbital phases. In all the spectra (Fig. 7) we can see the H$`\alpha `$ line in absorption from both components. The spectral subtraction reveals that both stars have excess H$`\alpha `$ emission. At some orbital phases, near to the conjunction, it is impossible to separate the contribution of both components. The excess H$`\alpha `$ emission of BF Lyn shows variations with the orbital phase for both components, but the hot star is the most active in H$`\alpha `$. In Fig 8 we have plotted for the McD 98 observing run the excess H$`\alpha `$ emission EW versus the orbital phase for the hot and cool components. The highest excess H$`\alpha `$ emission EW for the hot component has been reached at about 0.4 orbital phase and the lowest value is placed at about 0.9 orbital phase, whereas the cool component shows the highest excess H$`\alpha `$ emission EW at near 0.9 orbital phase and the lowest value between 0.2 and 0.4 orbital phases. The variations of the excess H$`\alpha `$ emission EW for both components are anti-correlated, which indicates that the chromospheric active regions are concentrated on faced hemispheres of both components, but at about 0.4 and 0.9 orbital phases for the hot and cool component, respectively. The same behaviour is also found in Ca ii IRT. The excess H$`\alpha `$ emission EW also shows seasonal variations, for instance, the values of the INT99 observing run are very different, specially for the cool component, from McD 98 values at very similar orbital phase.
The H$`\beta `$ line: Five spectra in the H$`\beta `$ line region are available. In all of them the H$`\beta `$ line appears in absorption from both components. The application of the spectral subtraction technique reveals clear excess H$`\beta `$ emission from both stars (see Fig. 9). The ($`\frac{E_{H\alpha }}{E_{H\beta }}`$) values that we have found for this star allow us to say that the emission comes from extended regions viewed off the limb.
The Ca ii H $`\&`$ K and H$`ϵ`$ lines: We took four spectra in the Ca ii H & K region during the NOT (96 & 98) observing runs (Fig. 10). Another spectrum was taken in 1993 with the 2.2 m telescope at the German Spanish Astronomical Observatory (CAHA) (Montes et al., 1995c). These spectra exhibit clear and strong Ca ii H & K and H$`ϵ`$ emission lines. At 0.02, 0.43, and 0.54 orbital phases the emission from both components can be deblended using a two-Gaussian fit (see Fig. 10). In the case of CAHA 93 run, the H$`ϵ`$ emission line from the hot component is overlapped with the Ca ii H emission of the cool component. The excess Ca ii H $`\&`$ K and H$`ϵ`$ emission changes with the orbital phase during the NOT 98 run in the same way as the corresponding excess Ca ii $`\lambda `$8542 and H$`\alpha `$ emission. The excess Ca ii H $`\&`$ K emission EW also shows seasonal variations, for instance, the values of CAHA 93 observing run are lower than NOT 96 & 98 values.
The Ca ii IRT lines: In all the spectra we can see the Ca ii IRT lines in emission from both components (Fig. 11). As in the case of H$`\alpha `$, the Ca ii IRT emission shows variations with the orbital phase for both components. In Fig. 8 we have plotted, for the McD 98 observing run, the excess Ca ii $`\lambda `$8542 emission EW versus the orbital phase for the hot and cool components. The variations of the excess Ca ii emission EW for both components are anti-correlated and they show the same behaviour as the excess H$`\alpha `$ emission EW.
### 4.7 IL Hya (HD 81410)
IL Hydrae is a typical RS CVn star with very strong Ca ii H & K emission (Bidelman & MacConnell BiMcC73 (1973)). The 12.86833-day photometric period derived by Raveendran et al. (1982) is very close to the orbital period. It was found to be an X-ray source and a microwave emitter (Mitrou et al. 1996). From a photometric analysis, Cutispoto (Cut95 (1995)) estimated the secondary to be a G8V star. Later, Donati et al. (Do97 (1997)) detected the secondary component in the optical range and they calculated a 1.0 R radius for it. Weber & Strassmeier (1998) gave a K0III-IV type for the primary and computed the first double-lined orbit of IL Hya. Later, Raveendran & Mekkaden (1998) gave a new orbital solution and just recently, Fekel et al. (F99 (1999)) have presented updated SB2 orbital elements which we have adopted and they are given in Table 2 (we have corrected the heliocentric Julian date on conjunction (T<sub>conj</sub>) so that the primary is in front). The multiwavelength Doppler images presented by Weber & Strassmeier (1998) revealed a cool polar spot and several features at low latitudes. These authors also found that the H$`\alpha `$ EW showed sinusoidal variation which was in phase with the photospheric light curve.
We have taken one spectrum of IL Hya (NOT98) with 0.02 orbital phase which is very close to the conjunction, so the very weak lines of the secondary can not be detected in any wavelength. In the observed spectrum (Fig. 12), the H$`\alpha `$ line can be seen as a filled-in absorption line with a superimposed 1.04 Å (48 km s<sup>-1</sup>) red-shifted absorption feature, as obtained from a two-Gaussian fit. After applying the spectral subtraction, clear excess emission is observed. The excess emission shows an asymmetric profile due to the red-shifted absorption feature. Similar H$`\alpha `$ profiles were observed by Weber & Strassmeier (1998) in this system, but the red-shifts measured in their spectra were larger (1.24 Å) and remained constant during a rotational cycle. This behaviour could be due to mass motions that change from one epoch to another, but a combination of several dynamical processes may be involved. A filling-in is also found in the H$`\beta `$ line. According to the value obtained for the corrected ratio of the excess emission EW of both lines, the emission may be ascribed to an extended region viewed off the limb. The He i D<sub>3</sub> line appears in absorption (Fig. 24), but no filling-in is detected in the Na i D<sub>1</sub>, D<sub>2</sub> lines. The Ca ii H $`\&`$ K lines are observed in emission. Furthermore, the Ca ii IRT lines show clear central emission reversal.
### 4.8 FG UMa (HD 89546)
FG UMa is the least studied star of our sample. Bidelman (Bi81 (1981)) included it in his Catalogue of stars with Ca ii H & K emission. This star is identified as a single-lined binary by CABS. A 21.50-day photometric period has been obtained from the automated monitoring that Henry et al. (1995a) carried out. From spectroscopic measurements, they confirmed V$`\mathrm{sin}i`$= (15$`\pm `$2) kms<sup>-1</sup> and a G8IV spectral type. These authors also mentioned that, according to an unpublished orbital analysis, the system is synchronized and circularized. We have also taken from them the orbital period and the radius. Some indication of possible eclipses is noted in the Hipparcos Catalogue (ESA ESA97 (1997)). Fluxes of the Ca ii H & K emission lines have been calculated by Strassmeier (1994b) and a filled-in and variable H$`\alpha `$ line has been reported by Henry et al. (1995a).
We have not got enough data to be able to compute the orbital phases corresponding to the two spectra (NOT98) that we present here. However, there is a change of 0.2 in the photometric phase between both observations. We have not found any evidence of the secondary star through the whole spectral range. Moreover, according to the appearance of some Ti i and Fe i lines (Paper II) we suggest that the observed spectra correspond to a luminosity class more evolved than subgiant (in agreement with the radius calculated by Henry et al. (1995a) who suggested a luminosity class III-IV). The presence of a notable He i D<sub>3</sub> absorption line (Fig. 24) encourages this conclusion. In the observed spectra of the H$`\alpha `$ and H$`\beta `$ lines (Fig. 13), both absorption lines show a clear filling-in. The spectral subtraction allows us to compare the two observations. As it can be read in Table 5 the ratio of the excess emission EW is typical of extended regions viewed off limb, and a significant variation, mainly due to H$`\beta `$, is obtained for the two different nights. Furthermore, excess emission is detected in the blue wing of the H$`\alpha `$ line. Similar behaviour was mentioned by Henry et al. (1995a). Strong filling-in is observed in the Ca ii IRT lines. The Ca ii H & K lines present strong emission with clear self-absorption with blue-shifted asymmetry in both observations.
### 4.9 LR Hya (HD 91816)
It is a double-lined spectroscopic binary, classified as a BY Dra-type system, that contains two almost identical K-type dwarf components (Bopp et al. Bo&84 (1984)). The orbital parameters were determined by Fekel et al. (F88 (1988)) who suggested a K0V spectral type for both components. The photometric period of 3.1448 days, given in CABS, was reported by Bopp et al. (Bo&84 (1984)), but following observation campaigns could not confirm that value. In fact, the results obtained are not consistent (Strassmeier 1989; Cutispoto Cut91 (1991), Cut93 (1993)) and point out low-amplitude rotational modulation due to the development and decline of small active regions at different longitudes of both components.
We only have got one spectrum (NOT98) of this system at 0.58 orbital phase, so that the chromospheric activity indicators from both components can be easily analyzed. The H$`\alpha `$ absorption line (Fig. 14) shows a weak filling-in for both components, as it has been previously mentioned by Strassmeier et al. (1990). No filling-in is detected in the H$`\beta `$ line. Although the S/N ratio in the Ca ii H & K region is very low in this observation and the synthetic and observed spectra are not well matched, we can clearly see moderate emission in the Ca ii H & K lines from both components. The Ca ii IRT lines of the two stars exhibit a clear filling-in. The measured excess emission in the different lines are very similar in both components, although a bit larger in the red-shifted one.
### 4.10 HU Vir (HD 106225)
HU Vir is a late-type star (K0III-IV) with strong Ca ii H $`\&`$ K emission noted by Bidelman (Bi81 (1981)) for the first time. Recently, Fekel et al. (F99 (1999)) have discovered that HU Vir is a triple system with a long period of about 6.3 years and we have taken from them the spectral type, T<sub>conj</sub> and the rotational period. The B–V colour index has been taken from Hipparcos Catalogue (ESA ESA97 (1997)) and V$`\mathrm{sin}i`$ from Fekel (F97 (1997)). Fekel et al. (1986) observed the H$`\alpha `$ line in emission and Strassmeier & Fekel (1990) found the H$`ϵ`$ line in emission. Such emission lines are seen in the most active RS CVn type systems. Strassmeier (1994a) found a big, cool polar spot from Doppler imaging and two hot plages 180<sup>o</sup> apart from the H$`\alpha `$ and Ca ii line-profile analysis. The chromospheric plages seemed to be spatially related to two large appendages of the polar spot. Broadening in the H$`\alpha `$ profile suggested mass flow in a coronal loop connecting the two plage regions. Hatzes (Hat98 (1998)) used the Doppler imaging technique to derive the cool spot distribution. He found a large spot at latitude 45<sup>o</sup> and a weak polar spot with an appendage. The polar spot was considerably smaller than similar features found on other RS CVn stars. From an analysis of the H$`\alpha `$ variations he also found evidence for a plage located at high ($``$ 70<sup>o</sup>) latitude, near the polar extension.
The H$`\alpha `$ line: Strassmeier (1994a) identified three distinctive features in the H$`\alpha `$ line: blue-shifted emission, central sharp absorption and red-shifted broad absorption. Hatzes (1998) found similar behaviour in this line. In our spectra (Fig. 15), the H$`\alpha `$ line always appears in emission, although its intensity is variable. Moreover, the emission is blue-shifted at some orbital phases. Hall & Ramsey (H&R92 (1992), H&R94 (1994)) explained the blue-shifted emission as prominence-like structures. We can see red-shifted broad absorption at the 0.54 orbital phase but we never observe central sharp absorption. Walter & Byrne (1998) said that there was growing evidence for continuous low-level mass in-fall, seen as red-shifted absorption in H$`\alpha `$ line profiles. We can also notice that the subtracted H$`\alpha `$ profiles have broad and variable extended wings which are not well matched using a single-Gaussian fit. These profiles have been fitted using two Gaussian components. The parameters of the broad and narrow components used in the two-Gaussian fit are given in Table 4 and the corresponding profiles are plotted in the left panel of Fig. 15. These broad wings are observed at different orbital phases and in different epochs. In some cases the blue wing is noticeable stronger than the red wing and the fit is better matched when the broad component is blue-shifted with respect to the narrow component. We have interpreted these broad components as microflaring activity that occurs in the chromosphere of this very active star (Paper I, II; Montes et al. 1998b ). The contribution of the broad component to the total EW of the line ranges from 32% to 66% which is in the range observed in the stars analysed in Paper I & II. Strassmeier (1994a), in 1991, and Hatzes (1998), in 1995, observed strong emission between the 0.27 and 0.51 orbital phases. We have obtained the strongest emission at the 0.44 orbital phase, in January 1998, and at 0.29, in April 1998 (see Table 3). Thus, our 0.29-0.44 orbital phase interval is similar to Strassmeier and Hatzes’s orbital phase interval, so we can conclude that HU Vir has an active longitude area (which corresponds to that orbital phase interval) since 1991. Similar active longitudes have been found by other authors in several chromospherically active binaries (Oláh et al. 1991; Henry et al. 1995b; Jetsu 1996; Berdyugina & Tuominen BerT98 (1998)).
The H$`\beta `$ line: A nearly complete filling-in is observed (Fig. 16 upper panel). After applying the spectral subtraction technique, clear excess H$`\beta `$ emission is obtained. We have obtained $`\frac{E_{H\alpha }}{E_{H\beta }}`$ values larger than three (see Table 5), so the emission can come from prominences.
The Ca ii H $`\&`$ K and H$`ϵ`$ lines: We can observe very strong Ca ii H $`\&`$ K emission and an important H$`ϵ`$ emission line superimposed to the wide Ca ii H line (Fig. 16 lower panel). The H$`ϵ`$ line in emission indicates that HU Vir is a very active system. Similar strong emission was found in our previous observation obtained in March 1993 at 0.71 orbital phase (Montes et al. 1995c). We can also notice that the largest values of the excess Ca ii H $`\&`$ K emission EW appear at the 0.29 and 0.38 orbital phases (see Table 6). It is in agreement with the H$`\alpha `$ line behaviour.
The Ca ii IRT lines: HU Vir shows the Ca ii IRT lines in emission above the continuum (Fig. 17). We observe that the emission is centered at its corresponding absorption. We also notice that the subtracted profiles have broad wings due to microflares according to Montes et al. (1997). The excess Ca ii $`\lambda `$8542 emission EW (see Table 7) behaves like the excess H$`\alpha `$ emission EW.
The He i D<sub>3</sub> line: We have not detected any significant absorption or emission for He i D<sub>3</sub> (Fig. 24), contrary to the absorption observed in other giants. This total filling-in of the He i D<sub>3</sub> line could be explained (Saar et al., 1997; Paper II) as a filling-in due to the low-level flaring (microflares) that takes place in this very active star according to the H$`\alpha `$ broad component that we have found.
The Li i $`\lambda `$6707.8 line: The Li i absorption line is clearly observed in the eight spectra of the McD98 observing run (see Fig.25). The mean EW obtained is 56 $`\pm `$11 (mÅ). At this spectral resolution the Li i line is blended with the nearby Fe i $`\lambda `$6707.41 Å line. We have corrected the total EW measured, EW(Li i+Fe i), by subtracting the EW of Fe i calculated from the empirical relationship with (B–V) given by Soderblom et al. (1990) (EW(Fe i)=24 (mÅ)). Finally, the corrected EW(Li i)<sub>corr</sub>=32 (mÅ) was converted into abundances by means of the curves of growth computed by Pallavicini et al. (1987), obtaining $`\mathrm{log}\mathrm{N}(\mathrm{Li}\mathrm{i})`$ =1.2 (on a scale where $`\mathrm{log}\mathrm{N}(\mathrm{H})`$ =12.0) with an accuracy of the $``$ 0.30 dex. This value is larger than the lower limit reported by Barrado et al. (Ba98 (1998)) for this star.
### 4.11 DK Dra (HD 106677)
This is a double-lined spectroscopic binary with almost identical components of spectral types K1III and Ca ii H & K emission from both components (Bopp et al. Bo&79 (1979); Fekel et al. F86 (1986); Strassmeier 1994b). Eker et al. (Ek95 (1995)) reported variable nature of H$`\alpha `$ and, using a subtracted spectrum, found emission of similar intensity from both components. In our previous observations (Montes et al. 1995a, b; FFMCC) we found a broad excess H$`\alpha `$ emission line and the Ca ii H & K lines in emission, but all of them were taken at orbital phases near to the conjunction, so it was impossible to distinguish the contribution of each component.
The spectra analysed in this paper were taken at 0.90 (NOT98) and 0.84 (NOT98) orbital phases (see Fig. 18). In both spectra we observe the H$`\alpha `$, H$`\beta `$ and Ca ii $`\lambda `$8662 lines filled-in, and the Ca ii H & K lines in emission. At 0.90, we cannot separate the contribution of each component in H$`\alpha `$, H$`\beta `$ and Ca ii H & K lines, but due to the large wavelength of the Ca ii $`\lambda `$8662 line, a double peak is clearly observed in the subtracted spectrum of this line. In the 0.84 spectrum, we can see that the H$`\alpha `$ line is filled-in for both stars. Moreover, we can observe the excess H$`\alpha `$ emission of the blue-shifted component is bigger than that of the red-shifted one. The excess H$`\beta `$ emission seems to come only from the blue-shifted component. Although the S/N ratio in the Ca ii H & K lines region is low, we can clearly see broad and two-piked emission in the H & K lines. In the subtracted spectrum of the Ca ii IRT $`\lambda `$$`\lambda `$ 8542 and 8662 lines the emission from both components is clearly separated, being the emission of the blue-shifted one slightly larger. All of this indicates that the blue-shifted component is slightly more active than the other component. Very small absorption is observed in the He i D<sub>3</sub> line in the expected position for both components (see Fig. 24), contrary to the notable absorption observed in other giants of the sample. This is probably due to the blend with other photospheric lines of both stars in this SB2 system.
### 4.12 BQ CVn (HD 112859)
A quarter of century ago Schild (1973) classified this star as a peculiar G8III-IV. Henry et al. (1995a) detected its Ca ii H & K emission and noticed that at red wavelengths the spectrum was double-lined although the intensity of the lines of both components was very different. They suggested a K0III primary and a late-F spectral type for the secondary. We have taken from these authors both the orbital and rotational period and the radius. On the other hand, Strassmeier (1994b) confirmed the strong Ca ii H & K emission.
Orbital parameters have not been published, so we cannot calculate the orbital phase of our spectrum (NOT98). However, some conclusions can be obtained looking at the observed spectrum (Fig. 19): at infrared and red wavelengths, the secondary lines are not clear. A very weak blue-shifted absorption feature can be ascribed to the secondary in the H$`\alpha `$ region. At shorter wavelengths, as in the Na i D<sub>1</sub>, D<sub>2</sub> lines region, the spectral lines of the secondary are more conspicuous. In the Ca ii H & K lines region the contribution of the F star is evident in the broad Ca ii absorption lines, where a clear red-shift of the emission can be seen. Taking into account what has been said above, we have calculated the synthesized spectrum using a F8V star as a template one for the secondary. In the observed spectrum, the H$`\alpha `$ line shows a filling-in for the primary star. The H$`\beta `$ line is slightly filled-in. The presence of the He i D<sub>3</sub> absorption line is detected (Fig. 24). In our spectrum we observe strong Ca ii H & K emission and a weak H$`ϵ`$ emission line from the cool star. Finally, the Ca ii IRT lines exhibit reversal emission.
### 4.13 IS Vir (HD 113816)
IS Vir is a single-lined spectroscopic binary classified as K0III by Henry et al. (1995a). We have taken from them the orbital and rotational periods and the radius given in Table 2. Strong Ca ii H & K emission was observed (Buckley et al. Bu87 (1987); Strassmeier 1994b; Montes et al. 1995c). In our previous observation of this system in the Ca ii H & K lines region in Mar-93 at 0.68 orbital phase (Montes et al. 1995c) we found strong emission in the H & K lines with intensity above the continuum at 3950 Å, but lower than reported by Strassmeier (1994b).
In the new spectrum (NOT 98) (Fig. 20), the H$`\alpha `$ line and the Ca ii IRT lines show intense filled-in absorption, whereas the H$`\beta `$ line only shows a slight filling-in. In the Ca ii H & K lines region the S/N ratio is low and the synthetic and observed spectra are not well matched, but a strong emission in the H & K lines well above the continuum is observed. The He i D<sub>3</sub> line appears in absorption (Fig. 24).
### 4.14 BL CVn (HD 115781)
This double-lined spectroscopy binary was classified by Hall (Hall90 (1990)) as K1III + FIV, and later as K1III + G5IV by Stawikowski & Glebocki (1994). Lines et al. (1985) found this system to have a photometric period of 9.31$`\pm `$0.06 days and an amplitude of 0.16 mag. They attributed the light variability to the ellipticity effect because the orbital period was twice the photometric period and times of maximum brightness occurred at times of maximum positive radial velocity. It is a nearly-synchronous binary: its orbital period is 18.6917$`\pm `$0.0011 days (Griffin & Fekel Gr&F88 (1988)) and its rotational period is 18.70 days (Stawikowski & Glebocki 1994). The orbital eccentricity cannot be far from zero (Griffin & Fekel Gr&F88 (1988)). Fekel et al. (F86 (1986)) found V$`\mathrm{sin}i`$ values of (35$`\pm `$2) km<sup>-1</sup> and (7$`\pm `$2) km<sup>-1</sup> for the primary and secondary, respectively. The great line broadening of BL CVn and the ellipsoidal light variations might suggest that the K giant star is close to filling its Roche lobe (Griffin & Fekel Gr&F88 (1988)). Moderate H$`\alpha `$ absorption is found by Fekel et al. (F86 (1986)). Strassmeier et al. (1990) found strong Ca ii H & K emission.
In our present observations (NOT98, Fig. 21) the H$`\alpha `$ and H$`\beta `$ lines appears in absorption in the observed spectrum. After applying the spectral subtraction technique, we only obtain small excess H$`\alpha `$ emission. Broad emission is observed in the Ca ii H & K lines. We observe small absorption in the He i D<sub>3</sub> line (Fig. 24). A small filling-in is obtained in the Ca ii $`\lambda `$8542 and $`\lambda `$8662 lines.
### 4.15 BM CVn (HD 116204)
This single-lined spectroscopy binary was classified by Keenan (1940) as K1III. It is a nearly-synchronous binary: its orbital period is (20.6252$`\pm `$0.0018) days (Griffin & Fekel Gr&F88 (1988)) and its rotational period is (20.66$`\pm `$0.03) days (Strassmeier et al. 1989a). The orbit is judged to be circular (Griffin & Fekel Gr&F88 (1988)). It is also a relatively fast rotator, Fekel et al. (F86 (1986)) found rotationally broadened lines with V$`\mathrm{sin}i`$= (15$`\pm `$2) km<sup>-1</sup>. This system shows strong Ca ii H & K emission (Bidelman Bi83 (1983)), together with an emission-line spectrum typical of RS CVn stars in the IUE ultraviolet region, but H$`\alpha `$ is an absorption line (Griffin & Fekel Gr&F88 (1988)).
In our present spectrum (NOT98, Fig. 22) we observe the H$`\alpha `$ line as nearly total filled-in absorption. After applying the spectral subtraction technique, strong excess H$`\alpha `$ emission is obtained. The H$`\beta `$ line appears as an absorption line in the observed spectrum and clear excess emission is observed in the subtracted spectrum. The $`\frac{E_{H\alpha }}{E_{H\beta }}`$ ratio obtained is larger than 3, which indicates that the emission would arise from extended regions viewed off the limb. The Ca ii H & K lines show strong emission, with a blue-shifted self-absorption feature, and small emission is also detected in the H$`ϵ`$ line. The He i D<sub>3</sub> line appears in absorption (Fig. 24). A clear emission reversal is observed in the Ca ii IRT $`\lambda `$8542 and $`\lambda `$8662 lines.
### 4.16 MS Ser (HD 143313)
Griffin (Gr78 (1978)) first observed the binary nature of MS Ser calculating the orbital elements for this system (see Table 2). Griffin gave its T<sub>0</sub> in MJD, and this yielded Strassmeier et al. (1993) to a bad calculation of the T<sub>0</sub> in HJD. Griffin also proposed K2V/K6V as spectral types of the components, based on photometric arguments for the secondary star. Bopp et al. (1981b ) observed a variable filling-in of the H$`\alpha `$ line and calculated a photometric period of 9.60 days, slightly different from the orbital period. Miller & Osborn (1996) confirmed the value of the photometric period and Strassmeier et al. (1990) observed strong emission in the Ca ii H & K composite spectrum. Dempsey et al. (1993a ) noted some filling-in in the Ca IRT lines for MS Ser, but not reverse emission. Alekseev (A99 (1999)) made a photometric and polarimetric study of MS Ser, calculating a spot area of 15% of the total stellar surface, and observed some seasonal variations. Finally, Osten & Saar (1998) revised the stellar parameters for MS Ser and suggested K2IV/G8V as a better classification. We have also found that the primary component may have a luminosity class IV or higher based on our analysis of some metallic lines, like the Ti i lines, the Hipparcos data and the Wilson Bappu effect (see Sanz-Forcada et al. 1999).
Two spectra of MS Ser were taken in the NOT98 observing run. Moreover, we analyse here another spectrum taken on 12th June 1995 with the 2.2 m telescope at the German Spanish Observatory (CAHA) in Calar Alto (Almería, Spain), using a Coudé spectrograph with the f/3 camera, CCD RCA #11 covering two ranges: H$`\alpha `$ (from 6510 to 6638 Å) and H$`\beta `$ (from 4807 to 4926 Å), with a resolution of $`\mathrm{\Delta }\lambda `$ 0.26 in both cases.
The H$`\alpha `$ and H$`\beta `$ lines: In both lines, we observe a nearly total filling-in in the 1995 spectra and small absorption in the spectra taken in the NOT98 observing run. After applying the spectral subtraction, a clear filling-in in the H$`\alpha `$ and H$`\beta `$ lines is observed in the three spectra (see Fig. 23). The H$`\alpha `$ line of this system is highly variable. We have found, in the present spectra, a variable filling-in whereas Bopp et al. (1981b ) found H$`\alpha `$ was a weak emission line. In the three new spectra we have taken of this system with the FOCES echelle spectrograph in July 1999 (forthcoming paper) we observe variable H$`\alpha `$ emission well above the continuum.
The Ca ii H $`\&`$ K and H$`ϵ`$ lines: Strong emission in these lines and the H$`ϵ`$ line in emission arising from the primary component was observed in our previous observation of this system in March 1993 at 0.16 orbital phase (Montes et al. 1995c). In the present observations (Fig. 23) we have deblended the emission arising from both components in the spectrum taken near to the quadratures ($`\phi `$ =0.21). The strongest emission, centered at the absorption line, arises from the K2IV component, which is the component with larger contribution to the continuum. The red-shifted and less intense emission corresponds to the G8V secondary component. In the $`\phi `$ =0.55 observation we cannot separate the contribution from both stars. The H$`ϵ`$ line appears in emission in both spectra. The emission intensity observed in our 1993 and 1998 spectra is larger than the emission intensity observed in the 1988 spectrum presented by Strassmeier et al. (1990).
The Ca ii IRT lines: A clear emission reversal is observed in the core of the Ca ii IRT absorption lines $`\lambda `$8542 and $`\lambda `$8662 in both spectra (Fig. 23). After applying the spectral subtraction technique, we can see small excess emission arising from the secondary component, in the 0.21 spectrum, as in the case of the Ca ii H & K lines. This emission reversal observed here clearly contrasts with the only filling-in in these lines reported by Dempsey et al. (1993a ).
The He i D<sub>3</sub> line: We can distinguish (Fig. 24) the He i D<sub>3</sub> line as a very small absorption from the primary star. A slight variation is observed between the spectra at 0.21 and 0.55 orbital phases. The luminosity class of this star (IV-III) and the SB2 nature of the system could be the reason of this small absorption in the He i D<sub>3</sub> line in relation with that observed in other giants of the sample.
## 5 Discussion and conclusions
In this paper we have analysed, using the spectral subtraction technique, high resolution echelle spectra of 16 chromospherically active binary systems. These spectra include all the optical chromospheric activity indicators from the Ca ii H & K to Ca ii IRT lines, and in some observing runs the Li i $`\lambda `$6707.8 line too.
H$`\alpha `$ emission above the continuum is observed in UX Ari and HU Vir, in the rest of the systems excess H$`\alpha `$ emission is clearly detected in the subtracted spectra. Filled-in absorption H$`\beta `$ line profiles have been found in all the stars of the sample except LR Hya and BL CVn, which also have a lower level of activity in the other chromospheric activity indicators.
Very broad and variable extended wings have been found in the subtracted H$`\alpha `$ profile of the very active star HU Vir. These line profiles are not well matched using a single-Gaussian fit, and have been fitted using two Gaussian components (narrow and broad). Similar behaviour has been found in other very active stars (Hatzes 1995; Jones et al. 1996; Stauffer et al. 1997; Paper I, II, Montes et al. 1998b ). We have interpreted this broad component emission as arising from microflaring (high turbulence emitting plasma) activity that takes place in the chromosphere by similarity with the broad components found in the chromospheric Mg ii h & k lines (Wood et al. 1996; Busà et al. Busa99 (1999)) and in several transition region lines of active stars observed with GHRS-HST (Linsky & Wood 1994; Linsky et al. 1995; Wood et al. 1996, 1997; Dempsey et al. 1996a , b; Robinson et al. 1996) and recently confirmed with STIS-HST observations (Pagano et al. 2000). In some cases the H$`\alpha `$ line is asymmetric and the fit is better matched when the broad component is blue-shifted or red-shifted with respect to the narrow component. These asymmetries are also observed during the impulsive and gradual phases of solar and stellar flares (Montes et al. 1996b , M99 (1999); Montes & Ramsey MR99 (1999)), and favour the interpretation of the broad component as due to upward and downward motions produced by microflaring in the chromosphere.
Excess absorption and emission is observed in the wings of the H$`\alpha `$ and H$`\beta `$ lines of several systems. Absorption features are detected by us in the blue wing of V1149 Ori and in the red wing of 12 Cam (also observed by Eker et al. (Ek95 (1995))), IL Hya (also detected by Weber & Strassmeier (1998)) and HU Vir (also seen by Strassmeier (1994a) and Hatzes (1998)). We detected excess emission in the blue wing of FG UMa too. Similar red-shifted absorption features were already seen by other authors in the single star OP And (Fekel et al. F86 (1986)) and in the binaries VY Ari (Bopp et al. Bo&89 (1989)); XX Tri (Bopp et al. Bo&93 (1993)) and IN Vir (Strassmeier 1997). Blue-shifted emission was also detected in the single giant YY Men (Vilhu et al. 1991). Several dynamical processes, or a combination of them, could be the origin of these blue-shifted or red-shifted absorption features: plage- and prominence-like structures (Hall & Ramsey H&R92 (1992), H&R94 (1994); Neff 1996; Eibe Eibe98 (1998)); continuous low-level mass infall (Walter & Byrne 1998; Walter 1999; Eibe et al. Eibe99 (1999)); local velocity fields and mass motions due to magnetic field inhomogeneities possibly coupled with a loop-like geometry (Strassmeier 1994a); fluctuations of both the column density and temperature gradient within the chromosphere (Smith & Dupree 1988); stellar winds and anti-winds (Linsky et al. 1995).
Prominence-like extended material viewed off the limb has been detected in many stars of the sample according to the high ratios E/E obtained. Prominences viewed against the disk seem to be present at some orbital phases in the dwarfs OU Gem and BF Lyn.
The application of the spectral subtraction reveals that the He i D<sub>3</sub> line appears as an absorption feature (Fig. 24) in mainly all the giants of the sample as in the case of the stars analysed in Paper I. Total filling-in of He i D<sub>3</sub> is observed in the very active star HU Vir, similar to the behaviour observed in EZ Peg (Paper II). These results are in agreement with the behaviour reported by Saar et al. (1997). These authors found that while for few active stars the He i D<sub>3</sub> line behaves ”normally”, increasing in absorption with increasing rotation, and showing consistent correlations with other activity indicators, the behaviour clearly diverges (large filling-in) when stars become very active, suggesting that the line could be filled-in due to frequent low-level flaring. In the most evolved stars the behaviour is different as a consequence of the lower chromospheric densities of these stars.
Ca ii H & K emission is observed in all the stars of the sample. Small emission in the close H$`ϵ`$ line is also detected in some of the more active stars. In some systems like OU Gem, BF Lyn, LR Hya, DK Dra and MS Ser the emission from both components is clearly deblended in these lines. Self-reversed absorption core with red asymmetry (I(K<sub>2V</sub>) $`<`$ I(K<sub>2R</sub>) is detected in the Ca ii H & K lines of the giants 12 Cam, FG UMa and BM CVn. The self-reversed feature is a consequence of the line formation process in the chromosphere (depth variation of the line source function in an atmosphere having a chromospheric temperature rise). Asymmetries in these profiles provide information on velocity fields in the line formation regions. In these three giants we observed a small red asymmetry (indicative of outward mass flux, wind), contrary to the blue asymmetry (indicative of upward propagating waves, but not large wind) normally observed in giant stars hotter than spectral type K3 (Stencel 1978) and also observed by us (FFMCC) in the giants (V1817 Cyg and V1764 Cyg).
The Ca ii IRT lines result to be a very useful chromospheric activity indicator too. We have found that all the stars analysed here show a clear filled-in absorption line profile or even notable emission reversal (UX Ari, OU Gem, BF Lyn, IL Hya, HU Vir, BQ CVn, BM CVn, MS Ser). Thanks to the higher resolution of our spectra, we were able to detect emission reversal in the Ca ii IRT lines in some systems in which previous studies (Linsky et al. 1979; Dempsey et al. 1993a ) only reported a filling-in. An increase in the level of activity of these stars could also be the cause of these detections. When both components of the binary system are active the excess emission detected in the Ca ii IRT lines is much more easily deblended thanks to the large wavelength of these lines (see e.g. DK Dra, Fig. 18). We have found E<sub>8542</sub>/E<sub>8498</sub> ratios in the range $``$ 1-2 which is indicative of optically thick emission in plage-like regions, in contrast with the prominence-like material inferred by the E/E ratios. These small E<sub>8542</sub>/E<sub>8498</sub> ratios, also found by Chester et al. (Ch94 (1994)) and Arévalo & Lázaro (A&L99 (1999)) in other active binaries, indicate the existence of distinct sources of Balmer and Ca ii IRT emission and suggest that the activity of these very active stars is not simply a scaled-up version of solar activity.
Some systems have been observed at different orbital phases and different epochs, covering all the orbital cycle and it was possible to study the variability of the chromospheric emission. The excess H$`\alpha `$ and Ca ii IRT emission of BF Lyn in the McD98 observing run shows anti-correlated variations with the orbital phase (see Fig 8). This anti-correlation could indicate that the chromospheric active regions are concentrated on faced hemispheres of both components but at about 0.4 and 0.9 orbital phases for the hot and cool component, respectively. Evidence of an active longitude area (0.29-0.44 orbital phase interval) has been found in HU Vir when we compare with the higher level of activity in this phase interval also reported by Strassmeier (1994a), in 1991, and Hatzes (1998), in 1995.
###### Acknowledgements.
We would like to thank Dr. L.W. Ramsey for collaborating in the McDonald observing run and the staff of McDonald observatory for their allocation of observing time and their assistance with our observations. We would like to thank Dr. B.H. Foing for allow us to use the ESA-MUSICOS spectrograph at Isaac Newton Telescope. We would also like to thank the referee S. Catalano for suggesting several improvements and clarifications. This work has been supported by the Universidad Complutense de Madrid and the Spanish Dirección General de Enseñanza Superior e Investigación Científica (DGESIC) under grant PB97-0259.
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# HIGH 𝐸_𝑇 PHOTOPRODUCTION AND THE PHOTON STRUCTURE
## 1 Introduction
The study of high transverse energy ($`E_T`$) dijet photoproduction provides both a test of pQCD and information on the structure of the photon. In events with jets of sufficiently high $`E_T`$, the contribution of non-perturbative effects, which are currently poorly understood, is suppressed. The effect of hadronisation is also reduced, allowing a more meaningful comparison of hadron-level data with parton-level next-to-leading-order (NLO) calculations. In the kinematic regime of the current measurements, the parton densities of the proton are well constrained, in contrast to the parton densities of the photon. The parameterisations currently fit $`F_2^\gamma `$ data from $`e^+e^{}`$ experiments which do not well constrain the gluon density of the photon or the quark densities at high $`x_\gamma `$, the fraction of momentum of the photon carried by the struck parton.
In this paper, events with an almost real photon (virtuality, $`Q^20\mathrm{GeV}^2`$) containing at least two jets reconstructed using the $`k_T`$ clustering algorithm $`^\mathrm{?}`$ are analysed. Measurements of events with a virtual photon and with jets of lower $`E_T`$ are discussed in detail elsewhere $`^{\mathrm{?},\mathrm{?}}`$. The H1 and ZEUS collaborations have analysed similar amounts of luminosity; 36 $`\mathrm{pb}^1`$ and 38 $`\mathrm{pb}^1`$ respectively, allowing measurements up to $`E_T^{\mathrm{jet}}90`$ GeV. Improvements in the understanding of pQCD for jet photoproduction has led to the agreement of many independent calculations to within $`510\%`$ $`^{\mathrm{?},\mathrm{?}}`$.
## 2 Results
Cross sections as a function of the transverse momentum of the leading jet and the average transverse momentum of the two jets are shown in Fig. 1 compared to expectations from the Pythia Monte Carlo $`^\mathrm{?}`$ program and an NLO calculation. The kinematic regime is shown in Fig. 1 where the leading jet was required to have $`P_T^{\mathrm{jet1}}>25`$ GeV. The transverse momentum of the jet falls by three orders of magnitude up to values of $``$ 90 GeV. The cross section measurements are well described by the pQCD calculation in both shape and normalisation.
As the pQCD calculations describe the measurements at very high values of transverse energy, lowering the energy and considering other quantities provides a test of the current photon parton density functions (PDF’s). It has already been demonstrated that at forward regions of pseudorapidity of the jet, $`\eta ^{\mathrm{jet}}`$, the NLO calculations underestimate the data for all choices of photon PDF $`^{\mathrm{?},\mathrm{?}}`$. Requiring that $`x_\gamma ^{\mathrm{obs}}`$, the fraction of the photon’s energy participating in the production of the two highest energy jets;
$$x_\gamma ^{\mathrm{obs}}=\frac{_{\mathrm{jet1},2}E_T^{\mathrm{jet}}e^{\eta ^{\mathrm{jet}}}}{2yE_e},$$
(1)
where $`yE_e`$ is the initial photon energy, satisfies $`x_\gamma ^{\mathrm{obs}}>0.75`$, enhances the direct photon component of the cross section. The NLO calculations were shown to agree with the data in this region justifying the validity of the calculations and indicating that the discrepancy when no $`x_\gamma ^{\mathrm{obs}}`$ cut is applied arises due to the inadequacies of the photon PDF’s.
This was investigated further by considering the cross section in $`x_\gamma ^{\mathrm{obs}}`$. The measured cross sections in $`x_\gamma ^{\mathrm{obs}}`$ for increasing slices in transverse energy of the leading jet are shown in Fig. 2 compared to an NLO calculation using the AFG-HO $`^\mathrm{?}`$ photon PDF. The measured data clearly lie up to $`5060\%`$ above the pQCD calculation for low values of $`x_\gamma ^{\mathrm{obs}}`$ and the trend is maintained up to high values of $`E_T^{\mathrm{jet1}}`$.
Having established a region of phase space in which the data is not described by the pQCD calculations, further studies were made to ascertain if the problem lies with the calculation or the photon PDF. This was done by considering the cross section in $`|\mathrm{cos}\theta ^{}|`$, the angle between the dijet and beam axes in the dijet centre-of-mass system. This measurement directly tests the parton-parton dynamics of the hard sub-process and was performed in two regions of $`x_\gamma ^{\mathrm{obs}}`$ ($`x_\gamma ^{\mathrm{obs}}<0.75`$ and $`x_\gamma ^{\mathrm{obs}}>0.75`$). Additional requirements were made to remove the biases in the distribution due to the $`E_T^{\mathrm{jet}}`$ and $`\eta ^{\mathrm{jet}}`$ cuts. Cuts on the dijet invariant mass, $`M_{jj}>39`$ GeV and the average pseudorapidity of the jets, $`0<\overline{\eta }<1`$, were made.
The cross section is shown in Fig. 3 compared to the absolute prediction and the shape of the NLO calculation. As has already been observed, the higher $`x_\gamma ^{\mathrm{obs}}`$ region is well described and the lower $`x_\gamma ^{\mathrm{obs}}`$ region poorly described in normalisation. However, the shape of the cross section for $`x_\gamma ^{\mathrm{obs}}<0.75`$ is well described by the NLO calculation indicating that the calculation describes the parton-parton dynamics of the hard sub-process.
## 3 Conclusions
Dijet photoproduction has been measured up to $`E_T^{\mathrm{jet}}`$ 90 GeV and compared to pQCD calculations. Clear evidence is seen that the calculations give a good description of the scattering process. However, inadequacies in the current parameterisations of the photon PDFs are evident. The use of the data presented here in future fits would greatly improve our understanding of the photon PDF.
## References
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# Clifford Homomorphisms and Higher Spin Dirac Operators
## 0 Introduction
Let $`M`$ be a $`n`$-dimensional spin manifold and $`\mathrm{𝐒𝐩𝐢𝐧}(M)`$ be its spin structure. The irreducible unitary representation $`(\pi _\rho ,V_\rho )`$ of $`Spin(n)`$ induces a vector bundle $`𝐒_\rho `$ associated to $`\mathrm{𝐒𝐩𝐢𝐧}(M)`$. We consider the canonical covariant derivative $``$ mapping from $`\mathrm{\Gamma }(M,𝐒_\rho )`$ to $`\mathrm{\Gamma }(M,𝐒_\rho T^{}(M))`$. If we decompose the tensor bundle $`𝐒_\rho T^{}(M)`$ into irreducible bundles, then we can construct first order differential operators associated to $``$,
$$D_{\lambda _k}^\rho :\mathrm{\Gamma }(M,𝐒_\rho )\stackrel{}{}\mathrm{\Gamma }(M,𝐒_\rho T^{}(M))\stackrel{\mathrm{\Pi }_{\lambda _k}^\rho }{}\mathrm{\Gamma }(M,𝐒_{\lambda _k}).$$
(0.1)
Here $`\mathrm{\Pi }_{\lambda _k}^\rho `$ is the orthogonal projection onto an irreducible bundle $`𝐒_{\lambda _k}`$ from $`𝐒_\rho (M)T^{}(M)`$. These operators are called the higher spin Dirac operators, the generalized gradient, or the Stein-Weiss operators, which are conformal invariant first order differential operators. Many basic and geometric operators on $`M`$ are constructed in this way; the Dirac operator, the twistor operator (see , , and ), the exterior derivative, its adjoint, the conformal killing operator (see and ), the Rarita-Schwinger operator (see ) and so on. As in these examples, the higher spin Dirac operators and their properties are closely related with geometry and topology of $`M`$ (see the references given above and , , , , etc.).
Now, we consider the Dirac operator $`D`$ as a basic example. We denote the spinor representation of $`Spin(n)`$ by $`(\pi _\mathrm{\Delta },V_\mathrm{\Delta })`$. Then the Dirac operator $`D`$ is nothing else but the higher spin Dirac operator $`D_\mathrm{\Delta }^\mathrm{\Delta }`$, and the projection $`\mathrm{\Pi }_\mathrm{\Delta }^\mathrm{\Delta }`$ is realized by using the Clifford algebra. The Clifford algebra gives us a lot of information about the Dirac operator. For instance, we can prove the ellipticity, the conformal invariance, and the Bochner identity of $`D`$ (see and ). As this example, the projection $`\mathrm{\Pi }_{\lambda _k}^\rho `$ in (0.1) is an essential tool to investigate the higher spin Dirac operators. Since the projection is defined fiberwise, we consider the tensor representation $`(\pi _\rho \pi _{\mathrm{Ad}},V_\rho 𝐑^n)`$ and the projection $`\mathrm{\Pi }_{\lambda _k}^\rho `$ onto an irreducible component $`V_{\lambda _k}`$. For any $`u`$ in $`𝐑^n`$, we have a homomorphism $`p_{\lambda _k}^\rho (u)`$ from $`V_\rho `$ to $`V_{\lambda _k}`$,
$$𝐑^n\times V_\rho (u,\varphi )p_{\lambda _k}^\rho (u)\varphi :=\mathrm{\Pi }_{\lambda _k}^\rho (\varphi u)V_{\lambda _k}.$$
(0.2)
We call this homomorphism the Clifford homomorphism. On the spinor space, $`p_\mathrm{\Delta }^\mathrm{\Delta }(u)`$ is the usual Clifford action $`u`$. So the Clifford homomorphism is a generalization of the usual Clifford action.
The aim of this paper is to study the Clifford homomorphisms and the higher spin Dirac operators. Since we can extend the Clifford homomorphisms to bundle homomorphisms, some relations among the Clifford homomorphisms induce the ones among the higher spin Dirac operators. In particular, we have general Bochner identities. Furthermore, these identities allow us to give a lower bound of the first eigenvalue for a Laplace type operator on each associated bundle.
In section 1, we review the Clifford algebras, the spin groups, and their representations. In particular, we introduce the Casimir operator and the conformal weights. In section 2, we decompose the representation space $`V_\rho 𝐑^n`$ into irreducible components as a $`Spin(n)`$-module (or $`𝔰𝔭𝔦𝔫(n)`$-module) and define the Clifford homomorphisms. We investigate these homomorphisms and give an explicit formula of the projection $`\mathrm{\Pi }_{\lambda _k}^\rho `$. As an example, we give explicit decompositions of $`V_\mathrm{\Delta }𝐑^n`$ and $`\mathrm{\Lambda }^k(𝐑^n)𝐑^n`$, where $`(\pi _\mathrm{\Delta },V_\mathrm{\Delta })`$ means the spinor representation and $`(\pi _{\mathrm{\Lambda }^k},\mathrm{\Lambda }^k(𝐑^n))`$ is the $`k`$-th exterior product representation. In section 3 and 4, we study the higher spin Dirac operators and give some relations among them. Then we have the Bochner identities for the higher spin Dirac operators. In the last section, we give a lower bound of the first eigenvalue of a Laplace type operator constructed in section 4. This is a generalization of the eigenvalue estimates known about the Dirac operator and the Laplace-Beltrami operator.
## 1 Preliminaries: representations of $`Spin(n)`$
In this section, we give a short review to the Clifford algebras, the spin groups, and their representations. Let $`𝐑^n`$ be the n-dimensional Euclidean space with orthonormal basis $`\{e_k\}_{k=1}^n`$. The Clifford algebra $`Cl_n`$ associated to $`𝐑^n`$ is an associative algebra with unit generated by $`\{e_k\}`$ under the relations $`e_ke_l+e_le_k=2\delta _{kl}`$. We denote the complexification of $`Cl_n`$ as $`𝐂l_n`$.
We define the Lie algebra $`𝔰𝔭𝔦𝔫(n)`$ in $`Cl_n`$ by
$$𝔰𝔭𝔦𝔫(n):=\mathrm{span}_𝐑\{[e_k,e_l]|\mathrm{\hspace{0.25em}1}k,ln\},$$
(1.1)
where $`[e_k,e_l]:=e_ke_le_le_k`$ and the Lie bracket is $`[a,b]:=abba`$ for $`a`$ and $`b`$ in $`𝔰𝔭𝔦𝔫(n)`$. The basis of $`𝔰𝔭𝔦𝔫(n)`$ is $`\{[e_k,e_l]\}_{k<l}=\{2e_ke_l\}_{k<l}`$. We put $`\mathrm{exp}(a):=\frac{a^n}{n!}`$ for $`a`$ in $`Cl_n`$ and obtain the spin group $`Spin(n)`$,
$$Spin(n):=\mathrm{exp}𝔰𝔭𝔦𝔫(n)Cl_n.$$
(1.2)
If we would like to think of the spin group as a double covering group of $`SO(n)`$, we use the adjoint representation $`\pi _{\mathrm{Ad}}`$ of $`Spin(n)`$ on $`𝐑^n`$,
$$Spin(n)\times 𝐑^n(g,x)\pi _{\mathrm{Ad}}(g)x=gxg^1𝐑^n.$$
(1.3)
Then we have the isomorphisms $`𝔰𝔬(n)𝔰𝔭𝔦𝔫(n)`$, where $`[e_k,e_l]`$ corresponds to the $`n\times n`$ matrix $`4E_{kl}+4E_{lk}`$. The Killing form gives an inner product on $`𝔰𝔭𝔦𝔫(n)`$ as follows:
$$[e_k,e_l],[e_i,e_j]=32\delta _{ik}\delta _{jl}\text{for }k<l\text{ and }i<j.$$
(1.4)
The irreducible unitary representations of $`Spin(n)`$ or $`𝔰𝔭𝔦𝔫(n)`$ are parametrized by dominant weights $`\rho =(\rho ^1,\mathrm{},\rho ^m)𝐙^m(𝐙+1/2)^m`$, satisfying that
$$\rho ^1\mathrm{}\rho ^{m1}|\rho ^m|,\text{for }n=2m,$$
(1.5)
$$\rho ^1\mathrm{}\rho ^{m1}\rho ^m0,\text{for }n=2m+1.$$
(1.6)
The dominant weight $`\rho `$ is the highest weight of the corresponding representation $`(\pi _\rho ,V_\rho )`$. Here, we use the same notation for the representation of $`Spin(n)`$ and its infinitesimal representation of $`𝔰𝔭𝔦𝔫(n)`$. When writing dominant weights, we denote a string of $`j`$ $`k`$’s for $`k`$ in $`𝐙(𝐙+1/2)`$ by $`k_j`$. For example, $`((\frac{3}{2})_p,(\frac{1}{2})_{mp})`$ is the weight whose first $`p`$ components are $`\frac{3}{2}`$ and others are $`\frac{1}{2}`$.
We shall introduce the Casimir operator and the conformal weights. Let $`(\pi _\rho ,V_\rho )`$ be an irreducible representation with highest weight $`\rho `$. The Casimir operator on $`V_\rho `$ is defined by
$$C_{\pi _\rho }=\frac{1}{32}\underset{i<j}{}\pi _\rho ([e_i,e_j])\pi _\rho ([e_i,e_j])=\frac{1}{64}\underset{i,j}{}\pi _\rho ([e_i,e_j])\pi _\rho ([e_i,e_j]).$$
(1.7)
This operator commutes with the action of $`Spin(n)`$. So the Casimir operator is a constant $`c(\rho )`$ on $`V_\rho `$:
$$c(\rho ):=\frac{1}{2}(\delta +\rho ^2\delta ^2).$$
(1.8)
Here the inner products on the weight space is the standard one, that is, $`\rho ^2=\rho ,\rho :=\rho ^k\rho ^k`$, and $`\delta `$ is half the sum of the positive roots.
We consider the tensor representation $`(\pi _\rho \pi _{\mathrm{Ad}},V_\rho 𝐑^n)(\pi _\rho \pi _{\mathrm{Ad}},V_\rho 𝐂^n)`$ and its irreducible decomposition $`\pi _\rho \pi _{\mathrm{Ad}}\pi _{\lambda _0}\pi _{\lambda _1}\mathrm{}\pi _{\lambda _N}`$. To identify the irreducible components, we need the Casimir operator or the operator $`\widehat{C}`$ given by
$$\widehat{C}:=C_{\pi _\rho \pi _{\mathrm{Ad}}}C_{\pi _\rho }11C_{\pi _{\mathrm{Ad}}}.$$
(1.9)
It is easy to show that
$$\widehat{C}=\frac{1}{32}\underset{i,j}{}\pi _\rho ([e_i,e_j])\pi _{\mathrm{Ad}}([e_i,e_j]).$$
(1.10)
This operator $`\widehat{C}`$ is a constant on the irreducible component $`V_{\lambda _k}`$. This constant is called the conformal weight for $`\lambda _k`$ (see ) and given by
$$\begin{array}{cc}\hfill m(\lambda _k):& =c(\lambda _k)c(\rho )c(\mathrm{Ad})\hfill \\ & =\frac{1}{2}(n\delta +\lambda _k^2+\delta +\rho ^21).\hfill \end{array}$$
(1.11)
## 2 Clifford homomorphisms
We shall discuss the Clifford homomorphisms. Before we define the Clifford homomorphisms, we recall the usual Clifford action on spinor spaces. We consider the complex spinor representation $`(\mathrm{\Pi }_\mathrm{\Delta },V_\mathrm{\Delta })`$ of $`Spin(n)`$. For $`n=2m`$, we remark that the spinor space $`V_\mathrm{\Delta }`$ splits to $`V_{\mathrm{\Delta }^+}V_\mathrm{\Delta }^{}`$ as a $`Spin(n)`$-module. We adopt the following definition of the Clifford action. When $`n`$ is $`2m+1`$, the representation space $`V_\mathrm{\Delta }𝐑^n`$ splits to $`V_TV_\mathrm{\Delta }`$ as a $`Spin(n)`$-module, where the highest weights $`\mathrm{\Delta }`$ and $`T`$ are $`((1/2)_m)`$ and $`(3/2,(1/2)_{m1})`$ respectively. If we have a vector $`\varphi u`$ in $`V_\mathrm{\Delta }𝐑^n`$, then we can project the vector onto $`V_\mathrm{\Delta }`$ along $`V_T`$ orthogonally. This projection gives the Clifford action of $`𝐑^n`$, $`u\varphi `$. This action satisfy the relation $`e_ie_j+e_je_i=2\delta _{ij}`$ and extends to the action of $`𝐂l_n`$. When $`n`$ is $`2m`$, the same discussion holds, so that we can define the Clifford action satisfying $`𝐑^n(V_{\mathrm{\Delta }^\pm })=V_\mathrm{\Delta }^{}`$. We apply this definition to other representations.
Let us consider an irreducible unitary representation $`(\pi _\rho ,V_\rho )`$ and the tensor representation $`\pi _\rho \pi _{\mathrm{Ad}}\pi _{\lambda _0}\mathrm{}\pi _{\lambda _N}`$. Since all weights of the adjoint representation $`\pi _{\mathrm{Ad}}`$ have multiplicity one, each irreducible constitute of $`(\pi _\rho \pi _{\mathrm{Ad}},V_\rho 𝐑^n)`$ has multiplicity one. In fact, Fegan shows the following fact in .
###### Lemma 2.1 ().
If $`V_\rho 𝐑^n=V_\lambda `$ is the irreducible decomposition as a $`Spin(n)`$-module (or a $`𝔰𝔭𝔦𝔫(n)`$-module), then the highest weight of irreducible components is given as follows:
1. When $`n`$ is $`2m`$, $`\lambda `$ is dominant weight and $`\lambda =\rho \pm \mu _i`$ $`(1im)`$.
2. When $`n`$ is $`2m+1`$ and $`\rho ^m\frac{1}{2}`$, $`\lambda `$ is dominant weight, and $`\lambda =\rho `$ or $`\lambda =\rho \pm \mu _i`$ $`(1im)`$.
3. When $`n`$ is $`2m+1`$ and $`\rho ^m`$ is $`0`$, $`\lambda `$ is dominant integral, and $`\lambda =\rho +\mu _m`$ or $`\lambda =\rho \pm \mu _i`$ $`(1im1)`$.
Here $`\rho =(\rho ^1,\rho ^2,\mathrm{},\rho ^m)`$ and $`\mu _i=(0_{i1},1,0_{mi})`$. The conformal weight of irreducible components is given as follows:
$$m(\rho +\mu _i)=i1\rho ^i,m(\rho \mu _i)=n+\rho ^ii1,m(\rho )=\frac{1}{2}(n1).$$
(2.1)
We remark that the highest weight $`\rho +\mu _1`$ called top term certainly occurs once in the decomposition. If we employ the lexicographical order for the weight space, then we arrange the highest weight $`\{\lambda _k\}_{0kN}`$ to satisfy $`\rho +\mu _1=\lambda _0>\lambda _1>\mathrm{}>\lambda _N`$. Lemma 2.1 implies that their conformal weights reverse as $`\rho ^1=m(\lambda _0)<m(\lambda _1)<\mathrm{}<m(\lambda _N)`$ except the following case: when $`n=2m`$, $`\rho ^{m1}1`$, and $`\rho ^m=0`$, we have $`m(\rho +\mu _m)=m(\rho \mu _m)`$. This case is said to be the exceptional case.
The inner products $`,`$ on $`V_\rho `$ and $`𝐑^n`$ give the one on $`V_\rho 𝐑^n`$ where irreducible components are orthogonal to each other. Then we have the orthogonal projection $`\mathrm{\Pi }_{\lambda _k}^\rho `$ onto irreducible component $`V_{\lambda _k}`$ from $`V_\rho 𝐑^n`$. This projection induces a homomorphism from $`V_\rho `$ to $`V_{\lambda _k}`$.
###### Definition 2.1.
Let $`V_\rho 𝐑^n=V_{\lambda _k}`$ be the irreducible decomposition as a $`Spin(n)`$-module. For each irreducible component $`V_{\lambda _k}`$, we have the bilinear mapping
$$𝐑^n\times V_\rho (u,\varphi )p_{\lambda _k}^\rho (u)\varphi :=\mathrm{\Pi }_{\lambda _k}^\rho (\varphi u)V_{\lambda _k}.$$
(2.2)
We call the linear mapping $`p_{\lambda _k}^\rho (u):V_\rho V_{\lambda _k}`$ for all $`u`$ in $`𝐑^n`$, the Clifford homomorphism from $`V_\rho `$ to $`V_{\lambda _k}`$.
We denote by $`(p_{\lambda _k}^\rho (u))^{}`$ the adjoint operator of $`p_{\lambda _k}^\rho (u)`$ with respect to the inner products on $`V_\rho `$ and $`V_{\lambda _k}`$. If we consider the tensor representation $`V_{\lambda _k}𝐑^n=V_\nu `$, then we find the irreducible component $`V_\rho `$ and the Clifford homomorphism $`p_\rho ^{\lambda _k}(u)`$ is $`(p_{\lambda _k}^\rho (u))^{}`$ up to normalization.
From now on, we shall investigate the Clifford homomorphisms. Since the Clifford homomorphisms is defined by a projection mapping, we can easily notice the following fact.
###### Proposition 2.2.
For $`u,v`$ in $`𝐑^n`$, the Clifford homomorphisms $`(p_{\lambda _k}^\rho (u))^{}p_{\lambda _k}^\rho (v)`$ satisfy that
$$\underset{k}{}(p_{\lambda _k}^\rho (u))^{}p_{\lambda _k}^\rho (v)=u,v\text{on }V_\rho .$$
(2.3)
###### Proof.
For $`\varphi v`$ and $`\psi v`$ in $`V_\rho 𝐑^n`$, we have
$$\varphi v,\psi u=\varphi ,\psi u,v=u,v\varphi ,\psi .$$
On the other hand, we have
$$\varphi v,\psi u=\underset{k}{}p_{\lambda _k}^\rho (v)\varphi ,p_{\lambda _k}^\rho (u)\psi =\underset{k}{}(p_{\lambda _k}^\rho (u))^{}p_{\lambda _k}^\rho (v)\varphi ,\psi .$$
Since these equations hold for all $`\psi `$ in $`V_\rho `$, we have proved the proposition. ∎
Next, we show an important lemma to give relations among the Clifford homomorphisms.
###### Lemma 2.3.
The Clifford homomorphism $`p_{\lambda _k}^\rho (u)`$ from $`V_\rho `$ to $`V_{\lambda _k}`$ satisfies the relation
$$\frac{1}{4}\underset{i}{}p_{\lambda _k}^\rho (e_i)\pi _\rho ([e_i,u])=m(\lambda _k)p_{\lambda _k}^\rho (u),$$
(2.4)
where $`\{e_i\}_{i=1}^n`$ is any orthonormal basis of $`𝐑^n`$.
###### Proof.
To prove the lemma, we use the operator $`\widehat{C}`$ given in section 1. For $`\varphi u`$ in $`V_\rho 𝐑^n`$, we have
$$\begin{array}{cc}\hfill \widehat{C}(\varphi u)& =\frac{1}{32}\underset{s,t,i}{}\pi _\rho ([e_s,e_t])\varphi [[e_s,e_t],e_i]e_i,u\hfill \\ & =\frac{1}{8}\underset{s,t}{}\pi _\rho ([e_s,e_t])\varphi (\delta _{si}e_t\delta _{ti}e_s)e_i,u\hfill \\ & =\frac{1}{4}\underset{l}{}\pi _\rho ([u,e_t])\varphi e_t\hfill \\ & =\frac{1}{4}\underset{k,t}{}p_{\lambda _k}^\rho (e_t)\pi _\rho ([u,e_t])\varphi .\hfill \end{array}$$
The operator $`\widehat{C}`$ is the constant $`m(\lambda _k)`$ on $`V_{\lambda _k}`$, so that $`\widehat{C}(\varphi u)=_km(\lambda _k)p_{\lambda _k}^\rho (u)\varphi `$. Hence, we conclude that
$$m(\lambda _k)p_{\lambda _k}^\rho (u)=\frac{1}{4}\underset{t}{}p_{\lambda _k}^\rho (e_t)\pi _\rho ([u,e_t]).$$
This lemma leads a relation among the Clifford homomorphisms.
###### Proposition 2.4.
We sum up the Clifford homomorphisms $`(p_{\lambda _k}^\rho (u))^{}p_{\lambda _k}^\rho (v)`$ for all $`\lambda _k`$ with its conformal weight. Then we have
$$\underset{k}{}m(\lambda _k)(p_{\lambda _k}^\rho (u))^{}p_{\lambda _k}^\rho (v)=\frac{1}{4}\pi _\rho ([u,v]).$$
(2.5)
###### Proof.
We substitute (2.4) into $`_km(\lambda _k)(p_{\lambda _k}^\rho (u))^{}p_\lambda ^\rho (v)`$ and use the relation (2.3). ∎
As this proposition, we give many relations among the Clifford homomorphisms.
###### Theorem 2.5.
For an irreducible representation $`(\pi _\rho ,V_\rho )`$ and any non-negative integer $`q`$, we define the bilinear mapping $`r_\rho ^q`$ from $`𝐑^n\times 𝐑^n`$ to $`\mathrm{End}(V_\rho )`$ as follows:
$$\begin{array}{c}r_\rho ^q:𝐑^n\times 𝐑^n(u,v)\hfill \\ \hfill \left(\frac{1}{4}\right)^q\underset{l_1,\mathrm{},l_{q1}}{}\pi _\rho ([u,e_{l_1}])\pi _\rho ([e_{l_1},e_{l_2}])\mathrm{}\pi _\rho ([e_{l_{q1}},v])\mathrm{End}(V_\rho ),\end{array}$$
(2.6)
and $`r_\rho ^0(u,v):=u,v`$. Then we have relations among the Clifford homomorphisms for any $`q`$:
$$\underset{k}{}m(\lambda _k)^q(p_{\lambda _k}^\rho (u))^{}p_{\lambda _k}^\rho (v)=r_\rho ^q(u,v).$$
(2.7)
As a corollary of this theorem, we give a description of general Casimir operators. Let $`U(𝔰𝔭𝔦𝔫(n))`$ be the enveloping algebra of $`𝔰𝔭𝔦𝔫(n)`$. Then the general Casimir operator is said to be an element of the center of $`U(𝔰𝔭𝔦𝔫(n))`$ or its image by $`\pi _\rho `$ for the representation $`(\pi _\rho ,V_\rho )`$. In , we know generators of the general Casimir operators and its action on the irreducible representations. The generators are given by
$$\left(\frac{1}{4}\right)^q\underset{l_1,\mathrm{},l_q}{}\pi _\rho ([e_{l_1},e_{l_2}])\pi _\rho ([e_{l_2},e_{l_3}])\mathrm{}\pi _\rho ([e_{l_q},e_{l_1}]).$$
(2.8)
###### Corollary 2.6.
The Clifford homomorphisms $`_{k,i}m(\lambda _k)^q(p_{\lambda _k}^\rho (e_i))^{}p_{\lambda _k}^\rho (e_i)`$ for $`q0`$ generate the general Casimir operators for $`(\pi _\rho ,V_\rho )`$ except the Pfaffian type Casimir operator in the remark below.
###### Remark 2.1.
For $`n=2m`$, we need the Pfaffian type Casimir operator to generate the center of $`U(𝔰𝔭𝔦𝔫(2m))`$, that is,
$$\mathrm{Pf}_\rho :=\underset{\sigma S_{2m}}{}\mathrm{sig}(\sigma )\pi _\rho ([e_{\sigma (1)},e_{\sigma (2)}])\pi _\rho ([e_{\sigma (3)},e_{\sigma (4)}])\mathrm{}\pi _\rho ([e_{\sigma (2m1)},e_{\sigma (2m)}]).$$
(2.9)
We can show a relation between this operator and the Clifford homomorphisms,
$$\begin{array}{c}\underset{k}{}p(\lambda _k)(p_{\lambda _k}^\rho (e_j))^{}p_{\lambda _k}^\rho (e_i)=\delta _{ij}p(\rho )+\hfill \\ \hfill 8m(1\delta _{ij})\mathrm{sgn}\left(\begin{array}{cccc}1& 2& \mathrm{}& 2m\\ i& j& \mathrm{}& 2m\end{array}\right)\underset{\sigma \stackrel{~}{S}_{2m}}{}\mathrm{sgn}(\sigma )\pi _\rho ([e_{\sigma (1)},e_{\sigma (2)}])\mathrm{}\pi _\rho ([e_{\sigma (2m1)},e_{\sigma (2m)}]),\end{array}$$
(2.10)
where $`\stackrel{~}{S}_{2m}`$ is the permutation of $`\{1,\mathrm{},2m\}\{i,j\}`$ and $`p(\rho )`$ is a scalar action of $`\mathrm{Pf}_\rho `$ on $`V_\rho `$ as
$$p(\rho )=(\sqrt{1})^m8^mm!(\rho ^1+m1)(\rho ^2+m2)\mathrm{}(\rho ^{m1}+1)\rho ^m.$$
(2.11)
If we define the bilinear mapping $`\mathrm{pf}_\rho (,)`$ from $`𝐑^n\times 𝐑^n`$ to $`\mathrm{End}(V_\rho )`$ by the right hand side of (2.10), then the trace of $`\mathrm{pf}_\rho (,)`$ is $`n\mathrm{Pf}_\rho `$.
Now, we know that the usual Clifford action satisfies that
$$(gug^1)\varphi =\pi _\mathrm{\Delta }(g)(u(\pi _\mathrm{\Delta }(g^1)\varphi )),$$
(2.12)
where $`\varphi `$ is in $`V_\mathrm{\Delta }`$ and $`g`$ in $`Spin(n)`$ (see ). We generalize this relation.
###### Proposition 2.7.
Let $`u`$ be in $`𝐑^n`$ and $`g`$ be in $`Spin(n)`$. The Clifford homomorphism $`p_{\lambda _k}^\rho (u)`$ satisfy that
$$p_{\lambda _k}^\rho (gug^1)=\pi _{\lambda _k}(g)p_{\lambda _k}^\rho (u)\pi _\rho (g^1).$$
(2.13)
###### Proof.
For $`\varphi `$ in $`V_\rho `$, we have
$$\begin{array}{cc}\hfill p_{\lambda _k}^\rho (gug^1)\varphi & =\mathrm{\Pi }_{\lambda _k}^\rho (\pi _\rho (g)\pi _\rho (g^1)\varphi \pi _{\mathrm{Ad}}(g)u)\hfill \\ & =\mathrm{\Pi }_{\lambda _k}^\rho (\pi _{\rho \mathrm{Ad}}(g)(\pi _\rho (g^1)\varphi u))\hfill \\ & =\pi _{\lambda _k}(g)\mathrm{\Pi }_{\lambda _k}^\rho (\pi _\rho (g^1)\varphi u)\hfill \\ & =\pi _{\lambda _k}(g)p_{\lambda _k}^\rho (u)\pi _\rho (g^1)\varphi .\hfill \end{array}$$
As a corollary of Proposition 2.7, we have a formula of the projection $`\mathrm{\Pi }_{\lambda _k}^\rho `$ by using the Clifford homomorphisms.
###### Corollary 2.8.
The orthogonal projection $`\mathrm{\Pi }_{\lambda _k}^\rho :V_\rho 𝐑^nV_{\lambda _k}`$ is realized as follows:
$$\mathrm{\Pi }_{\lambda _k}^\rho (\varphi u)=\underset{i}{}(p_{\lambda _k}^\rho (e_i))^{}p_{\lambda _k}^\rho (u)(\varphi )e_i.$$
(2.14)
###### Proof.
Let $`V_{\lambda _k}`$ be an irreducible component of $`V_\rho 𝐑^n`$. We consider the following embedding from $`V_{\lambda _k}`$ to $`V_\rho 𝐑^n`$, which dose not always preserve their inner products:
$$i_{\lambda _k}:V_{\lambda _k}\psi \underset{i}{}(p_{\lambda _k}^\rho (e_i))^{}(\psi )e_iV_\rho 𝐑^n.$$
The map $`i_{\lambda _k}`$ is independent of the orthonormal basis of $`𝐑^n`$ which we chose, and commutes with the action of $`Spin(n)`$ . Since the irreducible component of $`V_\rho 𝐑^n`$ has multiplicity one and $`i_{\lambda _k}`$ is not zero, we prove that the map $`i_{\lambda _k}`$ is a well-defined embedding.
For $`\varphi u`$ in $`V_\rho 𝐑^n`$, we have
$$\begin{array}{cc}\hfill \varphi u& =\underset{i}{}u,e_i\varphi e_i\hfill \\ & =\underset{k,i}{}(p_{\lambda _k}^\rho (e_i))^{}p_{\lambda _k}^\rho (u)(\varphi )e_i\text{( from Proposition }\text{2.2}\text{ )}\hfill \\ & =\underset{k}{}i_{\lambda _k}(p_{\lambda _k}^\rho (u)(\varphi )).\hfill \end{array}$$
This completes the proof of Corollary 2.8. ∎
This corollary and Theorem 2.5 allow us to have an explicit formula of the projection $`\mathrm{\Pi }_{\lambda _k}^\rho `$. We define the $`(N+1)\times (N+1)`$ matrix $`𝐌=(m_{ij})_{0i,jN}`$ by $`m_{ij}:=m(\lambda _j)^i`$, which is a Vandermonde matrix. Since the conformal weights are different from each other except the exceptional case, the inverse matrix of $`𝐌`$ exists. From (2.7) and (2.14), we have a formula of $`\mathrm{\Pi }_{\lambda _k}^\rho `$,
$$\mathrm{\Pi }_{\lambda _k}^\rho (\varphi u)=\underset{i,q}{}n_{kq}r_\rho ^q(e_i,u)\varphi e_i.$$
(2.15)
Here, $`n_{ij}`$ is the $`(i,j)`$-component of $`𝐌^1`$. If we set $`S_j(x_0,\mathrm{},\widehat{x_i},\mathrm{},x_N)`$ as the $`j`$-th fundamental symmetric polynomial of $`\{x_0,x_1,\mathrm{},x_N\}\{x_i\}`$, then
$$n_{ij}=(1)^{Nj}\frac{S_{Nj}(m(\lambda _0),\mathrm{},\widehat{m(\lambda _i)},\mathrm{},m(\lambda _N))}{_{ki}(m(\lambda _i)m(\lambda _k))}.$$
(2.16)
In the exceptional case that $`n`$ is even, $`\rho ^m1`$, and $`\rho ^m=0`$, the conformal weight of $`\lambda _+:=\rho +\mu _m`$ coincides with the one of $`\lambda _{}:=\rho \mu _m`$. So we use $`\mathrm{pf}_\rho `$ given in (2.10) to obtain a formula of $`\mathrm{\Pi }_{\lambda _\pm }^\rho `$. It is easy to show that
$$p(\lambda _+)(p_{\lambda _+}^\rho (e_j))^{}p_{\lambda _+}^\rho (e_i)+p(\lambda _{})(p_\lambda _{}^\rho (e_j))^{}p_\lambda _{}^\rho (e_i)=\mathrm{pf}_\rho (e_j,e_i).$$
(2.17)
It follows that we realize the projection of $`\mathrm{\Pi }_{\lambda _\pm }^\rho `$ by using $`r_\rho ^q`$ and $`\mathrm{pf}_\rho `$.
We shall state some examples of the Clifford homomorphisms.
###### Example 2.1 (Spinor).
We discuss only the case of $`n=2m+1`$ and leave the case of $`n=2m`$ to the reader. Let $`V_\mathrm{\Delta }`$ be spinor space which is an irreducible representation space with highest weight $`((1/2)_m)`$. It follows from Lemma 2.1 that the vector space $`V_\mathrm{\Delta }𝐑^n`$ is isomorphic to $`V_\mathrm{\Delta }V_T`$. Here, $`V_T`$ is the representation space with highest weight $`(3/2,(1/2)_{m1})`$. From Theorem 2.5, we have
$$(p_T^\mathrm{\Delta }(u))^{}p_T^\mathrm{\Delta }(v)+(p_\mathrm{\Delta }^\mathrm{\Delta }(u))^{}p_\mathrm{\Delta }^\mathrm{\Delta }(v)=u,v,$$
(2.18)
$$\frac{1}{2}(p_T^\mathrm{\Delta }(u))^{}p_T^\mathrm{\Delta }(v)+\frac{n1}{2}(p_\mathrm{\Delta }^\mathrm{\Delta }(u))^{}p_\mathrm{\Delta }^\mathrm{\Delta }(v)=\frac{1}{4}\pi _\mathrm{\Delta }([u,v]),$$
(2.19)
for $`u`$ and $`v`$ in $`𝐑^n`$. It follows that we give the usual Clifford relation
$$(p_\mathrm{\Delta }^\mathrm{\Delta }(u))^{}p_\mathrm{\Delta }^\mathrm{\Delta }(v)+(p_\mathrm{\Delta }^\mathrm{\Delta }(v))^{}p_\mathrm{\Delta }^\mathrm{\Delta }(u)=\frac{2}{n}u,v.$$
(2.20)
Since $`p_\mathrm{\Delta }^\mathrm{\Delta }`$ is a constant multiple of the Clifford action $`u`$ satisfying $`u\varphi ,\psi =\varphi ,u\psi `$, we have $`p_\mathrm{\Delta }^\mathrm{\Delta }(u)=\frac{1}{\sqrt{n}}u`$ and $`(p_\mathrm{\Delta }^\mathrm{\Delta }(u))^{}=\frac{1}{\sqrt{n}}u`$. Furthermore the representation of $`𝔰𝔭𝔦𝔫(n)`$ on spinor space is realized as $`\pi _\mathrm{\Delta }([e_i,e_j])=[e_i,e_j]`$. The orthogonal projections $`\mathrm{\Pi }_\mathrm{\Delta }^\mathrm{\Delta }`$ and $`\mathrm{\Pi }_T^\mathrm{\Delta }`$ are
$$\begin{array}{cc}\hfill \mathrm{\Pi }_\mathrm{\Delta }^\mathrm{\Delta }(\varphi u)& =\{\frac{1}{n}e_i,u\frac{1}{2n}\pi _\mathrm{\Delta }([e_i,u])\}(\varphi )e_i\hfill \\ & =\frac{1}{n}e_iu\varphi e_i,\hfill \end{array}$$
(2.21)
$$\begin{array}{cc}\hfill \mathrm{\Pi }_T^\mathrm{\Delta }(\varphi u)& =\{\frac{n1}{n}e_i,u+\frac{1}{2n}\pi _\mathrm{\Delta }([e_i,u])\}(\varphi )e_i\hfill \\ & =\varphi u+\frac{1}{n}e_iu\varphi e_i.\hfill \end{array}$$
(2.22)
###### Example 2.2 (Exterior algebra).
We shall discuss the exterior tensor product representation $`(\pi _{\mathrm{\Lambda }^k},\mathrm{\Lambda }^k(𝐑^n))`$ with highest weight $`\rho =(1_k,0_{mk})`$ for $`1km`$. There are three irreducible representations $`\lambda _0=(2,1_{k1},0_{mk})`$, $`\lambda _1=(1_{k+1},0_{mk1})`$, and $`\lambda _2=(1_{k1},0_{mk+1})`$ in $`\mathrm{\Lambda }^k(𝐑^n)𝐑^n`$. We realize the Clifford homomorphisms as follows:
$`(p_{\lambda _0}^\rho (e_j))^{}p_{\lambda _0}^\rho (e_i)`$ $`={\displaystyle \frac{1}{(k+1)(nk+1)}}(k(nk)\delta _{ij}nr_\rho ^1(e_j,e_i)+r_\rho ^2(e_j,e_i)),`$ (2.23)
$`(p_{\lambda _1}^\rho (e_j))^{}p_{\lambda _1}^\rho (e_i)`$ $`={\displaystyle \frac{1}{(k+1)(n2k)}}((nk)\delta _{ij}+(nk1)r_\rho ^1(e_j,e_i)r_\rho ^2(e_j,e_i)),`$ (2.24)
$`(p_{\lambda _2}^\rho (e_j))^{}p_{\lambda _2}^\rho (e_i)`$ $`={\displaystyle \frac{1}{(n2k)(nk+1)}}(k\delta _{ij}(k1)r_\rho ^1(e_j,e_i)+r_\rho ^2(e_j,e_i)).`$ (2.25)
In general, it holds that $`r_\rho ^2(u,v)r_\rho ^2(v,u)=(n2)r_\rho ^1(u,v)`$. Then we have
$$(k+1)(p_{\lambda _1}^\rho (e_j))^{}p_{\lambda _1}^\rho (e_i)+(nk+1)(p_{\lambda _2}^\rho (e_i))^{}p_{\lambda _2}^\rho (e_j)=\delta _{ij},$$
(2.26)
and
$$\begin{array}{cc}& (k+1)\{(p_{\lambda _1}^\rho (e_j))^{}p_{\lambda _1}^\rho (e_i)(p_{\lambda _1}^\rho (e_i))^{}p_{\lambda _1}^\rho (e_j)\}\hfill \\ \hfill =& (nk+1)\{(p_{\lambda _2}^\rho (e_j))^{}p_{\lambda _2}^\rho (e_i)(p_{\lambda _2}^\rho (e_i))^{}p_{\lambda _2}^\rho (e_j)\}\hfill \\ \hfill =& r_\rho ^1(e_j,e_i)=\frac{1}{4}\pi _\rho ([e_j,e_i]).\hfill \end{array}$$
(2.27)
These equations correspond to the following relations, respectively:
$$i(e_j)e_i+e_ji(e_i)=\delta _{ij},$$
(2.28)
$$i(e_i)e_ji(e_i)e_j=e_ji(e_i)e_ii(e_j)=\frac{1}{4}\pi _\rho ([e_j,e_i]),$$
(2.29)
where $`i(u)`$ is the interior product of $`u`$. So we put
$$p_{\lambda _1}^\rho (e_i):=\frac{1}{\sqrt{k+1}}e_i,p_{\lambda _2}^\rho (e_i):=\frac{1}{\sqrt{nk+1}}i(e_i).$$
(2.30)
Then the projections are realized as
$$\mathrm{\Pi }_{\lambda _1}^\rho (\varphi u)=\frac{1}{k+1}i(e_i)u_{}\varphi e_i,$$
(2.31)
$$\mathrm{\Pi }_{\lambda _2}^\rho (\varphi u)=\frac{1}{nk+1}e_ii(u)\varphi e_i,$$
(2.32)
$$\mathrm{\Pi }_{\lambda _0}^\rho (\varphi u)=\varphi u\mathrm{\Pi }_1^\rho (\varphi u)\mathrm{\Pi }_2^\rho (\varphi u).$$
(2.33)
If we consider the (exceptional) case that $`n`$ is $`2m`$ and the highest weight $`\rho `$ is $`(1_{m1},0)`$, then we have four irreducible representations, $`\lambda _0=(2,1_{m1},0)`$, $`(\lambda _1)_+=(1_m)`$, $`(\lambda _1)_{}=(1_{m1},1)`$, and $`\lambda _2=(1_{m2},0_2)`$ in $`\mathrm{\Lambda }^{m1}(𝐑^{2m})𝐑^{2m}`$. We denote the Hodge star operator by the asterisk $``$ and have
$$\mathrm{\Pi }_{(\lambda _1)_+}^\rho (\varphi u)=\frac{1}{m}i(e_i)\frac{1}{2}(1+)u_{}\varphi e_i,$$
(2.34)
$$\mathrm{\Pi }_{(\lambda _1)_{}}^\rho (\varphi u)=\frac{1}{m}i(e_i)\frac{1}{2}(1)u_{}\varphi e_i.$$
(2.35)
## 3 Higher spin Dirac operators
In this section we shall discuss the higher spin Dirac operators. Let $`M`$ be a $`n`$-dimensional spin manifold and $`\mathrm{𝐒𝐩𝐢𝐧}(M)`$ be its spin structure (If $`M`$ dose not have a spin structure, we should consider only the representations of $`SO(n)`$). For an irreducible unitary representation $`(\pi _\rho ,V_\rho )`$ of the structure group $`Spin(n)`$, we have a vector bundle $`𝐒_\rho `$ associated to $`\mathrm{𝐒𝐩𝐢𝐧}(M)`$,
$$𝐒_\rho :=\mathrm{𝐒𝐩𝐢𝐧}(M)\times _\rho V_\rho .$$
(3.1)
We introduce a covariant derivative on $`\mathrm{\Gamma }(M,𝐒_\rho )`$ associated to the spin connection or the Levi-Civita connection,
$$:\mathrm{\Gamma }(M,𝐒_\rho )\mathrm{\Gamma }(M,𝐒_\rho T^{}(M)).$$
(3.2)
If we decompose $`𝐒_\rho T^{}(M)`$ into the direct sum of irreducible components, $`𝐒_{\lambda _k}`$, then we construct a first order differential operator from $`\mathrm{\Gamma }(M,𝐒_\rho )`$ to $`\mathrm{\Gamma }(M,𝐒_{\lambda _k})`$ called the higher spin Dirac operator,
$$D_{\lambda _k}^\rho :\mathrm{\Gamma }(M,𝐒_\rho )\stackrel{}{}\mathrm{\Gamma }(M,𝐒_\rho T^{}(M))\stackrel{\mathrm{\Pi }_{\lambda _k}^\rho }{}\mathrm{\Gamma }(M,𝐒_{\lambda _k}).$$
(3.3)
From Proposition 2.7, we show that the bundle homomorphism $`p_{\lambda _k}^\rho (X)`$ from $`𝐒_\rho `$ to $`𝐒_{\lambda _k}`$ is well-defined, where $`X`$ is in $`\mathrm{\Gamma }(M,T(M))`$. Then we have a formula of $`D_{\lambda _k}^\rho `$,
$$D_{\lambda _k}^\rho =p_{\lambda _k}^\rho (e_i)_{e_i}:\mathrm{\Gamma }(M,𝐒_\rho )\mathrm{\Gamma }(M,𝐒_{\lambda _k}).$$
(3.4)
Here, $`\{e_i\}`$ is an orthonormal frame of $`T(M)T^{}(M)`$. We set a formal adjoint operator of $`D_{\lambda _k}^\rho `$ by $`(D_{\lambda _k}^\rho )^{}:=(p_{\lambda _k}^\rho (e_i))^{}_{e_i}`$.
We shall extend the relations among the Clifford homomorphisms to the ones among higher spin Dirac operators. First, we show that the sum of the higher spin Dirac operator $`(D_{\lambda _k}^\rho )^{}D_{\lambda _k}^\rho `$ for $`0kN`$ is the connection Laplacian.
###### Proposition 3.1.
Let $`^{}`$ is the connection Laplacian on $`\mathrm{\Gamma }(M,𝐒_\rho )`$ defined by $`_{e_i}_{e_i}`$, and $`𝐒_{\lambda _k}`$ be the irreducible decomposition of $`𝐒_\rho T^{}(M)`$. Then
$$^{}=\underset{k}{}(D_{\lambda _k}^\rho )^{}D_{\lambda _k}^\rho .$$
(3.5)
###### Proof.
For any $`x`$ in $`M`$, we can find a local orthonormal frame $`\{e_i\}`$ such that $`(_{e_i}e_j)_x=0`$ for any $`i`$ and $`j`$, so that $`(_{e_i}p_{\lambda _k}^\rho (e_j))_x=0`$. Then we have
$$\begin{array}{cc}\hfill \underset{k}{}(D_{\lambda _k}^\rho )^{}D_{\lambda _k}^\rho & =\underset{i,j}{}\underset{k}{}(p_{\lambda _k}^\rho (e_j))^{}p_{\lambda _k}^\rho (e_i)_{e_j}_{e_i}\hfill \\ & =\underset{i,j}{}\delta _{ij}_{e_j}_{e_i}\hfill \\ & =\underset{i}{}_{e_i}_{e_i}.\hfill \end{array}$$
Next, we introduce a curvature transformation on $`𝐒_\rho `$. For $`X`$ and $`Y`$ in $`\mathrm{\Gamma }(M,T(M))`$, the curvature $`R_\rho (X,Y)`$ is defined by
$$R_\rho (X,Y)=_X_Y_Y_X_{[X,Y]}\mathrm{\Gamma }(M,\mathrm{End}(𝐒_\rho )).$$
(3.6)
This curvature leads the curvature transformation $`R_\rho ^1`$ in $`\mathrm{\Gamma }(M,\mathrm{End}(𝐒_\rho ))`$,
$$R_\rho ^1:=\frac{1}{8}\underset{i,j}{}\pi _\rho ([e_i,e_j])R_\rho (e_i,e_j).$$
(3.7)
For example, $`R_\mathrm{\Delta }^1`$ is $`\kappa /8`$ and $`R_{\mathrm{\Lambda }^1}^1`$ is $`\mathrm{Ric}`$, where $`\kappa `$ is the scalar curvature and $`\mathrm{Ric}`$ is the Ricci curvature transformation. We show that the sum of $`(D_{\lambda _k}^\rho )^{}D_{\lambda _k}^\rho `$ with conformal weight is the curvature transformation.
###### Proposition 3.2.
Let $`R_\rho ^1`$ be the curvature transformation on $`𝐒_\rho `$ as above. Then we have
$$R_\rho ^1=\underset{k}{}m(\lambda _k)(D_{\lambda _k}^\rho )^{}D_{\lambda _k}^\rho .$$
(3.8)
Now, we define a family of self adjoint differential operators on $`\mathrm{\Gamma }(M,𝐒_\rho )`$ by $`R_\rho ^q:=_{i,j}r_\rho ^q(e_i,e_j)_i_j`$ for $`q0`$. If $`q`$ is even and $`M`$ is compact, then $`R_\rho ^q`$ is a non-negative operator and $`\mathrm{ker}R_\rho ^q`$ is $`\mathrm{ker}`$, that is, the space of parallel sections.
###### Theorem 3.3.
The sum of higher spin Dirac operators $`(D_{\lambda _k}^\rho )^{}D_{\lambda _k}^\rho `$ with the $`q`$-th power of conformal weight is $`R_\rho ^q`$:
$$R_\rho ^q=\underset{k}{}m(\lambda _k)^q(D_{\lambda _k}^\rho )^{}D_{\lambda _k}^\rho .$$
(3.9)
This implies an explicit formula of $`(D_{\lambda _k}^\rho )^{}D_{\lambda _k}^\rho `$ as follows:
$$(D_{\lambda _k}^\rho )^{}D_{\lambda _k}^\rho =\underset{q}{}n_{kq}R_\rho ^q,$$
(3.10)
where $`n_{ij}`$ is a constant given in (2.16) and, for the exceptional case, we think of the direct sum of $`𝐒_{\rho +\mu _m}`$ and $`𝐒_{\rho \mu _m}`$ as only one component of $`𝐒_{\lambda _k}`$.
###### Remark 3.1.
We can define a differential operator $`\mathrm{P}_\rho `$ by $`\mathrm{P}_\rho :=\mathrm{pf}_\rho (e_i,e_j)_{e_i}_{e_j}`$. Since $`\mathrm{pf}_\rho (e_i,e_j)=\mathrm{pf}_\rho (e_j,e_i)`$ for the exceptional case, the operator $`P_\rho `$ is a curvature transformation.
###### Example 3.1 (Spinor).
We have two higher spin Dirac operators on spinor bundle $`𝐒_\mathrm{\Delta }`$, that is, the Dirac operator $`D`$ and the twistor operator $`T`$. Our higher spin Dirac operators $`D_\mathrm{\Delta }^\mathrm{\Delta }:\mathrm{\Gamma }(M,𝐒_\mathrm{\Delta })\mathrm{\Gamma }(M,𝐒_\mathrm{\Delta })`$ and $`D_T^\mathrm{\Delta }:\mathrm{\Gamma }(M,𝐒_\mathrm{\Delta })\mathrm{\Gamma }(M,𝐒_T)`$ satisfy that
$$(D_\mathrm{\Delta }^\mathrm{\Delta })^{}D_\mathrm{\Delta }^\mathrm{\Delta }+(D_T^\mathrm{\Delta })^{}D_T^\mathrm{\Delta }=^{},$$
(3.11)
$$\frac{1}{2}(D_T^\mathrm{\Delta })^{}D_T^\mathrm{\Delta }+\frac{n1}{2}(D_\mathrm{\Delta }^\mathrm{\Delta })^{}D_\mathrm{\Delta }^\mathrm{\Delta }=R_\mathrm{\Delta }^1=\frac{1}{8}\kappa .$$
(3.12)
So we set $`D_\mathrm{\Delta }^\mathrm{\Delta }:=\frac{1}{\sqrt{n}}D`$ and $`D_T^\mathrm{\Delta }:=\frac{\sqrt{n1}}{\sqrt{n}}T`$. Then we have Bochner identities,
$$D^2=^{}+\frac{1}{4}\kappa ,T^{}T=^{}\frac{1}{4(n1)}\kappa .$$
(3.13)
It is from the projection formula (2.22) that twistor operator is realized as follows (see and ):
$$D_T^\mathrm{\Delta }(\varphi )=\underset{i}{}p_T^\mathrm{\Delta }(e_i)_{e_i}\varphi =(_{e_i}\varphi +\frac{1}{n}e_iD\varphi )e_i.$$
(3.14)
Now, we assume that $`M`$ is compact and $`\lambda ^0`$ is the first eigenvalue of $`D^2`$. To give a lower bound of $`\lambda ^0`$, we rewrite the identity (3.12) as
$$D^2=\frac{n}{4(n1)}\kappa +\frac{n1}{2n}T^{}T.$$
(3.15)
This identity indicates that $`\lambda ^0\frac{n}{4(n1)}\kappa _0`$, where $`\kappa _0`$ is $`\mathrm{min}_{xM}\kappa (x)`$ (see and ).
###### Example 3.2 (Differential from).
We consider the $`k`$-th differential forms $`\mathrm{\Lambda }^k(M)`$. Then we obtain three higher spin Dirac operators,
$$d:=\sqrt{k+1}D_{\lambda _1}^\rho :\mathrm{\Gamma }(M,\mathrm{\Lambda }^k(M))\mathrm{\Gamma }(M,\mathrm{\Lambda }^{k+1}(M)),$$
(3.16)
$$d^{}:=\sqrt{nk+1}D_{\lambda _2}^\rho :\mathrm{\Gamma }(M,\mathrm{\Lambda }^k(M))\mathrm{\Gamma }(M,\mathrm{\Lambda }^{k1}(M)),$$
(3.17)
$$C:=D_{\lambda _0}^\rho :\mathrm{\Gamma }(M,\mathrm{\Lambda }^k(M))\mathrm{\Gamma }(M,𝐒_{\lambda _0}).$$
(3.18)
Here $`C`$ is called the conformal killing operator whose kernel is the space of conformal killing forms. It follows from Theorem 3.3 that
$$C^{}C+\frac{1}{k+1}d^{}d+\frac{1}{nk+1}dd^{}=^{},$$
(3.19)
$$C^{}C+\frac{k}{k+1}d^{}d+\frac{nk}{nk+1}dd^{}=R_\rho ^1,$$
(3.20)
$$C^{}C+\frac{k^2}{k+1}d^{}d+\frac{(nk)^2}{nk+1}dd^{}=R_\rho ^2.$$
(3.21)
The first and second identities give the Bochner identity for differential forms. The projection formula (2.33) leads a formula of $`C`$ (see ):
$$C(\varphi )=(_i\varphi \frac{1}{k+1}i(e_i)d\varphi +\frac{1}{nk+1}e_id^{}\varphi )e_i.$$
(3.22)
###### Example 3.3 (Higher spin fields on anti self-dual $`4`$-manifolds).
Let $`M`$ be an anti self-dual $`4`$-dimensional spin manifold and $`𝐒_{k,l}`$ be the vector bundle corresponding to the highest weight $`\rho :=(k+l,kl)`$ for $`k,l`$ in $`𝐙(𝐙+1/2)`$. We remark that $`𝐒_{k,l}`$ corresponds to the representation of $`\pi _k\widehat{}\pi _l`$ of $`Spin(4)=SU(2)\times SU(2)`$, where $`\pi _k`$ is a spin $`k`$ representation of $`SU(2)`$. We consider the vector bundle $`𝐒_{k,0}`$ and have two higher spin Dirac operators,
$$D_0:\mathrm{\Gamma }(M,𝐒_{k,0})\mathrm{\Gamma }(M,𝐒_{k+1/2,1/2}),$$
(3.23)
$$D_1:\mathrm{\Gamma }(M,𝐒_{k,0})\mathrm{\Gamma }(M,𝐒_{k1/2,1/2}).$$
(3.24)
The relations between $`D_0`$ and $`D_1`$ are
$$(D_0)^{}D_0+(D_1)^{}D_1=^{},$$
(3.25)
$$k(D_0)^{}D_0+(k+1)(D_1)^{}D_1=R_{k,0}^1.$$
(3.26)
Since $`M`$ is an anti self-dual manifold, we can show that the curvature transformation $`R_{k,0}^1`$ is $`\frac{k(k+1)}{6}\kappa `$. We set $`\widehat{D}_1:=\sqrt{(2k+1)/k}D_1`$ and give the Bochner identity for $`\widehat{D}_1`$ (see and ):
$$(\widehat{D}_1)^{}\widehat{D}_1=^{}+\frac{k+1}{6}\kappa .$$
(3.27)
The identity (3.26) gives a lower bound of the first eigenvalue $`\lambda ^0`$ of $`(\widehat{D}_1)^{}\widehat{D}_1`$, that is, $`\lambda ^0\frac{2k+1}{6}\kappa _0`$.
## 4 General Bochner identities
We have considered some Bochner identities in the previous section. It allows us to give a general Bochner identity on any associated bundle $`𝐒_\rho `$. We have known the following two identities:
$$\underset{0kN}{}(D_{\lambda _k}^\rho )^{}D_{\lambda _k}^\rho =^{},$$
(4.1)
$$\underset{0kN}{}m(\lambda _k)(D_{\lambda _k}^\rho )^{}D_{\lambda _k}^\rho =R_\rho ^1,$$
(4.2)
where we order $`\{\lambda _k\}`$ to satisfy that $`\lambda _0>\lambda _1>\mathrm{}>\lambda _N`$. We eliminate $`(D_{\lambda _0}^\rho )^{}D_{\lambda _0}^\rho `$ from above two equations and obtain a general Bochner identity on $`𝐒_\rho `$,
$$\underset{1kN}{}(1\frac{m(\lambda _k)}{m(\lambda _0)})(D_{\lambda _k}^\rho )^{}D_{\lambda _k}^\rho =^{}+\frac{1}{m(\lambda _0)}R_\rho ^1.$$
(4.3)
Here, $`m(\lambda ^0)`$ is $`\rho ^1`$ and $`(1\frac{m(\lambda _k)}{m(\lambda _0)})`$ is positive.
###### Theorem 4.1.
We define the differential operator $`\widehat{D}_{\lambda _k}^\rho `$ for $`1kN`$ by
$$\widehat{D}_{\lambda _k}^\rho :=\sqrt{1\frac{m(\lambda _k)}{m(\lambda _0)}}D_{\lambda _k}^\rho .$$
(4.4)
Then the following identity holds:
$$\mathrm{\Delta }_\rho :=\underset{k1}{}(\widehat{D}_{\lambda _k}^\rho )^{}\widehat{D}_{\lambda _k}^\rho =^{}+\frac{1}{\rho ^1}R_\rho ^1.$$
(4.5)
This Laplace type operator $`\mathrm{\Delta }_\rho `$ is an elliptic second order operator on $`\mathrm{\Gamma }(M,𝐒_\rho )`$. If $`M`$ is compact, then $`\mathrm{\Delta }_\rho `$ is non-negative and $`\mathrm{ker}\mathrm{\Delta }_\rho `$ is the intersection of $`\mathrm{ker}\widehat{D}_{\lambda _k}^\rho `$ for $`1kN`$.
## 5 A lower bound of $`\mathrm{\Delta }_\rho `$ on manifolds of positive curvature
We assume that $`M`$ is a compact manifold of positive curvature in this section. In other words, there exists a constant $`r>0`$ such that the curvature $`R_{\mathrm{Ad}}(,)`$ on the tangent bundle $`T(M)`$ satisfies
$$R_{\mathrm{Ad}}(u,v)v,u2ru,uv,v\text{for any }u\text{ and }v\text{ in }T(M)\text{ }.$$
(5.1)
Then the curvature transformation $`R_\rho ^1`$ on $`𝐒_\rho `$ has a lower bound as follows: for $`\varphi `$ in $`𝐒_\rho `$,
$$\begin{array}{cc}\hfill R_\rho ^1\varphi ,\varphi & =\frac{1}{4}\underset{i,j}{}\pi _\rho ([e_i,e_j])R_\rho (e_i,e_j)\varphi ,\varphi \hfill \\ & =\frac{1}{64}\underset{i,j,k,l}{}R_{\mathrm{Ad}}(e_i,e_j)e_k,e_l\pi _\rho ([e_i,e_j])\pi _\rho ([e_k,e_l])\varphi ,\varphi \hfill \\ & \frac{r}{32}\pi _\rho ([e_i,e_j])\pi _\rho ([e_j,e_i])\varphi ,\varphi \hfill \\ & =2rc(\rho )\varphi ,\varphi ,\hfill \end{array}$$
(5.2)
where $`c(\rho )`$ is a negative constant due to the Casimir operator on $`V_\rho `$ in (1.8).
###### Theorem 5.1.
Let $`M`$ be a compact manifold of positive curvature as above and $`\mathrm{\Delta }_\rho `$ be the Laplace type operator on $`\mathrm{\Gamma }(M,𝐒_\rho )`$ given in (4.5). The first eigenvalue $`\lambda ^0`$ of $`\mathrm{\Delta }_\rho `$ has a lower bound,
$$\lambda ^0\frac{m(\lambda _0)m(\lambda _N)}{m(\lambda _0)m(\lambda _N)}(2rc(\rho ))>0.$$
(5.3)
If the equality holds in (5.3), then the eigensections with the first eigenvalue $`\lambda ^0`$ are in $`\mathrm{ker}D_{\lambda _0}^\rho `$ and $`\mathrm{ker}\widehat{D}_{\lambda _k}^\rho `$ for $`1kN1`$.
###### Proof.
We remark that $`m(\lambda _N)>0`$ and, for any $`k`$,
$$\frac{m(\lambda _0)m(\lambda _N)}{m(\lambda _0)m(\lambda _N)}\frac{m(\lambda _0)m(\lambda _k)}{m(\lambda _0)m(\lambda _k)}.$$
We have known the following identities:
$$\mathrm{\Delta }_\rho =(\widehat{D}_{\lambda _k}^\rho )^{}\widehat{D}_{\lambda _k}^\rho ,$$
$$R_\rho ^1+\rho ^1(D_{\lambda _0}^\rho )^{}D_{\lambda _0}^\rho =\underset{1kN}{}\frac{m(\lambda _0)m(\lambda _k)}{m(\lambda _0)m(\lambda _k)}(\widehat{D}_{\lambda _k}^\rho )^{}\widehat{D}_{\lambda _k}^\rho .$$
Then
$$\begin{array}{cc}\hfill (\mathrm{\Delta }_\rho \varphi ,\varphi )& =(\widehat{D}_{\lambda _k}^\rho \varphi ,\widehat{D}_{\lambda _k}^\rho \varphi )\hfill \\ & =\frac{m(\lambda _0)m(\lambda _N)}{m(\lambda _0)m(\lambda _N)}\frac{m(\lambda _0)m(\lambda _N)}{m(\lambda _0)m(\lambda _N)}\widehat{D}_{\lambda _k}^\rho \varphi ^2\hfill \\ & \frac{m(\lambda _0)m(\lambda _N)}{m(\lambda _0)m(\lambda _N)}\frac{m(\lambda _0)m(\lambda _k)}{m(\lambda _0)m(\lambda _k)}\widehat{D}_{\lambda _k}^\rho \varphi ^2\hfill \\ & =\frac{m(\lambda _0)m(\lambda _N)}{m(\lambda _0)m(\lambda _N)}((R_\rho ^1\varphi ,\varphi )+\rho ^1D_{\lambda _0}^\rho \varphi ^2)\hfill \\ & \frac{m(\lambda _0)m(\lambda _N)}{m(\lambda _0)m(\lambda _t)}(R_\rho ^1\varphi ,\varphi )=\frac{m(\lambda _0)m(\lambda _N)}{m(\lambda _0)m(\lambda _N)}_MR_\rho ^1\varphi ,\varphi 1\hfill \\ & \frac{m(\lambda _0)m(\lambda _N)}{m(\lambda _0)m(\lambda _N)}(2rc(\rho ))(\varphi ,\varphi ).\hfill \end{array}$$
This inequality leads us to the proposition. ∎
###### Example 5.1 (Differential form).
We give a lower bound of the first eigenvalue $`\lambda ^0`$ of $`dd^{}+d^{}d`$ on $`\mathrm{\Lambda }^k(M)`$, where $`M`$ is a compact $`n`$-dimensional manifold of positive curvature. From the above theorem, we have $`\lambda ^0k(nk+1)r`$ for $`1k[\frac{n}{2}]`$. If the equality holds in the equation, the eigensections with the eigenvalue $`\lambda ^0`$ is in $`\mathrm{ker}C\mathrm{ker}d`$. This lower bound coincides with the one given in .
## Acknowledgements
The author is partially supported by Waseda University Grant for Special Research Project 2000A-880.
## Appendix A Appendix: Conformal invariance of the higher spin Dirac operator
In , Fegan shows that all the higher spin Dirac operators are conformal invariant first order differential operators. We show the conformal invariance of $`D_\lambda ^\rho `$ explicitly by using the Clifford homomorphisms. Let $`(M,g)`$ be a spin manifold with Riemannian metric $`g`$. We deform the metric $`g`$ conformally to $`g^{}:=e^{2\sigma (x)}g`$, where $`\sigma (x)`$ is a scalar function on $`M`$. We denote the objects associated to $`g^{}`$ by adding a symbol “ ” to them: for example, $`\mathrm{𝐒𝐩𝐢𝐧}^{}(M)`$, $`𝐒_\rho ^{}`$ and so on.
The isomorphism between the orthonormal frame bundles $`\mathrm{𝐒𝐎}(M)`$ and $`\mathrm{𝐒𝐎}^{}(M)`$ is realized by the mapping
$$\mathrm{\Psi }:\mathrm{𝐒𝐎}(M)\{e_i\}_i\{e_i^{}:=e^\sigma e_i\}_i\mathrm{𝐒𝐎}^{}(M).$$
(A.1)
This mapping is lifted to the isomorphism for spin bundles and induces a bundle isometry $`\psi _\rho :𝐒_\rho 𝐒_\rho ^{}`$ for each associated bundle such that $`\psi _\lambda p_\lambda ^\rho (e_i)=e^\sigma p_\lambda ^\rho (e_i)\psi _\rho `$. The covariant derivatives $``$ and $`^{}`$ on $`\mathrm{\Gamma }(M,𝐒_\rho )`$ and $`\mathrm{\Gamma }(M,𝐒_\rho ^{})`$ are related as follows (for the spinor case, see and ):
$$_X^{}=\psi _\rho \{_X\varphi +\frac{1}{4}\pi _\rho ([\mathrm{grad}(\sigma ),X])\}\psi _\rho ^1,$$
(A.2)
where $`X`$ is any vector field on $`M`$. This relation implies a relation for the higher spin Dirac operator.
###### Lemma A.1.
The higher spin Dirac operators $`D_\lambda ^\rho `$ and $`D^{}_\lambda ^\rho `$ associated to metric $`g`$ and $`g^{}`$ respectively are related as follows:
$$D^{}{}_{\lambda }{}^{\rho }=e^\sigma \psi _\lambda \{D_\lambda ^\rho +m(\lambda )p_\lambda ^\rho (\mathrm{grad}(\sigma ))\}\psi _\rho ^1$$
(A.3)
###### Proof.
For $`\varphi `$ in $`\mathrm{\Gamma }(M,𝐒_\rho )`$, we have
$$\begin{array}{cc}\hfill D^{}{}_{\lambda }{}^{\rho }\psi _\rho (\varphi )& =p_\lambda ^\rho (e_i^{})_{e_i^{}}^{}\psi _\rho (\varphi )\hfill \\ & =p_\lambda ^\rho (e_i)_{e_i}^{}\psi (\varphi )\hfill \\ & =p_\lambda ^\rho (e_i)\psi _\rho \{_{e_i}\varphi +\frac{1}{4}\pi _\rho ([\mathrm{grad}(\sigma ),e_i])\varphi \}\hfill \\ & =e^\sigma \psi _\lambda \{p_\lambda ^\rho (e_i)(_{e_i}\varphi +\frac{1}{4}\pi _\rho ([\mathrm{grad}(\sigma ),e_i])\varphi )\}\hfill \\ & =e^\sigma \psi _\lambda \{D_\lambda ^\rho \varphi +m(\lambda )p_\lambda ^\rho (\mathrm{grad}(\sigma ))\varphi \}\text{( from (}\text{2.4}\text{))}.\hfill \end{array}$$
The usual Dirac operator $`D`$ satisfies that $`Df=\mathrm{grad}(f)+fD`$, where $`f`$ is any smooth function on $`M`$. The higher spin Dirac operator has the similar property.
###### Lemma A.2.
For any smooth function $`f`$ on $`M`$, we have
$$D_\lambda ^\rho f=p_\lambda ^\rho (\mathrm{grad}(f))+fD_\lambda ^\rho .$$
(A.4)
The above two lemma gives us the conformal invariance of the operator $`D_\lambda ^\rho `$:
$$D^{}{}_{\lambda }{}^{\rho }=(e^{(m(\lambda )+1)\sigma }\psi _\lambda )D_\lambda ^\rho (e^{m(\lambda )\sigma }\psi _\rho )^1.$$
(A.5)
In particular, we have $`dim\mathrm{ker}D^{}{}_{\lambda }{}^{\rho }=dim\mathrm{ker}D_\lambda ^\rho `$.
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# Freiburg-THEP 00/11 TTP00–15 hep-ph/0007088 Three-loop Three-Linear Vertices and Four-Loop ("mom")̃ 𝛽 functions in massless QCD
## 1 Introduction
Momentum subtraction schemes provide a possibility to define a renormalization description for QCD in a regularization independent way. The concept of momentum subtraction is very old. It played an important role in the discussion of renormalization description dependence of physical quantities long ago . Recently, the momentum subtraction approach has been heavily used to relate lattice results for quark masses and coupling constants to their perturbatively determined $`\overline{\text{ms}}`$ counterparts (see, e.g. ). In it has been argued that the knowledge of three-loop coefficients for the corresponding $`\beta `$-functions is necessary and even the four-loop contributions should be taken in to account. This is because the accessible energy ranges in these calculations are just reaching a level where perturbative QCD calculations start to be valid approximations.
A large subclass (in fact infinitely many) of mom schemes can be defined by subtracting vertices at the asymmetric point where one external momentum vanishes. We will call this point the zero point (ZP). These schemes are referred to as $`\stackrel{~}{\text{mom}}`$ schemes and were introduced and discussed in some detail at two-loop order in . A crucial fact for $`\stackrel{~}{\text{mom}}`$ schemes is that setting one external momentum to zero for all three 3-vertices of massless QCD never will produce infrared divergencies.
In this paper we present the perturbative calculation of the gluon, ghost and quark self-energies and all fundamental<sup>1</sup><sup>1</sup>1That is appearing in the QCD Lagrangian 3-vertices with one vanishing external momentum for massless QCD in a general covariant gauge. The one-loop triple gluon vertex was obtained in in Feynman gauge, the result for general gauge can be found in . At two loops, the triple gluon and ghost gluon vertices at the ZP have been determined in general gauge in . References to earlier relevant publications and results for various momentum configurations can also be found in this work. The quark-gluon vertex can be found to two-loop order in Feynman gauge in .
These three-loop results at hand allow one to relate the coupling constants of any $`\stackrel{~}{\text{mom}}`$-like scheme to the $`\overline{\text{ms}}`$ scheme at three-loop order. Recently, in even the four-loop term of the $`\overline{\text{ms}}`$ $`\beta `$-function has been computed. Using this result, we can also determine the $`\beta `$-functions of any such $`\stackrel{~}{\text{mom}}`$-scheme up to (and including) four loops.
The paper is organized as follows: In Section 2 the definitions of the ghost and quark self-energy and the gluon polarization are introduced and a brief outline of their calculation is given. In Section 3 we introduce the triple gluon, quark gluon and ghost gluon vertices at the ZP and our notation for these along with a description of their calculation. The calculation of the triple gluon vertex is performed directly and additionally in an independent way using the Ward-Slavnov-Taylor (WST) identity relating the triple gluon vertex to the ghost gluon vertex. This is content of Section 4. In Section 5 the $`\overline{\text{ms}}`$ renormalization procedure is described. In Sections 6 and 7 we use the results of the previous Sections to obtain the coupling constants and $`\beta `$-functions for 4 different $`\stackrel{~}{\text{mom}}`$ schemes of particular interest from their $`\overline{\text{ms}}`$ counterparts. Finally, in Section 8 we discuss the numerical importance of our results on the example of a recent lattice computation of the momentum dependence (running) of the three-gluon asymmetrical vertex.
The complete results are given in the appendix and also will be made available as input file for the algebraic programs form and mathematica at:
http://www-ttp.physik.uni-karlsruhe.de/Progdata/
## 2 The Gluon Polarization and the Ghost and Quark Self-Energies
The QCD Lagrangian with $`n_f`$ massless quark flavors in the covariant gauge is:
$``$ $`=`$ $`{\displaystyle \frac{1}{4}}G_{\mu \nu }^aG^{a\mu \nu }+\mathrm{i}{\displaystyle \underset{f=1}{\overset{n_f}{}}}\overline{\psi }_i^f[/D]_{ij}\psi _j^f{\displaystyle \frac{1}{2\xi _L}}(^\mu A_\mu ^a)^2+^\mu \overline{\eta }^a(\eta ^agf^{abc}\eta ^bA_\mu ^c),`$ (1)
$`G_{\mu \nu }^a`$ $`=`$ $`_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^aA_\nu ^b,[D_\mu ]_{ij}=\delta _{ij}_\mu \mathrm{i}gA_\mu ^aT_{ij}^a.`$
The quark fields $`\psi _i^f`$ transform as the fundamental representation and the gluon fields $`A_\mu ^a`$ as the adjoint representation of the gauge group $`\mathrm{SU}(3)`$. $`T_{ij}^a`$ and $`f^{abc}`$ are the generators of the fundamental and adjoint representation of the corresponding Lie algebra. The $`\eta ^a`$ are the ghost fields and $`\xi _L`$ is the gauge parameter ($`\xi _L=0`$ corresponds to Landau gauge).
From this Lagrangian one can derive three types of 2-point functions
$`G_{\mu \nu }^{(2)ab}(q)`$ $`=`$ $`D_{\mu \nu }^{ab}(q)=\mathrm{i}{\displaystyle dx\mathrm{e}^{\mathrm{i}qx}T[A_\mu ^a(x)A_\nu ^b(0)]},`$
$`G^{(2)ab}(q)`$ $`=`$ $`\mathrm{\Delta }^{ab}(q)=\mathrm{i}{\displaystyle dx\mathrm{e}^{\mathrm{i}qx}T[\eta ^a(x)\overline{\eta }^b(0)]},`$ (2)
$`G_{ij}^{(2)}(q)`$ $`=`$ $`S_{ij}(q)=\mathrm{i}{\displaystyle dx\mathrm{e}^{\mathrm{i}qx}T[\psi _i(x)\overline{\psi }_j(0)]},`$
where as usual $`T[AB]`$ is the time ordered product of $`A`$ and $`B`$ and from now on we will skip the flavor index of the quarks. The propagators in eq. (2) can be expressed in terms of the corresponding self-energies in the following way:
$`D_{\mu \nu }^{ab}(q)`$ $`=`$ $`{\displaystyle \frac{\delta ^{ab}}{(q^2)}}\left[(g_{\mu \nu }+{\displaystyle \frac{q_\mu q_\nu }{q^2}}){\displaystyle \frac{1}{1+\mathrm{\Pi }(q^2)}}\xi _L{\displaystyle \frac{q_\mu q_\nu }{q^2}}\right],`$
$`\mathrm{\Delta }^{ab}(q)`$ $`=`$ $`{\displaystyle \frac{\delta ^{ab}}{(q^2)}}{\displaystyle \frac{1}{1+\stackrel{~}{\mathrm{\Pi }}(p^2)}},`$ (3)
$`S_{ij}(q)`$ $`=`$ $`{\displaystyle \frac{\delta ^{ij}}{(q^2)}}{\displaystyle \frac{/q}{1+\mathrm{\Sigma }_V(q^2)}}.`$
The self-energies $`\mathrm{\Pi }(q^2)`$, $`\stackrel{~}{\mathrm{\Pi }}(q^2)`$ and $`\mathrm{\Sigma }_V(q^2)`$ can be calculated by applying the following projections to all one particle irreducible (1PI) diagrams with two external legs of the corresponding type:
$`\mathrm{\Pi }(q^2)`$ $`=`$ $`{\displaystyle \frac{\delta ^{ab}}{N^21}}\left[{\displaystyle \frac{1}{D1}}(g_{\mu \nu }+{\displaystyle \frac{q_\mu q_\nu }{q^2}}){\displaystyle \frac{1}{q^2}}\right]\times \text{},`$
$`\stackrel{~}{\mathrm{\Pi }}(q^2)`$ $`=`$ $`{\displaystyle \frac{\delta ^{ab}}{N^21}}{\displaystyle \frac{1}{q^2}}\times \text{},`$ (4)
$`\mathrm{\Sigma }_V(q^2)`$ $`=`$ $`{\displaystyle \frac{\delta ^{ab}}{N}}\mathrm{Tr}\left[{\displaystyle \frac{/q}{4q^2}}\times \text{}\right].`$
These projectors take into account that we are using dimensional regularization with $`D=42ϵ`$ and are valid for a general $`\mathrm{SU}(N)`$ gauge group. The trace has to be taken over Dirac matrices. The resulting scalar integrals are of the massless propagator type and can be evaluated with standard methods. For details about the technical setup used to perform these calculations we refer to the appendix.
## 3 The Triple Gluon, Ghost and Quark Vertex Functions
In massless QCD there are the following 3-point functions:
$`G_{\mu \nu \rho }^{(3)abc}(p,q)`$ $`=`$ $`\mathrm{i}^2{\displaystyle dxdy\mathrm{e}^{\mathrm{i}(px+qy)}T[A_\mu ^a(x)A_\nu ^b(y)A_\rho ^c(0)]},`$
$`G_\mu ^{(3)abc}(p,q)`$ $`=`$ $`\mathrm{i}^2{\displaystyle dxdy\mathrm{e}^{\mathrm{i}(px+qy)}T[\eta ^b(x)\overline{\eta }^c(y)A_\mu ^a(0)]},`$ (5)
$`G_{\mu ij}^{(3)a}(p,q)`$ $`=`$ $`\mathrm{i}^2{\displaystyle dxdy\mathrm{e}^{\mathrm{i}(px+qy)}T[\psi _i(x)\overline{\psi }_j(y)A_\mu ^a(0)]}.`$
These are related to the vertex functions by
$`G_{\mu \nu \rho }^{(3)abc}(p,q)`$ $`=`$ $`D_{\mu \mu ^{}}^{ad}(p)D_{\nu \nu ^{}}^{be}(q)D_{\rho \rho ^{}}^{cf}(pq)\mathrm{\Gamma }_{\mu ^{}\nu ^{}\rho ^{}}^{def}(p,q,pq),`$
$`G_\mu ^{(3)abc}(p,q)`$ $`=`$ $`\mathrm{\Delta }^{ad}(p)\stackrel{~}{\mathrm{\Gamma }}_\mu ^{}^{def}(p,q;pq)\mathrm{\Delta }^{eb}(q)D_{\mu ^{}\mu }^{cf}(q+p),`$ (6)
$`G_{\mu ij}^{(3)a}(p,q)`$ $`=`$ $`S_{ii^{}}(p)\mathrm{\Lambda }_{\mu ^{}i^{}j^{}}^d(p,q;qp)S_{j^{}j}(q)D_{\mu ^{}\mu }^{ad}(p+q).`$
Setting one external momentum to zero leaves only one momentum which can be used to construct tensor decompositions of the vertex functions. The color and Lorentz structure of all vertices at the ZP and our conventions are discussed in the next subsections.
### 3.1 The Triple Gluon Vertex
For symmetry reasons setting any of the three external gluon momenta of the triple gluon vertex to zero will lead to the same scalar functions. Up to two loops, the triple gluon vertex is known to be proportional to the totally antisymmetric color structure functions $`f^{abc}`$. A color structure proportional to the totally symmetric $`d^{abc}`$ will not have a divergent part. In this work we only consider the $`f^{abc}`$ part.
The triple gluon vertex needs to be totally symmetric under exchange of any two of the (bosonic) gluons and we are interested in the part proportional to $`f^{abc}`$. So the Lorentz structure for this part in the case of one vanishing external momentum is limited to the following three tensor structures, which are all antisymmetric under exchange of the two gluons with non vanishing momentum:
$`\mathrm{\Gamma }_{\mu \nu \rho }^{abc}(q,q,0)`$ $`=`$ $`\mathrm{i}gf^{abc}(\text{}(2g_{\mu \nu }q_\rho g_{\mu \rho }q_\nu g_{\rho \nu }q_\mu )T_1(q^2)`$ (7)
$`(g_{\mu \nu }{\displaystyle \frac{q_\mu q_\nu }{q^2}})q_\rho T_2(q^2)+q_\mu q_\nu q_\rho T_3(q^2)).`$
Due to the WST identity for the triple gluon vertex(for details see Section 4) the third function needs to vanish and finding $`T_3(q^2)=0`$ is a check for the correctness of the result. The functions $`T_i(q^2)`$ can directly be calculated by applying the following projectors to the 1PI diagrams with 3 external gluon legs of which one carries zero momentum:
$`P_{1\mu \nu \rho }^{abc}(q)`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}f^{abc}}{N(N^21)}}\left[{\displaystyle \frac{1}{D1}}\left({\displaystyle \frac{q_\mu q_\nu q_\rho }{q^4}}{\displaystyle \frac{g_{\nu \rho }q_\mu }{2q^2}}{\displaystyle \frac{g_{\mu \rho }q_\nu }{2q^2}}\right)\right],`$
$`P_{2\mu \nu \rho }^{abc}(q)`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}f^{abc}}{N(N^21)}}\left[{\displaystyle \frac{1}{D1}}\left({\displaystyle \frac{3q_\mu q_\nu q_\rho }{q^4}}g_{\mu \nu }{\displaystyle \frac{q_\rho }{q^2}}g_{\mu \rho }{\displaystyle \frac{q_\nu }{q^2}}g_{\nu \rho }{\displaystyle \frac{q_\mu }{q^2}}\right)\right],`$
$`P_{3\mu \nu \rho }^{abc}(q)`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}f^{abc}}{N(N^21)}}\left[{\displaystyle \frac{q_\mu q_\nu q_\rho }{q^6}}\right],`$ (8)
$$T_i(q^2)=P_{i\mu \nu \rho }^{abc}(q)\times \text{}.$$
(9)
The explicit calculation shows that indeed $`T_3`$ vanishes.
### 3.2 The Ghost Gluon Vertex
The ghost gluon vertex is at tree level proportional to the momentum of the outgoing ghost. Since this vertex is the only interaction of the ghost field, it is clear that the vertex is also proportional to this momentum at any order in perturbation theory. This leaves two possibilities of one zero external momentum and due to the simple Lorentz structure also only two scalar functions to be determined (again we only consider only the $`f^{abc}`$ color structure):
$`\stackrel{~}{\mathrm{\Gamma }}_\mu ^{abc}(q,0;q)`$ $`=`$ $`\mathrm{i}gf^{abc}q_\mu \stackrel{~}{\mathrm{\Gamma }}_\mathrm{h}(q^2),`$
$`\stackrel{~}{\mathrm{\Gamma }}_\mu ^{abc}(q,q;0)`$ $`=`$ $`\mathrm{i}gf^{abc}q_\mu \stackrel{~}{\mathrm{\Gamma }}_\mathrm{g}(q^2).`$ (10)
The subscripts g and h stand for the external line that carries zero momentum: g for the gluon and h for the ghost. Again the $`\stackrel{~}{\mathrm{\Gamma }}_i(q^2)`$ can be directly computed from the 1PI diagrams with two external ghosts and one external gluon:
$`\stackrel{~}{\mathrm{\Gamma }}_\mathrm{h}(q^2)`$ $`=`$ $`+{\displaystyle \frac{\mathrm{i}f^{abc}}{N(N^21)}}{\displaystyle \frac{q_\mu }{q^2}}\times \text{},`$ (11)
$`\stackrel{~}{\mathrm{\Gamma }}_\mathrm{g}(q^2)`$ $`=`$ $`+{\displaystyle \frac{\mathrm{i}f^{abc}}{N(N^21)}}{\displaystyle \frac{q_\mu }{q^2}}\times \text{}.`$ (12)
### 3.3 The Quark Gluon Vertex
Finally the quark gluon vertex is proportional to the color structure $`T_{ij}^a`$ and for one vanishing external momentum there are two different possibilities (here setting either of the quark momenta to zero gives the same scalar functions up to three loops). A useful tensor decomposition is:
$`\mathrm{\Lambda }_{\mu ij}^a(q,0;q)`$ $`=`$ $`gT_{ij}^a\left[\gamma _\mu \mathrm{\Lambda }_\mathrm{q}(q^2)+\gamma _\nu \left(g_{\mu \nu }{\displaystyle \frac{q_\mu q_\nu }{q^2}}\right)\mathrm{\Lambda }_\mathrm{q}^T(q^2)\right],`$
$`\mathrm{\Lambda }_{\mu ij}^a(q,q;0)`$ $`=`$ $`gT_{ij}^a\left[\gamma _\mu \mathrm{\Lambda }_\mathrm{g}(q^2)+\gamma _\nu \left(g_{\mu \nu }{\displaystyle \frac{q_\mu q_\nu }{q^2}}\right)\mathrm{\Lambda }_\mathrm{g}^T(q^2)\right],`$ (13)
where the additional subscript q corresponds to the case of a vanishing external quark momentum and $`T`$ marks the transversal part. The $`\mathrm{\Lambda }_i`$ can directly be calculated in the following way:
$`\mathrm{\Lambda }_\mathrm{q}(q^2)`$ $`=`$ $`+{\displaystyle \frac{T_{ij}^a}{NC_F}}Tr\left[{\displaystyle \frac{/qq_\mu }{4q^2}}\times \text{}\right],`$
$`\mathrm{\Lambda }_\mathrm{q}^T(q^2)`$ $`=`$ $`{\displaystyle \frac{T_{ij}^a}{NC_F}}Tr\left[{\displaystyle \frac{1}{4(D1)}}\left(\gamma _\mu D{\displaystyle \frac{/qq_\mu }{q^2}}\right)\times \text{}\right],`$ (14)
$`\mathrm{\Lambda }_\mathrm{g}(q^2)`$ $`=`$ $`+{\displaystyle \frac{T_{ij}^a}{NC_F}}Tr\left[{\displaystyle \frac{/qq_\mu }{4q^2}}\times \text{}\right],`$
$`\mathrm{\Lambda }_\mathrm{g}^T(q^2)`$ $`=`$ $`{\displaystyle \frac{T_{ij}^a}{NC_F}}Tr\left[{\displaystyle \frac{1}{4(D1)}}\left(\gamma _\mu D{\displaystyle \frac{/qq_\mu }{q^2}}\right)\times \text{}\right].`$ (15)
## 4 The Ward-Slavnov-Taylor Identity
By using the WST identity for the triple gluon vertex, one can determine the scalar functions $`T_1`$ and $`T_2`$ from the gluon and ghost self-energies and functions related to the ghost gluon vertex. This check has also been performed in other determinations of the gluon vertex . The general WST identity was determined in and can be found in the following form for the triple gluon vertex in :
$`k^\rho \mathrm{\Gamma }_{\mu \nu \rho }(p,q,k)`$ $`=`$ $`J(p^2)G(k^2)(g_\mu ^\rho p^2p^\rho p_\mu )\stackrel{~}{\mathrm{\Gamma }}_{\rho \nu }(p,k;q)`$ (16)
$`+J(q^2)G(k^2)(g_\nu ^\rho q^2q^\rho q_\nu )\stackrel{~}{\mathrm{\Gamma }}_{\rho \mu }(q,k;p).`$
Here $`\stackrel{~}{\mathrm{\Gamma }}_{\nu \mu }(p,q;k)`$ is related to the proper ghost gluon vertex $`\stackrel{~}{\mathrm{\Gamma }}_\mu ^{abc}(p,q;k)`$ by:
$$\stackrel{~}{\mathrm{\Gamma }}_\mu ^{abc}(p,q;k)=\mathrm{i}gf^{abc}p^\nu \stackrel{~}{\mathrm{\Gamma }}_{\nu \mu }(p,q;k).$$
(17)
and we have introduced the functions
$$J(p^2)=1+\mathrm{\Pi }(p^2)\mathrm{and}G(p^2)=\frac{1}{1+\stackrel{~}{\mathrm{\Pi }}(p^2)}.$$
In ref. the following tensor decomposition for $`\stackrel{~}{\mathrm{\Gamma }}_{\mu \nu }(p,q;k)`$ is given:
$`\stackrel{~}{\mathrm{\Gamma }}_{\mu \nu }(p,k;q)`$ $`=`$ $`+g_{\mu \nu }a(q,k,p)q_\mu k_\nu b(q,k,p)+p_\mu q_\nu c(q,k,p)`$ (18)
$`+q_\mu p_\nu d(q,k,p)+p_\mu p_\nu e(q,k,p).`$
The analytic determination of the full momentum dependence of the generalized ghost gluon vertex $`\stackrel{~}{\mathrm{\Gamma }}_{\nu \mu }`$ at three loops is impossible with the current calculation techniques. But of course we do not need to know the full momentum dependence for checking the WST identity at the ZP. Still, even the determination of all possible tensor structures for the case of one vanishing momentum and the corresponding expansions to first order in the vanishing momentum would be a very demanding project. To avoid the calculation of unnecessary parts of the ghost gluon vertex we follow closely the approach of .
For the vertex function there are two independent momenta. They can be chosen in such a way that at the ZP one of them vanishes. Contracting the vertex in the WST identity (16) with the momentum that does not vanish at the ZP is relatively simple. In this case one can safely set the other momentum to zero on the right hand side and is left with functions of just one momentum. Inserting the tensor decompositions (7) and (18) gives:
$$T_1(p^2)=a_3(p^2)G(p^2)J(p^2),a_3(p^2)=a(0,p,p).$$
(19)
There is also the possibility to contract the vertex with the momentum that vanishes in the limit of the ZP. This case gives a differential WST identity at the ZP. If the right hand side is expanded to first order in $`k_\rho `$ the constant lowest order term cancels. Dividing by $`k_\rho `$, setting $`k`$ to $`0`$ on both sides and taking into account that for massless quarks $`G(0)=1`$<sup>2</sup><sup>2</sup>2$`G(0)=1`$ since all contributing diagrams are massless tadpoles which vanish in dimensional regularization. results in the following representation of this differential identity:
$`\mathrm{\Gamma }_{\mu \nu \rho }(p,p,0)`$ $`=`$ $`+a_2(p^2)J(p^2)\left(2g_{\mu \nu }p_\rho g_{\mu \rho }p_\nu g_{\nu \rho }p_\mu \right)`$ (20)
$`+a_2(p^2)p^2{\displaystyle \frac{\mathrm{d}J(p^2)}{\mathrm{d}p^2}}\left(2{\displaystyle \frac{p_\mu p_\nu p_\rho }{p^2}}+2g_{\mu \nu }p_\rho \right)`$
$`+d_2(p^2)p^2J(p^2)\left(2{\displaystyle \frac{p_\mu p_\nu p_\rho }{p^2}}+g_{\mu \rho }p_\nu +g_{\nu \rho }p_\mu \right)`$
$`+J(p^2)\left(a_{23}^{}(p^2)a_{21}^{}(p^2)\right)\left({\displaystyle \frac{p_\mu p_\nu p_\rho }{p^2}}g_{\mu \nu }p_\rho \right).`$
where we have expanded $`a(p,k,q)`$ to first order in small $`k`$ and use the following shortcuts:
$`d(p,0,p)=d_2(p^2)`$ , $`a(p,0,p)=a(p,0,p)=a_2(p^2),`$
$`a(pk,k,p)`$ $`=`$ $`a_2(p^2)+{\displaystyle \frac{kp}{p^2}}a_{23}^{}(p^2)+𝒪(k^2),`$ (21)
$`a(p,k,pk)`$ $`=`$ $`a_2(p^2)+{\displaystyle \frac{kp}{p^2}}a_{21}^{}(p^2)+𝒪(k^2).`$
Contracting eq. (20) with $`p^\mu `$ once more and projecting out the structure $`T_1(p^2)`$ gives another possibility to reproduce $`T_1(p^2)`$. It can be used to produce the simple relation
$$a_2(p^2)p^2d_2(p^2)=G(p^2)a_3(p^2).$$
(22)
Using this relation and projecting out the structure proportional to $`T_2(p^2)`$ we can also obtain $`T_2(p^2)`$ in a completely independent way:
$`T_2(p^2)`$ $`=`$ $`2T_1(p^2)2a_(p^2){\displaystyle \frac{\mathrm{d}}{\mathrm{d}p^2}}(p^2J(p^2))`$ (23)
$`+J(p^2)\left(a_{23}^{}(p^2)a_{21}^{}(p^2)\right).`$
Finally, applying the projector $`P_3`$ of eq. (3.1) to the right hand side of eq. (20) should give an alternative representation of $`T_3(p^2)`$ and it is indeed found to be zero from this identity, as has been mentioned in the last Section.
The scalar functions $`a_2`$, $`a_3`$, $`a_{23}^{}`$ and $`a_{21}^{}`$ can be calculated from 1PI diagrams with the following operations:
$$P_{\mu \nu }^{abc}(q)=\frac{f^{abc}}{N(N^21)}\frac{1}{D1}\left(q^2g_{\mu \nu }q_\mu q_\nu \right),$$
(24)
$$a_2(q^2)=1+P_{\mu \nu }^{abc}(q)\times \left[\text{}\right],$$
(25)
$$a_3(q^2)=1+P_{\mu \nu }^{abc}(q)\times \left[\text{}\right],$$
(26)
$$a_{23}^{}(q^2)=\mathrm{}_k\left(qkP_{\mu \nu }^{abc}(q)𝒯_k^{(1)}\times \left[\text{}\right]\right),$$
(27)
$$a_{21}^{}(q^2)=\mathrm{}_k\left(qkP_{\mu \nu }^{abc}(q)𝒯_k^{(1)}\times \left[\text{}\right]\right).$$
(28)
The additional gluon line is to visualize that at the outgoing ghost vertex we do not contract with the external momentum, which leaves an additional open index $`\nu `$ as shown in eq. (17). The shortcut $`𝒯_k^{(1)}`$ stands for a first order Taylor expansion in $`k`$ and $`\mathrm{}_k=g_{\alpha \beta }\frac{}{k_\alpha }\frac{}{k_\beta }`$. Both are to be applied to the integrand before any of the integrations are performed. These manipulations are necessary because the mincer package can only handle scalar products of internal and one external momentum for more than one loop. Since there are only the two independent external momenta $`p`$ and $`k`$, the first order terms in $`k`$ (of a scalar function) have to be proportional to $`kp`$. This guarantees that the above manipulations will extract the scalars $`a_{2i}^{}`$.
Using this strategy we found complete agreement for the bare expression with the ones from the direct computation of the triple gluon vertex.
## 5 Renormalization
For a generic renormalization scheme<sup>3</sup><sup>3</sup>3In the following all quantities without explicit superscript refer to the $`\overline{\text{ms}}`$ scheme, the superscript $`\mathrm{B}`$ marks the bare quantities $`\mathrm{R}`$, the following relations between bare and renormalized quantities hold in dimensional regularization:
$`(A^\mathrm{B})_\nu ^a=\sqrt{Z_3^\mathrm{R}(h^\mathrm{R},\mu ,ϵ)}(A^\mathrm{R})_\nu ^a(\mu ),`$ $`\eta ^\mathrm{B}=\sqrt{\stackrel{~}{Z}_3^\mathrm{R}(h^\mathrm{R},\mu ,ϵ)}\eta ^\mathrm{R}(\mu ),`$
$`\psi _i^\mathrm{B}=\sqrt{Z_2^\mathrm{R}(h^\mathrm{R},\mu ,ϵ)}\psi _i^\mathrm{R}(\mu ),`$ $`h^\mathrm{B}=\mu ^{2ϵ}Z_h^\mathrm{R}(h^\mathrm{R},\mu ,ϵ)h^\mathrm{R}(\mu ).`$ (29)
Here $`h=\alpha _s/(4\pi )=g_s^2/(16\pi ^2)`$, $`\mu `$ is the ’t Hooft unit of mass which is the renormalization point in the ms scheme.
The self-energies and vertex functions are renormalized as follows
$`1+\mathrm{\Pi }^\mathrm{R}(q^2,h^R,\mu ,ϵ)`$ $`=`$ $`Z_3^\mathrm{R}(h^\mathrm{R},\mu ,ϵ)\left[1+\mathrm{\Pi }^\mathrm{B}(q^2,h^\mathrm{B},ϵ)\right]_{h^\mathrm{B}=\mu ^{2ϵ}Z_hh,\xi _L^\mathrm{B}=Z_3\xi _L},`$
$`1+\stackrel{~}{\mathrm{\Pi }}^\mathrm{R}(q^2,h^R,\mu ,ϵ)`$ $`=`$ $`\stackrel{~}{Z}_3^\mathrm{R}(h^\mathrm{R},\mu ,ϵ)\left[1+\stackrel{~}{\mathrm{\Pi }}^\mathrm{B}(q^2,h^\mathrm{B},ϵ)\right]_{h^\mathrm{B}=\mu ^{2ϵ}Z_hh,\xi _L^\mathrm{B}=Z_3\xi _L},`$
$`1+\mathrm{\Sigma }_V^\mathrm{R}(q^2,h^R,\mu ,ϵ)`$ $`=`$ $`Z_2^\mathrm{R}(h^\mathrm{R},\mu ,ϵ)\left[1+\mathrm{\Sigma }_V^\mathrm{B}(q^2,h^\mathrm{B},ϵ)\right]_{h^\mathrm{B}=\mu ^{2ϵ}Z_hh,\xi _L^\mathrm{B}=Z_3\xi _L},`$
$`T_{1,2}^\mathrm{R}(q^2,h^R,\mu ,ϵ)`$ $`=`$ $`Z_1^\mathrm{R}(h^\mathrm{R},\mu ,ϵ)T_{1,2}^\mathrm{B}(q^2,h^\mathrm{B},ϵ)|_{h^\mathrm{B}=\mu ^{2ϵ}Z_hh,\xi _L^\mathrm{B}=Z_3\xi _L},`$
$`\stackrel{~}{\mathrm{\Gamma }}_{h,g}^\mathrm{R}(q^2,h^R,\mu ,ϵ)`$ $`=`$ $`\stackrel{~}{Z}_1^\mathrm{R}(h^\mathrm{R},\mu ,ϵ)\stackrel{~}{\mathrm{\Gamma }}_{h,g}^\mathrm{B}(q^2,h^\mathrm{B},ϵ)|_{h^\mathrm{B}=\mu ^{2ϵ}Z_hh,\xi _L^\mathrm{B}=Z_3\xi _L},`$
$`\mathrm{\Lambda }_{g,q}^\mathrm{R}(q^2,h^R,\mu ,ϵ)`$ $`=`$ $`\overline{Z}_1^\mathrm{R}(h^\mathrm{R},\mu ,ϵ)\mathrm{\Lambda }_{g,q}^\mathrm{B}(q^2,h^\mathrm{B},ϵ)|_{h^\mathrm{B}=\mu ^{2ϵ}Z_hh,\xi _L^\mathrm{B}=Z_3\xi _L}.`$ (30)
The renormalization constants for the self-energies can be read of eqs.(2,2). Likewise, from eqs. (3,3) and the QCD Lagrangian the renormalization factors for the trilinear operators can be related to the renormalization of the coupling constant and the field renormalizations
$$\sqrt{Z_a^RZ_3^R}=\frac{Z_1^R}{Z_3^R}=\frac{\stackrel{~}{Z}_1^R}{\stackrel{~}{Z}_3^R}=\frac{\overline{Z}_1^R}{Z_2^R}.$$
(31)
QCD is known to be a renormalizable theory. This means that $`Z_x^R`$ can be found which make the renormalized Greens functions eqs. (2,3) finite when taking the limit $`ϵ0`$ and at the same time fulfill the WST identities, respectively eq. (31), to any order in perturbation theory.
Using all these formulas, renormalization in the $`\overline{\text{ms}}`$-scheme is straightforward. It requires the $`Z_x=Z_x^{\overline{\text{ms}}}(h,ϵ)`$ to only contain poles in $`ϵ`$ and thus to be of the following form:
$$Z_x(h,ϵ)=1+\underset{i>0}{}\frac{1}{ϵ^i}Z_x^{(i)}(h),Z_x^{(i)}(h)=1+\underset{ji}{}h^jZ_x^{(i,j)}$$
(32)
The $`\mu `$-dependence of the fields and the coupling $`h`$ is then described in the usual way by the renormalization group equations:
$`\mu ^2{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\mu ^2}}h(\mu )=ϵh+\beta (h)`$ , $`\beta (h)={\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}h^{(i+2)}\beta _i=h^2{\displaystyle \frac{}{h}}Z_h^{(1)},`$
$`2\mu ^2{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\mu ^2}}A_\nu ^a(\mu )=\gamma _3(h)A_\nu ^a(\mu )`$ , $`\gamma _3(h)={\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}h^{(i+1)}\gamma _{3i}=h{\displaystyle \frac{}{h}}Z_3^{(1)},`$
$`2\mu ^2{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\mu ^2}}\eta (\mu )=\stackrel{~}{\gamma }_3(h)\eta (\mu )`$ , $`\stackrel{~}{\gamma }_3(h)={\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}h^{(i+1)}\stackrel{~}{\gamma }_{3i}=h{\displaystyle \frac{}{h}}\stackrel{~}{Z}_3^{(1)},`$
$`2\mu ^2{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\mu ^2}}\psi _i(\mu )=\gamma _2(h)\psi _i(\mu )`$ , $`\gamma _2(h)={\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}h^{(i+1)}\gamma _{2i}=h{\displaystyle \frac{}{h}}Z_2^{(1)},`$ (33)
where the last equalities can be derived from eqs. (5), the $`D`$ dimensional $`\beta `$-function $`[ϵh+\beta (h)]`$ and the fact that the bare quantities must be independent of $`\mu `$. As the functions in eq. (5) $`\beta `$ and the anomalous dimensions are here defined in $`D`$ dimensions but are finite for $`ϵ0`$.
The $`\overline{\text{ms}}`$ renormalization constants and anomalous dimensions are well known up to three loops and for the $`\beta `$-function and the quark field anomalous dimension $`\gamma _2`$ even the four-loop terms are known. Of course up to three loops the $`Z_i`$ can also be determined independently from the bare results for the self-energies and vertex functions. The corresponding anomalous dimensions are given in the appendix.
Performing this standard renormalization procedure one arrives at the $`\overline{\text{ms}}`$ renormalized expressions for all self-energies and vertex functions. The expressions for these in the limit $`ϵ0`$ at the point $`p^2=\mu ^2`$ are given in the appendix for generic color factors and a generic gauge parameter. The momentum dependence can be obtained from the anomalous dimensions. The fact that indeed all poles in $`ϵ`$ cancel and eq. (31) holds is another check for the two-loop finite part and at least for the pole terms of any of the three-loop results.
## 6 Four Particular $`\stackrel{~}{\text{mom}}`$ like Scheme Definitions
Momentum subtraction schemes are defined by setting some of the 2 and 3-point functions to their tree values for a fixed configuration of the states of the external particles (that is momenta and polarization state) in a certain gauge. This fixes the field renormalization constants at the renormalization point $`\mu `$:
$`1+\mathrm{\Pi }^{\stackrel{~}{\text{mom}}}(\mu ^2)=1`$ $`=`$ $`Z_3^{\stackrel{~}{\text{mom}}}(\mu ^2,ϵ)\left[1+\mathrm{\Pi }^\mathrm{B}(\mu ^2,ϵ)\right],`$
$`\left[1+\stackrel{~}{\mathrm{\Pi }}^{\stackrel{~}{\text{mom}}}(\mu ^2)\right]=1`$ $`=`$ $`\stackrel{~}{Z}_3^{\stackrel{~}{\text{mom}}}(\mu ^2,ϵ)\left[1+\stackrel{~}{\mathrm{\Pi }}^\mathrm{B}(\mu ^2,ϵ)\right],`$
$`\left[1+\mathrm{\Sigma }_V^{\stackrel{~}{\text{mom}}}(\mu ^2,ϵ)\right]=1`$ $`=`$ $`Z_2^{\stackrel{~}{\text{mom}}}(\mu ^2,ϵ)\left[1+\mathrm{\Sigma }_V^\mathrm{B}(\mu ^2,ϵ)\right].`$ (34)
For the renormalization of the coupling constant there are infinitely many possibilities to define a momentum subtraction renormalization scheme, even when considering only the ZP. Not only is there an ambiguity in which vertex to subtract, but also there is the freedom to use a certain linear combination of the scalar functions appearing in the gluon and quark vertices, which can be related to fixed polarization states of the external particles . If one considers the generalized ghost gluon vertex as defined in eq. (17) there are even more possibilities.
Each of these choices defines a different renormalization scheme just as a different choice of the arbitrary renormalization scale. In general it is not possible to fix more than one vertex to it’s tree value at the ZP with $`p^2=\mu ^2`$, since they all are related by Ward identities.
“The” $`\stackrel{~}{\text{mom}}`$ scheme originally was defined in as a scheme in which all three triple-vertices are subtracted at the ZP without violating the Ward identities. This can be achieved by choosing a special set of the appearing scalar functions at one-loop order. Of course this universality does not hold anymore when considering the two- and three-loop order. This is why even in the original publication a generalization of this scheme for two and more loops is given. It subtracts the ghost vertex and corresponds to the original definition at one-loop order. Here we also consider the other two schemes that coincide with the $`\stackrel{~}{\text{mom}}`$ scheme at one-loop order. They are defined by subtracting just the scalar functions that appear at tree level of the triple gluon (we will refer to it as the $`\stackrel{~}{\text{mom}}\mathrm{g}`$ scheme) and quark gluon vertex (quark momentum set to zero, $`\stackrel{~}{\text{mom}}\mathrm{q}`$), respectively.
Another $`\stackrel{~}{\text{mom}}`$ like renormalization scheme definition ($`\stackrel{~}{\text{mom}}\mathrm{gg}`$) that uses the gluon vertex lately has been used to relate lattice results of the triple gluon vertex and the gluon propagator to perturbative calculations. The $`\stackrel{~}{\text{mom}}\mathrm{q}`$ scheme that subtracts the quark gluon vertex has also been used in lattice calculations , but is not used beyond one-loop order in this work.
### 6.1 Subtracting the Ghost Gluon Vertex
A generalization of the original $`\stackrel{~}{\text{mom}}`$ scheme to more than one loop is given by subtracting the ghost gluon vertex with the moment of the incoming ghost set to zero and a longitudinal polarization of the external gluon (the case of a transversally polarized gluon will vanish at the ZP). The renormalization condition is:
$$\stackrel{~}{\mathrm{\Gamma }}_h^{\stackrel{~}{\text{mom}}}(\mu ^2)=1=\stackrel{~}{Z}_1^{\stackrel{~}{\text{mom}}}(\mu ^2)\stackrel{~}{\mathrm{\Gamma }}^\mathrm{B}(\mu ^2).$$
(35)
Using eqs. (5-31,6), we can connect the coupling constant in this scheme to quantities calculable in the $`\overline{\text{ms}}`$ scheme through the following chain of equations:
$`g_s^{\stackrel{~}{\text{mom}}}(\mu )`$ $`=`$ $`\mu ^ϵg_s^b{\displaystyle \frac{1}{\stackrel{~}{Z}_g^{\stackrel{~}{\text{mom}}}(\mu )}}=\mu ^ϵg_s^b{\displaystyle \frac{\stackrel{~}{Z}_3^{\stackrel{~}{\text{mom}}}(\mu )\sqrt{Z_3^{\stackrel{~}{\text{mom}}}(\mu )}}{\stackrel{~}{Z}_1^{\stackrel{~}{\text{mom}}}(\mu )}}`$ (36)
$`=`$ $`g_s^{\overline{\text{ms}}}(\mu )Z_g^{\overline{\text{ms}}}(\mu )\left[{\displaystyle \frac{\stackrel{~}{\mathrm{\Gamma }}_h^\mathrm{B}(p^2)}{(1+\stackrel{~}{\mathrm{\Pi }}^\mathrm{B}(p^2))\sqrt{1+\mathrm{\Pi }^\mathrm{B}(p^2)}}}\right]_{p^2=\mu ^2}`$
$`=`$ $`g_s^{\overline{\text{ms}}}(\mu )\left[{\displaystyle \frac{\stackrel{~}{\mathrm{\Gamma }}_h^{\overline{\text{ms}}}(p^2)}{(1+\stackrel{~}{\mathrm{\Pi }}^{\overline{\text{ms}}}(p^2))\sqrt{1+\mathrm{\Pi }^{\overline{\text{ms}}}(p^2)}}}\right]_{p^2=\mu ^2}.`$
Squaring this and inserting the $`\overline{\text{ms}}`$ expressions for $`\mathrm{\Pi }`$, $`\stackrel{~}{\mathrm{\Pi }}`$ and $`\stackrel{~}{\mathrm{\Gamma }}_h`$ we arrive at the following relation between the coupling in this scheme and the $`\overline{\text{ms}}`$ coupling, expressed as an expansion in $`h=h^{\overline{\text{ms}}}`$ and<sup>4</sup><sup>4</sup>4see Section 7 for details $`\xi _L=\xi _L^{\overline{\text{ms}}}`$:
$`h^{\stackrel{~}{\text{mom}}}`$ $`=`$ $`h+h^2[+{\displaystyle \frac{169}{12}}{\displaystyle \frac{10}{9}}n_f+{\displaystyle \frac{9}{2}}\xi _L+{\displaystyle \frac{3}{4}}\xi _L^2]+h^3[+{\displaystyle \frac{76063}{144}}{\displaystyle \frac{4}{3}}n_f\zeta _35n_f\xi _L`$ (37)
$`{\displaystyle \frac{1913}{27}}n_f+{\displaystyle \frac{100}{81}}n_f^2+{\displaystyle \frac{117}{8}}\zeta _3\xi _L{\displaystyle \frac{351}{8}}\zeta _3+{\displaystyle \frac{1719}{16}}\xi _L+{\displaystyle \frac{549}{16}}\xi _L^2+{\displaystyle \frac{81}{16}}\xi _L^3]`$
$`+h^4[+{\displaystyle \frac{42074947}{1728}}159n_f\zeta _3\xi _L+{\displaystyle \frac{3}{4}}n_f\zeta _3\xi _L^2+{\displaystyle \frac{8362}{27}}n_f\zeta _3+{\displaystyle \frac{2320}{9}}n_f\zeta _5{\displaystyle \frac{6931}{16}}n_f\xi _L`$
$`{\displaystyle \frac{757}{16}}n_f\xi _L^2{\displaystyle \frac{769387}{162}}n_f+{\displaystyle \frac{16}{3}}n_f^2\zeta _3\xi _L+{\displaystyle \frac{28}{9}}n_f^2\zeta _3+{\displaystyle \frac{38}{9}}n_f^2\xi _L+{\displaystyle \frac{199903}{972}}n_f^2`$
$`{\displaystyle \frac{1000}{729}}n_f^3+{\displaystyle \frac{4893}{4}}\zeta _3\xi _L+{\displaystyle \frac{1485}{16}}\zeta _3\xi _L^2{\displaystyle \frac{1341}{32}}\zeta _3\xi _L^3{\displaystyle \frac{117}{32}}\zeta _3\xi _L^4{\displaystyle \frac{60675}{16}}\zeta _3`$
$`{\displaystyle \frac{8505}{16}}\zeta _5\xi _L{\displaystyle \frac{4635}{32}}\zeta _5\xi _L^2+{\displaystyle \frac{405}{16}}\zeta _5\xi _L^3+{\displaystyle \frac{315}{64}}\zeta _5\xi _L^4{\displaystyle \frac{70245}{64}}\zeta _5+{\displaystyle \frac{290371}{64}}\xi _L`$
$`+{\displaystyle \frac{22287}{16}}\xi _L^2+{\displaystyle \frac{21141}{64}}\xi _L^3+{\displaystyle \frac{2547}{64}}\xi _L^4].`$
### 6.2 Subtracting the Quark Gluon Vertex
Using again the renormalization condition that at one-loop level is identical to the original $`\stackrel{~}{\text{mom}}`$ scheme we have:
$$\mathrm{\Lambda }_q^{\stackrel{~}{\text{mom}}\mathrm{q}}(\mu ^2)=1.$$
(38)
This subtracts just the structure proportional to $`\gamma _\mu `$ which is present at the tree level and corresponds to a longitudinally polarized gluon. Similar operations as above then lead to
$$h^{\stackrel{~}{\text{mom}}\mathrm{q}}(\mu )=h^{\overline{\text{ms}}}(\mu )\left[\frac{\mathrm{\Lambda }_q^{\overline{\text{ms}}}(\mu ^2)}{(1+\mathrm{\Sigma }_V^{\overline{\text{ms}}}(\mu ^2))\sqrt{1+\mathrm{\Pi }^{\overline{\text{ms}}}(\mu ^2)}}\right]^2$$
(39)
and inserting the $`\overline{\text{ms}}`$ results gives:
$`h^{\stackrel{~}{\text{mom}}\mathrm{q}}`$ $`=`$ $`h+h^2[+{\displaystyle \frac{169}{12}}{\displaystyle \frac{10}{9}}n_f+{\displaystyle \frac{9}{2}}\xi _L+{\displaystyle \frac{3}{4}}\xi _L^2]+h^3[+{\displaystyle \frac{77035}{144}}{\displaystyle \frac{4}{3}}n_f\zeta _35n_f\xi _L`$ (40)
$`{\displaystyle \frac{1913}{27}}n_f+{\displaystyle \frac{100}{81}}n_f^2+18\zeta _3\xi _L54\zeta _3+{\displaystyle \frac{1647}{16}}\xi _L+{\displaystyle \frac{549}{16}}\xi _L^2+{\displaystyle \frac{81}{16}}\xi _L^3]`$
$`+h^4[+{\displaystyle \frac{42735097}{1728}}{\displaystyle \frac{333}{2}}n_f\zeta _3\xi _L+{\displaystyle \frac{3}{4}}n_f\zeta _3\xi _L^2+{\displaystyle \frac{37579}{108}}n_f\zeta _3+{\displaystyle \frac{2320}{9}}n_f\zeta _5{\displaystyle \frac{6759}{16}}n_f\xi _L`$
$`{\displaystyle \frac{757}{16}}n_f\xi _L^2{\displaystyle \frac{386759}{81}}n_f+{\displaystyle \frac{16}{3}}n_f^2\zeta _3\xi _L+{\displaystyle \frac{28}{9}}n_f^2\zeta _3+{\displaystyle \frac{38}{9}}n_f^2\xi _L+{\displaystyle \frac{199903}{972}}n_f^2`$
$`{\displaystyle \frac{1000}{729}}n_f^3+{\displaystyle \frac{20067}{16}}\zeta _3\xi _L+{\displaystyle \frac{2511}{32}}\zeta _3\xi _L^2{\displaystyle \frac{225}{8}}\zeta _3\xi _L^3{\displaystyle \frac{117}{32}}\zeta _3\xi _L^4{\displaystyle \frac{35385}{8}}\zeta _3`$
$`{\displaystyle \frac{8325}{16}}\zeta _5\xi _L{\displaystyle \frac{3195}{32}}\zeta _5\xi _L^2+{\displaystyle \frac{255}{16}}\zeta _5\xi _L^3+{\displaystyle \frac{315}{64}}\zeta _5\xi _L^4{\displaystyle \frac{66765}{64}}\zeta _5+{\displaystyle \frac{281059}{64}}\xi _L`$
$`+{\displaystyle \frac{43311}{32}}\xi _L^2+{\displaystyle \frac{20925}{64}}\xi _L^3+{\displaystyle \frac{2547}{64}}\xi _L^4].`$
### 6.3 Subtracting the Triple Gluon Vertex I
Subtracting just the Lorentz structure that is present at the tree level the gluon vertex gives the following renormalization condition
$$T_1^{\stackrel{~}{\text{mom}}\mathrm{g}}(\mu ^2)=1.$$
(41)
This subtracts the triple gluon vertex with the zero momentum gluon and one of the others being polarized longitudinal (in direction of their momenta) and the last being polarized transversal and parallel to the plane defined by the gluon momenta (for details see ). A similar derivation as for the ghost gluon vertex leads then to the following relation of the coupling in the $`\stackrel{~}{\text{mom}}\mathrm{g}`$ scheme to the $`\overline{\text{ms}}`$ coupling:
$$h^{\stackrel{~}{\text{mom}}\mathrm{g}}(\mu )=h^{\overline{\text{ms}}}(\mu )\frac{\left(T_1^{\overline{\text{ms}}}(\mu ^2)\right)^2}{\left(1+\mathrm{\Pi }^{\overline{\text{ms}}}(\mu ^2)\right)^3}.$$
(42)
This gives the following relations between the coupling constants of the two schemes:
$`h^{\stackrel{~}{\text{mom}}\mathrm{g}}`$ $`=`$ $`h+h^2[+{\displaystyle \frac{169}{12}}{\displaystyle \frac{10}{9}}n_f+{\displaystyle \frac{9}{2}}\xi _L+{\displaystyle \frac{3}{4}}\xi _L^2]+h^3[+{\displaystyle \frac{38261}{72}}{\displaystyle \frac{4}{3}}n_f\zeta _35n_f\xi _L`$ (43)
$`{\displaystyle \frac{7571}{108}}n_f+{\displaystyle \frac{100}{81}}n_f^2+{\displaystyle \frac{81}{4}}\zeta _3\xi _L{\displaystyle \frac{207}{4}}\zeta _3+{\displaystyle \frac{795}{8}}\xi _L+{\displaystyle \frac{549}{16}}\xi _L^2+{\displaystyle \frac{81}{16}}\xi _L^3]`$
$`+h^4[+{\displaystyle \frac{84464417}{3456}}{\displaystyle \frac{1375}{8}}n_f\zeta _3\xi _L+{\displaystyle \frac{3}{4}}n_f\zeta _3\xi _L^2+{\displaystyle \frac{72161}{216}}n_f\zeta _3+{\displaystyle \frac{2320}{9}}n_f\zeta _5`$
$`{\displaystyle \frac{39197}{96}}n_f\xi _L{\displaystyle \frac{187}{4}}n_f\xi _L^2{\displaystyle \frac{6098639}{1296}}n_f+{\displaystyle \frac{16}{3}}n_f^2\zeta _3\xi _L+{\displaystyle \frac{28}{9}}n_f^2\zeta _3+{\displaystyle \frac{38}{9}}n_f^2\xi _L`$
$`+{\displaystyle \frac{197149}{972}}n_f^2{\displaystyle \frac{1000}{729}}n_f^3+{\displaystyle \frac{87729}{64}}\zeta _3\xi _L+{\displaystyle \frac{6261}{64}}\zeta _3\xi _L^2{\displaystyle \frac{1371}{64}}\zeta _3\xi _L^3{\displaystyle \frac{117}{32}}\zeta _3\xi _L^4`$
$`{\displaystyle \frac{279861}{64}}\zeta _3{\displaystyle \frac{26805}{64}}\zeta _5\xi _L{\displaystyle \frac{4155}{64}}\zeta _5\xi _L^2+{\displaystyle \frac{885}{64}}\zeta _5\xi _L^3+{\displaystyle \frac{315}{64}}\zeta _5\xi _L^4{\displaystyle \frac{3285}{4}}\zeta _5`$
$`+{\displaystyle \frac{267071}{64}}\xi _L+{\displaystyle \frac{83511}{64}}\xi _L^2+{\displaystyle \frac{41445}{128}}\xi _L^3+{\displaystyle \frac{2547}{64}}\xi _L^4].`$
### 6.4 Subtracting the Triple Gluon Vertex II
The $`\stackrel{~}{\text{mom}}\mathrm{gg}`$ scheme used in and is defined by the following renormalization conditions:
$$T_1^{\stackrel{~}{\text{mom}}\mathrm{gg}}(\mu ^2)\frac{1}{2}T_2^{\stackrel{~}{\text{mom}}\mathrm{gg}}(\mu ^2)=1.$$
(44)
This subtracts the triple gluon vertex with the zero momentum gluon being polarized transversal and parallel to the plane defined by the gluon momenta. The other two gluon are also polarized transversal but perpendicular to the vertex plane. Another possibility to understand this linear combination is to realize that it corresponds to subtracting the transversal part of the 3-point function, which is the only one to survive at the ZP (see also ):
$$G_{\mu \nu \rho }^{(3)abc}(p,p,0)=(g_{\mu \nu }\frac{p_\mu p_\nu }{p^2})p_\rho G^{(3)}(p^2)$$
(45)
Inserting eqs. (3,2,7) shows that indeed:
$`G_{\mu _1\mu _2\mu _3}^{(3)abc}(p,p,0)`$ $`=`$ $`(g_{\mu _1\mu _2}{\displaystyle \frac{p_{\mu _1}p_{\mu _2}}{p^2}})p_{\mu _3}{\displaystyle \frac{\left(T_1(p^2)\frac{1}{2}T_2(p^2)\right)}{(1+\mathrm{\Pi }^{(2)}(p^2))^2(1+\mathrm{\Pi }^{(2)}(0))}}.`$ (46)
From this follows the relation of the coupling in the $`\stackrel{~}{\text{mom}}\mathrm{gg}`$ scheme to the $`\overline{\text{ms}}`$ coupling:
$$h^{\stackrel{~}{\text{mom}}\mathrm{gg}}(\mu )=h^{\overline{\text{ms}}}(\mu )\frac{\left(T_1^{\overline{\text{ms}}}(\mu ^2)\frac{1}{2}T_2^{\overline{\text{ms}}}(\mu ^2)\right)^2}{\left(1+\mathrm{\Pi }^{\overline{\text{ms}}}(\mu ^2)\right)^3}.$$
(47)
Inserting $`T_1`$, $`T_2`$ and $`\mathrm{\Pi }`$ gives
$`h^{\stackrel{~}{\text{mom}}\mathrm{gg}}`$ $`=`$ $`h+h^2[+{\displaystyle \frac{70}{3}}{\displaystyle \frac{22}{9}}n_f]+h^3[+{\displaystyle \frac{516217}{576}}{\displaystyle \frac{4}{3}}n_f\zeta _3{\displaystyle \frac{7}{4}}n_f\xi _L+{\displaystyle \frac{1}{2}}n_f\xi _L^2`$
$`{\displaystyle \frac{8125}{54}}n_f+{\displaystyle \frac{376}{81}}n_f^2+{\displaystyle \frac{63}{4}}\zeta _3\xi _L{\displaystyle \frac{153}{4}}\zeta _3+{\displaystyle \frac{225}{8}}\xi _L+{\displaystyle \frac{249}{32}}\xi _L^2+{\displaystyle \frac{45}{16}}\xi _L^3+{\displaystyle \frac{9}{64}}\xi _L^4]`$
$`+h^4[+{\displaystyle \frac{304676635}{6912}}{\displaystyle \frac{1197}{8}}n_f\zeta _3\xi _L{\displaystyle \frac{21}{8}}n_f\zeta _3\xi _L^2+{\displaystyle \frac{13339}{27}}n_f\zeta _315n_f\zeta _5\xi _L`$
$`+{\displaystyle \frac{1885}{9}}n_f\zeta _5{\displaystyle \frac{16663}{96}}n_f\xi _L{\displaystyle \frac{23}{48}}n_f\xi _L^2{\displaystyle \frac{3}{32}}n_f\xi _L^3{\displaystyle \frac{7}{32}}n_f\xi _L^4{\displaystyle \frac{13203725}{1296}}n_f`$
$`+{\displaystyle \frac{8}{3}}n_f^2\zeta _3\xi _L+{\displaystyle \frac{40}{9}}n_f^2\zeta _3+{\displaystyle \frac{44}{9}}n_f^2\xi _L{\displaystyle \frac{8}{9}}n_f^2\xi _L^2+{\displaystyle \frac{580495}{972}}n_f^2{\displaystyle \frac{5680}{729}}n_f^3`$
$`+{\displaystyle \frac{138171}{64}}\zeta _3\xi _L{\displaystyle \frac{63}{64}}\zeta _3\xi _L^2+{\displaystyle \frac{423}{64}}\zeta _3\xi _L^3{\displaystyle \frac{27}{32}}\zeta _3\xi _L^4{\displaystyle \frac{299961}{64}}\zeta _3{\displaystyle \frac{47295}{64}}\zeta _5\xi _L`$
$`+{\displaystyle \frac{3825}{64}}\zeta _5\xi _L^2{\displaystyle \frac{1665}{64}}\zeta _5\xi _L^3{\displaystyle \frac{81825}{64}}\zeta _5+{\displaystyle \frac{188523}{128}}\xi _L+{\displaystyle \frac{118591}{256}}\xi _L^2+{\displaystyle \frac{5175}{32}}\xi _L^3`$
$`+{\displaystyle \frac{5829}{256}}\xi _L^4{\displaystyle \frac{27}{64}}\xi _L^5{\displaystyle \frac{27}{256}}\xi _L^6].`$
Following these examples and using the results in the appendix, it should be possible to construct the relations to relate the coupling constant in any $`\stackrel{~}{\text{mom}}`$ renormalization scheme to the $`\overline{\text{ms}}`$ one.
## 7 The $`\stackrel{~}{\text{mom}}`$ $`\beta `$-functions
The first two coefficients of the $`\beta `$-function in any massless renormalization scheme are scheme independent. When considering renormalization schemes in which the coupling constant depends on the generic covariant gauge parameter, this statement is not true, since the additional gauge dependence spoils these universality just as a mass parameter . This means that we need not only to take into account the gauge parameter when writing down renormalization group equations, but also that there is a difference between the gauge parameter defined in the $`\overline{\text{ms}}`$ scheme and in any of the $`\stackrel{~}{\text{mom}}`$ schemes we are considering:
$$\xi _L^{\overline{\text{ms}}}=\frac{Z_3^{\stackrel{~}{\text{mom}}}}{Z_3^{\overline{\text{ms}}}}\xi _L^{\stackrel{~}{\text{mom}}}.$$
(49)
In Landau gauge we do not need to consider the additional $`\mu `$ dependence introduced by the gauge parameter nor the difference between $`\xi _L^{\overline{\text{ms}}}`$ and $`\xi _L^{\stackrel{~}{\text{mom}}}`$. So in Landau gauge the $`\beta `$-function in any $`\stackrel{~}{\text{mom}}`$ like scheme can be obtained in the following simple way once the $`\stackrel{~}{\text{mom}}`$ coupling is expressed as a series in the $`\overline{\text{ms}}`$ coupling
$$\beta ^{\stackrel{~}{\text{mom}}}(h^{\stackrel{~}{\text{mom}}})=\mu ^2\frac{\mathrm{d}h^{\stackrel{~}{\text{mom}}}}{\mathrm{d}\mu ^2}=\frac{h^{\stackrel{~}{\text{mom}}}}{h^{\overline{\text{ms}}}}\mu ^2\frac{\mathrm{d}h^{\overline{\text{ms}}}}{\mathrm{d}\mu ^2}=\frac{h^{\stackrel{~}{\text{mom}}}}{h^{\overline{\text{ms}}}}\beta ^{\overline{\text{ms}}}(h^{\overline{\text{ms}}}).$$
(50)
where one has to invert the series $`h^{\stackrel{~}{\text{mom}}}(h^{\overline{\text{ms}}})`$ and insert this series on the right hand side to express the $`\stackrel{~}{\text{mom}}`$ $`\beta `$-function as a series in $`h^{\stackrel{~}{\text{mom}}}`$. The four-loop QCD $`\beta `$-function in the $`\overline{\text{ms}}`$ scheme was obtained in and reads:
$`\beta ^{\overline{\text{ms}}}`$ $`=`$ $`h^2\left[11+{\displaystyle \frac{2}{3}}n_f\right]+h^3\left[102+{\displaystyle \frac{38}{3}}n_f\right]+h^4\left[{\displaystyle \frac{2857}{2}}+{\displaystyle \frac{5033}{18}}n_f{\displaystyle \frac{325}{54}}n_f^2\right]`$ (51)
$`+h^5[{\displaystyle \frac{149753}{6}}+{\displaystyle \frac{6508}{27}}n_f\zeta _3+{\displaystyle \frac{1078361}{162}}n_f{\displaystyle \frac{6472}{81}}n_f^2\zeta _3{\displaystyle \frac{50065}{162}}n_f^2`$
$`{\displaystyle \frac{1093}{729}}n_f^33564\zeta _3],`$
Combining this result with the relations between the coupling constants, we can finally obtain the first four coefficients of the $`\beta `$-functions in the above introduced four $`\stackrel{~}{\text{mom}}`$ schemes, where we use the definition:
$$\beta (h)=\underset{i=0}{\overset{3}{}}h^{(i+2)}\beta _i$$
(52)
and in Landau gauge $`\beta _0`$ and $`\beta _1`$ are independent of the renormalization scheme. The three- and four-loop contribution in Landau gauge for the four schemes under consideration are:
$`\beta _3^{\stackrel{~}{\text{mom}}\mathrm{g}}`$ $`=`$ $`\left[+{\displaystyle \frac{58491}{16}}{\displaystyle \frac{2277}{4}}\zeta _3\right]+n_f\left[{\displaystyle \frac{15283}{24}}+{\displaystyle \frac{119}{6}}\zeta _3\right]+n_f^2\left[+{\displaystyle \frac{481}{27}}+{\displaystyle \frac{8}{9}}\zeta _3\right],`$ (53)
$`\beta _4^{\stackrel{~}{\text{mom}}\mathrm{g}}`$ $`=`$ $`\left[+{\displaystyle \frac{10982273}{64}}{\displaystyle \frac{1425171}{32}}\zeta _3{\displaystyle \frac{36135}{2}}\zeta _5\right]+n_f\left[{\displaystyle \frac{3830167}{96}}+{\displaystyle \frac{1075423}{144}}\zeta _3+{\displaystyle \frac{60895}{9}}\zeta _5\right]`$ (54)
$`+n_f^2\left[+{\displaystyle \frac{724445}{324}}{\displaystyle \frac{12959}{54}}\zeta _3{\displaystyle \frac{9280}{27}}\zeta _5\right]+n_f^3\left[{\displaystyle \frac{788}{27}}+{\displaystyle \frac{16}{9}}\zeta _3\right],`$
$`\beta _3^{\stackrel{~}{\text{mom}}}`$ $`=`$ $`\left[+{\displaystyle \frac{28965}{8}}{\displaystyle \frac{3861}{8}}\zeta _3\right]+n_f\left[{\displaystyle \frac{7715}{12}}+{\displaystyle \frac{175}{12}}\zeta _3\right]+n_f^2\left[+{\displaystyle \frac{989}{54}}+{\displaystyle \frac{8}{9}}\zeta _3\right],`$ (55)
$`\beta _4^{\stackrel{~}{\text{mom}}}`$ $`=`$ $`\left[+{\displaystyle \frac{1380469}{8}}{\displaystyle \frac{625317}{16}}\zeta _3{\displaystyle \frac{772695}{32}}\zeta _5\right]+n_f\left[{\displaystyle \frac{970819}{24}}+{\displaystyle \frac{516881}{72}}\zeta _3+{\displaystyle \frac{1027375}{144}}\zeta _5\right]`$ (56)
$`+n_f^2\left[+{\displaystyle \frac{736541}{324}}{\displaystyle \frac{6547}{27}}\zeta _3{\displaystyle \frac{9280}{27}}\zeta _5\right]+n_f^3\left[{\displaystyle \frac{800}{27}}+{\displaystyle \frac{16}{9}}\zeta _3\right],`$
$`\beta _3^{\stackrel{~}{\text{mom}}\mathrm{q}}`$ $`=`$ $`\left[+{\displaystyle \frac{29559}{8}}594\zeta _3\right]+n_f\left[{\displaystyle \frac{7769}{12}}+{\displaystyle \frac{64}{3}}\zeta _3\right]+n_f^2\left[+{\displaystyle \frac{989}{54}}+{\displaystyle \frac{8}{9}}\zeta _3\right],`$ (57)
$`\beta _4^{\stackrel{~}{\text{mom}}\mathrm{q}}`$ $`=`$ $`\left[+{\displaystyle \frac{2795027}{16}}{\displaystyle \frac{174207}{4}}\zeta _3{\displaystyle \frac{734415}{32}}\zeta _5\right]+n_f\left[{\displaystyle \frac{487751}{12}}+{\displaystyle \frac{67939}{9}}\zeta _3+{\displaystyle \frac{1016935}{144}}\zeta _5\right]`$ (58)
$`+n_f^2\left[+{\displaystyle \frac{737837}{324}}{\displaystyle \frac{6709}{27}}\zeta _3{\displaystyle \frac{9280}{27}}\zeta _5\right]+n_f^3\left[{\displaystyle \frac{800}{27}}+{\displaystyle \frac{16}{9}}\zeta _3\right],`$
$`\beta _3^{\stackrel{~}{\text{mom}}\mathrm{gg}}`$ $`=`$ $`\left[+{\displaystyle \frac{186747}{64}}{\displaystyle \frac{1683}{4}}\zeta _3\right]+n_f\left[{\displaystyle \frac{35473}{96}}+{\displaystyle \frac{65}{6}}\zeta _3\right]+n_f^2\left[{\displaystyle \frac{829}{54}}+{\displaystyle \frac{8}{9}}\zeta _3\right]+n_f^3\left[+{\displaystyle \frac{8}{9}}\right]`$ (59)
,
$`\beta _4^{\stackrel{~}{\text{mom}}\mathrm{gg}}`$ $`=`$ $`\left[+{\displaystyle \frac{20783939}{128}}{\displaystyle \frac{1300563}{32}}\zeta _3{\displaystyle \frac{900075}{32}}\zeta _5\right]`$ (60)
$`+n_f\left[{\displaystyle \frac{2410799}{64}}+{\displaystyle \frac{1323259}{144}}\zeta _3+{\displaystyle \frac{908995}{144}}\zeta _5\right]`$
$`+n_f^2\left[+{\displaystyle \frac{1464379}{648}}{\displaystyle \frac{12058}{27}}\zeta _3{\displaystyle \frac{7540}{27}}\zeta _5\right]`$
$`+n_f^3\left[{\displaystyle \frac{3164}{27}}+{\displaystyle \frac{64}{9}}\zeta _3\right]+n_f^4\left[+{\displaystyle \frac{320}{81}}\right].`$
## 8 Discussion
In numerical form the coupling constant relations and the corresponding $`\beta `$-functions for these four schemes in Landau gauge read:
$`h^{\stackrel{~}{\text{mom}}}`$ $`=`$ $`h+h^2\left[+14.08331.11111n_f\right]+h^3\left[+475.47572.4546n_f+1.23457n_f^2\right]`$ (61)
$`+h^4\left[+18652.44109.72n_f+209.401n_f^21.37174n_f^3\right],`$
$`h^{\stackrel{~}{\text{mom}}\mathrm{q}}`$ $`=`$ $`h+h^2\left[+14.08331.11111n_f\right]+h^3\left[+470.05472.4546n_f+1.23457n_f^2\right]`$ (62)
$`+h^4\left[+18332.44089.25n_f+209.401n_f^21.37174n_f^3\right],`$
$`h^{\stackrel{~}{\text{mom}}\mathrm{g}}`$ $`=`$ $`h+h^2\left[+14.08331.11111n_f\right]+h^3\left[+469.19671.7046n_f+1.23457n_f^2\right]`$ (63)
$`+h^4\left[+183324036.86n_f+206.568n_f^21.37174n_f^3\right],`$
$`h^{\stackrel{~}{\text{mom}}\mathrm{gg}}`$ $`=`$ $`h+h^2\left[+23.33332.44444n_f\right]`$ (64)
$`+h^3\left[+850.231152.066n_f+4.64198n_f^2\right]`$
$`+h^4\left[+37119.79377.02n_f+602.56n_f^27.7915n_f^3\right],`$
$`\beta ^{\stackrel{~}{\text{mom}}}`$ $`=`$ $`\left(h^{\stackrel{~}{\text{mom}}}\right)^2\left[11+0.666667n_f\right]+\left(h^{\stackrel{~}{\text{mom}}}\right)^3\left[102+12.6667n_f\right]`$ (65)
$`+\left(h^{\stackrel{~}{\text{mom}}}\right)^4\left[3040.48+625.387n_f19.3833n_f^2\right]`$
$`+\left(h^{\stackrel{~}{\text{mom}}}\right)^5\left[100541+24423.3n_f1625.4n_f^2+27.4926n_f^3\right],`$
$`\beta ^{\stackrel{~}{\text{mom}}\mathrm{g}}`$ $`=`$ $`\left(h^{\stackrel{~}{\text{mom}}\mathrm{g}}\right)^2\left[11+0.666667n_f\right]+\left(h^{\stackrel{~}{\text{mom}}\mathrm{g}}\right)^3\left[102+12.6667n_f\right]`$ (66)
$`+\left(h^{\stackrel{~}{\text{mom}}\mathrm{g}}\right)^4\left[2971.42+612.951n_f18.8833n_f^2\right]`$
$`+\left(h^{\stackrel{~}{\text{mom}}\mathrm{g}}\right)^5\left[99327.8+23904.4n_f1591.07n_f^2+27.0482n_f^3\right],`$
$`\beta ^{\stackrel{~}{\text{mom}}\mathrm{q}}`$ $`=`$ $`\left(h^{\stackrel{~}{\text{mom}}\mathrm{q}}\right)^2\left[11+0.666667n_f\right]+\left(h^{\stackrel{~}{\text{mom}}\mathrm{q}}\right)^3\left[102+12.6667n_f\right]`$ (67)
$`+\left(h^{\stackrel{~}{\text{mom}}\mathrm{q}}\right)^4\left[2980.85+621.773n_f19.3833n_f^2\right]`$
$`+\left(h^{\stackrel{~}{\text{mom}}\mathrm{q}}\right)^5\left[98539.5+24249n_f1622.19n_f^2+27.4926n_f^3\right],`$
$`\beta ^{\stackrel{~}{\text{mom}}\mathrm{gg}}`$ $`=`$ $`\left(h^{\stackrel{~}{\text{mom}}\mathrm{gg}}\right)^2\left[11+0.666667n_f\right]+\left(h^{\stackrel{~}{\text{mom}}\mathrm{gg}}\right)^3\left[102+12.6667n_f\right]`$
$`+\left(h^{\stackrel{~}{\text{mom}}\mathrm{gg}}\right)^4\left[2412.16+356.488n_f+14.2834n_f^20.888889n_f^3\right]`$
$`+\left(h^{\stackrel{~}{\text{mom}}\mathrm{gg}}\right)^5\left[84353.8+20077.1n_f1433.44n_f^2+108.637n_f^33.95062n_f^4\right].`$
The corresponding numbers for the $`\overline{\text{ms}}`$ $`\beta `$-function are (as an expansion in $`h=h^{\overline{\text{ms}}}`$):
$`\beta ^{\overline{\text{ms}}}`$ $`=`$ $`h^2\left[11+0.666667n_f\right]+h^3\left[102+12.6667n_f\right]`$ (69)
$`+h^4\left[1428.5+279.611n_f6.01852n_f^2\right]`$
$`+h^5\left[29243+6946.29n_f405.089n_f^21.49931n_f^3\right].`$
It has been argued that momentum subtraction schemes will lead to couplings that are only slightly different. Comparing the $`\stackrel{~}{\text{mom}}\mathrm{gg}`$ scheme with the other three schemes it is easy to see that this of course depends a lot on the linear combination used in the renormalization condition. On the other hand it is clearly seen that the couplings as well as the $`\beta `$-functions of the three schemes that are identical at one-loop order will also give very close-by values for the coupling constant at higher orders. Also the $`\stackrel{~}{\text{mom}}\mathrm{gg}`$ scheme is still much closer to the other $`\stackrel{~}{\text{mom}}`$ schemes than to the $`\overline{\text{ms}}`$ one.
Inserting the values $`1,\mathrm{},6`$ for $`n_f`$ shows that the coefficients for all the values of $`n_f`$ are smaller for the $`\overline{\text{ms}}`$ scheme than in any of the $`\stackrel{~}{\text{mom}}`$ schemes. The difference between the $`\overline{\text{ms}}`$ scheme and the $`\stackrel{~}{\text{mom}}`$ schemes is only small for $`n_f=6`$ and gets quite large for the small values of $`n_f`$.
Recently the momentum dependence (running) of the three-gluon asymmetrical vertex corresponding to the combination $`T_{gg}=T_1(p^2)\frac{1}{2}T_2`$ was computed within the lattice approach in . The resulting (nonperturbative!) behaviour of
$$\alpha _s(p^2)=4\pi h^{\stackrel{~}{\text{mom}}\mathrm{gg}}(p^2)$$
in the flavorless QCD has been found to be best described by an ansatz<sup>5</sup><sup>5</sup>5To be consistent to we use the Euclidean metrics below.
$$\alpha _s(p^2)=\alpha _s^{Pert}(p^2)+\frac{c}{p^2},$$
(70)
with
$$\mathrm{\Lambda }_{\overline{\text{ms}}}=237\pm 3\text{MeV},c=0.63\pm 0.03\text{GeV}^2.$$
(71)
Here $`\mathrm{\Lambda }_{\overline{\text{ms}}}=\mathrm{e}xp(70/66)\mathrm{\Lambda }_{\stackrel{~}{\text{mom}}\mathrm{gg}}`$. The above results have been obtained by working at three-loop level and employing the $`\stackrel{~}{\text{mom}}\mathrm{gg}`$ scheme. The authors of have also investigated the dependence of the result on the (then unknown) four-loop contribution to $`\beta ^{\stackrel{~}{\text{mom}}\mathrm{gg}}`$. Their findings are summarized in Table 1, which we have copied from with adding one more row obtained with the help of the linear extrapolation and corresponding to the true value of the parameter
$$\frac{b_3}{b_2}=\frac{\beta _3^{\stackrel{~}{\text{mom}}\mathrm{gg}}/(4\pi )}{\beta _2^{\stackrel{~}{\text{mom}}\mathrm{gg}}}=2.78284\mathrm{}\text{(for }n_f=0\text{)}.$$
Table 1 clearly demonstrates that taking into account the four-loop term in the $`\beta ^{\stackrel{~}{\text{mom}}\mathrm{gg}}`$-function leads to a significant decrease (around 30%) of the value of the non-perturbative $`1/p^2`$ correction.
## 9 Conclusions
In this paper we have analytically computed the full set of the three-loops propagators and fundamental three-linear vertexes with one of three external momenta set to zero in the massless QCD. The results have been used to find the NNNLO conversion factors transforming the $`\overline{\mathrm{MS}}`$ coupling constant to the ones defined corresponding to a set of regularization independent renormalization schemes based on momentum subtractions ($`\stackrel{~}{\text{mom}}`$-schemes). Then we have used the conversion factors to evaluate the four-loop $`\beta `$ functions in these schemes.
The newly computed corrections to the coupling constant running in $`\stackrel{~}{\text{mom}}`$-schemes prove to be numerically significant. They should be taken into account when confronting the running obtained with the help of lattice simulations with the pQCD predictions.
Acknowledgments
K. G. Ch. is grateful A. I. Davydychev and J.H. Kühn for useful discussions. A. R. wants to thank Thorsten Seidensticker, Timo van Ritbergen and Dominik Stöckinger for valuable discussions. Thorsten Seidensticker also provided some preliminary results of work in progress for crosschecks. This work was supported by the DFG under Contract Ku 502/8-1 (DFG-Forschergruppe “Quantenfeldtheorie, Computeralgebra und Monte-Carlo-Simulationen” ) and the Graduiertenkolleg “Elementarteilchenphysik an Beschleunigern” at the Universität Karlsruhe.
## Appendix A Technical Remarks
To evaluate the self-energies and vertex function at three loop order we made heavy use of computer programs. The program QGRAF was used to generate the diagrams. For the self-energies and the vertex functions at the ZP we arrive at scalar integrals of the massless propagator type after applying some projectors. An algorithm which allows one to analytically evaluate divergent as well as finite parts of such integrals was elaborated in . It has been implemented in an efficient way in the MINCER package written for the symbolic manipulation program FORM . The huge number of diagrams in the three-loop calculations requires a complete automation of the whole procedure which should include also the calculation of the color factors. This has been implemented and is described in more detail along with other similar installations in . This setup also gives the user tools to perform series expansions in small parameters, a feature that was used to calculate the naive first order expansion in a small external momentum for the ghost vertex which is necessary to check the WST identity as described in Section 4. We had to evaluate the following diagrams:
| Number of diagrams for | 1 loops | 2 loops | 3 loops | approximate runtime per function |
| --- | --- | --- | --- | --- |
| $`\stackrel{~}{\mathrm{\Pi }}`$,$`\mathrm{\Sigma }_V`$ | 1 | 7 | 106 | 6 hours |
| $`\mathrm{\Pi }`$ | 4 | 27 | 494 | 12 hours |
| $`2\times \stackrel{~}{\mathrm{\Gamma }}_i+4\times a_i^{()}`$ | 2 | 40 | 1022 | 12 hours |
| $`2\times (\mathrm{\Lambda }_i,\mathrm{\Lambda }_i^T)`$ | 2 | 40 | 1022 | 2 days |
| $`T_i`$ | 10 | 189 | 5526 | 4 weeks |
All calculations were done on workstations using the Alpha 21164 processors running at 600 MHz. The approximate runtimes are only rough estimates and show that the calculation of the triple gluon vertex was by far the most demanding part. This is not only because of the number of diagrams but also because the complicated vertex and propagator structures generate a huge number of intermediate terms.
Wherever possible, we have cross-checked the following expressions with the known two- and three-loop results, the gluon, ghost and quark field anomalous dimensions can also be found in and . In most cases we found complete agreement. We only find different values than for the two-loop contributions to $`\mathrm{\Lambda }_g`$ and $`\mathrm{\Lambda }_g^T`$ (their $`\mathrm{\Gamma }_3`$ and $`\mathrm{\Gamma }_4`$). For the two-loop triple gluon vertex we find complete agreement with , which contains a list of miss-prints in earlier publications, among them the gauge dependence of the one-loop $`T_2`$ in .
The latex code for all results in this paper has been created automatically. The code has then been retransformed to form input files and cross-checked against the original expressions. We hope that by doing these checks and avoiding any hand editing of the formulas we have circumvent the common problem of miss-prints in otherwise correct results. The price is of course that the layout of some of the formulas is not as fancy as it could be. Also note that these expressions will be made available in the World-Wide-Web.
## Appendix B Conventions
### B.1 Color Factors
The following results are given for generic color factors and are valid for any semi-simple compact Lie group, where $`T_{ij}^a`$ are the generators of the fundamental representation and $`f^{abc}`$ are the structure constants of the Lie algebra.
| | Definition | $`\mathrm{SU}(N)`$ | QCD | QED |
| --- | --- | --- | --- | --- |
| $`C_A`$ | $`f^{acd}f^{bcd}=C_A\delta ^{ab}`$ | $`N`$ | $`3`$ | 0 |
| $`C_F`$ | $`[T^aT^a]_{ij}=C_F\delta _{ij}`$ | $`\frac{N^21}{2N}`$ | $`\frac{4}{3}`$ | $`1`$ |
| $`T`$ | $`\mathrm{Tr}(T^aT^b)=T\delta ^{ab}`$ | $`\frac{1}{2}`$ | $`\frac{1}{2}`$ | $`1`$ |
### B.2 Momentum Dependence
We give the $`\overline{\text{ms}}`$ renormalized result for all self-energies and vertex functions at the point $`p^2=\mu ^2`$. The correct $`p^2`$ dependence for any of the self-energies and vertex functions $`\mathrm{\Gamma }`$ can be restored from the renormalization group equations, which are with our conventions:
$$\mu ^2\frac{\mathrm{d}}{\mathrm{d}\mu ^2}\mathrm{\Gamma }(h,q,\mu ,\xi _L)=\left(\mu ^2\frac{}{\mu ^2}+\beta \frac{}{h}+\gamma _3\xi _L\frac{}{\xi _L}\right)\mathrm{\Gamma }(h,q,\mu ,\xi _L)=\gamma _\mathrm{\Gamma }\mathrm{\Gamma }(h,q,\mu ,\xi _L)$$
(B.1)
For massless QCD one cam write
$$\mathrm{\Gamma }=\mathrm{\Gamma }_0+\underset{i>0}{}\mathrm{\Gamma }_i\mathrm{log}^i(\frac{\mu ^2}{q^2}),$$
(B.2)
and as a direct consequence of (B.1) one gets (for $`n>0)`$
$$\mathrm{\Gamma }_n=\left(\beta \frac{}{h}\gamma _3\xi _L\frac{}{\xi _L}+\gamma _\mathrm{\Gamma }\right)\frac{1}{n}\mathrm{\Gamma }_{n1}$$
(B.3)
The equation can be directly used to reconstruct the $`q`$-dependence of the scalar functions $`1+\mathrm{\Pi }`$, $`1+\stackrel{~}{\mathrm{\Pi }}`$ and $`1+\mathrm{\Sigma }_V`$. For the vertex functions one needs to multiply the functions $`T_i`$, $`\stackrel{~}{\mathrm{\Gamma }}_i`$ and $`\mathrm{\Lambda }_i`$ by $`g=4\pi \sqrt{h}`$. The anomalous dimensions of any $`\mathrm{\Gamma }`$ is:
$$\gamma _\mathrm{\Gamma }=\left(\frac{n_g}{2}\gamma _3+\frac{n_h}{2}\stackrel{~}{\gamma }_3+\frac{n_q}{2}\gamma _2\right)$$
(B.4)
where $`n_g`$,$`n_h`$ and $`n_q`$ are the number of external gluon, ghost and quark fields, respectively, of the diagrams contributing to $`\mathrm{\Gamma }`$. The QCD $`\beta `$-function and the field anomalous dimensions are given in part D of this appendix. Note that in our definitions for the anomalous dimensions a relative sign and a factor 2 is different from many other publications.
## Appendix C Propagators and Vertices in massless QCD
### C.1 The Gluon Self-Energy
$`\mathrm{\Pi }^{\overline{\text{ms}}}`$ $`=`$ $`hC_A\left[{\displaystyle \frac{97}{36}}{\displaystyle \frac{1}{2}}\xi _L{\displaystyle \frac{1}{4}}\xi _L^2\right]+hTn_f\left[+{\displaystyle \frac{20}{9}}\right]`$ (C.1)
$`+h^2C_ATn_f\left[+{\displaystyle \frac{59}{4}}+8\zeta _3{\displaystyle \frac{10}{9}}\xi _L{\displaystyle \frac{10}{9}}\xi _L^2\right]+h^2C_FTn_f\left[+{\displaystyle \frac{55}{3}}16\zeta _3\right]`$
$`+h^2C_A^2\left[{\displaystyle \frac{2381}{96}}2\zeta _3\xi _L+3\zeta _3+{\displaystyle \frac{463}{288}}\xi _L+{\displaystyle \frac{95}{144}}\xi _L^2{\displaystyle \frac{1}{16}}\xi _L^3+{\displaystyle \frac{1}{16}}\xi _L^4\right]`$
$`+h^3C_A^3[{\displaystyle \frac{10221367}{31104}}{\displaystyle \frac{12071}{288}}\zeta _3\xi _L{\displaystyle \frac{161}{96}}\zeta _3\xi _L^2+{\displaystyle \frac{149}{96}}\zeta _3\xi _L^3+{\displaystyle \frac{13}{96}}\zeta _3\xi _L^4+{\displaystyle \frac{1549}{24}}\zeta _3+{\displaystyle \frac{3}{8}}\zeta _4\xi _L`$
$`+{\displaystyle \frac{3}{32}}\zeta _4\xi _L^2+{\displaystyle \frac{9}{32}}\zeta _4+{\displaystyle \frac{115}{8}}\zeta _5\xi _L+{\displaystyle \frac{385}{96}}\zeta _5\xi _L^2{\displaystyle \frac{5}{24}}\zeta _5\xi _L^3{\displaystyle \frac{35}{192}}\zeta _5\xi _L^4+{\displaystyle \frac{7025}{192}}\zeta _5`$
$`+{\displaystyle \frac{13141}{1152}}\xi _L+{\displaystyle \frac{30835}{10368}}\xi _L^2{\displaystyle \frac{2813}{1152}}\xi _L^3{\displaystyle \frac{29}{48}}\xi _L^4+{\displaystyle \frac{1}{16}}\xi _L^5{\displaystyle \frac{1}{64}}\xi _L^6]`$
$`+h^3C_A^2Tn_f[+{\displaystyle \frac{1154561}{3888}}+{\displaystyle \frac{202}{9}}\zeta _3\xi _L{\displaystyle \frac{25}{6}}\zeta _3\xi _L^2+{\displaystyle \frac{871}{18}}\zeta _39\zeta _4{\displaystyle \frac{160}{3}}\zeta _5{\displaystyle \frac{241}{24}}\xi _L`$
$`{\displaystyle \frac{6137}{1296}}\xi _L^2{\displaystyle \frac{5}{18}}\xi _L^3+{\displaystyle \frac{5}{12}}\xi _L^4]`$
$`+h^3C_AT^2n_f^2\left[{\displaystyle \frac{10499}{243}}{\displaystyle \frac{64}{9}}\zeta _3\xi _L{\displaystyle \frac{256}{9}}\zeta _3+{\displaystyle \frac{16}{9}}\xi _L{\displaystyle \frac{100}{81}}\xi _L^2\right]`$
$`+h^3C_AC_FTn_f\left[+{\displaystyle \frac{96809}{324}}+8\zeta _3\xi _L+8\zeta _3\xi _L^2{\displaystyle \frac{1492}{9}}\zeta _3+12\zeta _480\zeta _5{\displaystyle \frac{55}{6}}\xi _L{\displaystyle \frac{55}{6}}\xi _L^2\right]`$
$`+h^3C_FT^2n_f^2\left[{\displaystyle \frac{7402}{81}}+{\displaystyle \frac{608}{9}}\zeta _3\right]+h^3C_F^2Tn_f\left[{\displaystyle \frac{286}{9}}{\displaystyle \frac{296}{3}}\zeta _3+160\zeta _5\right].`$
### C.2 The Ghost Self-Energy
$`\stackrel{~}{\mathrm{\Pi }}^{\overline{\text{ms}}}`$ $`=`$ $`hC_A\left[1\right]+h^2C_ATn_f\left[+{\displaystyle \frac{95}{24}}\right]`$ (C.2)
$`+h^2C_A^2\left[{\displaystyle \frac{1751}{192}}{\displaystyle \frac{3}{8}}\zeta _3\xi _L+{\displaystyle \frac{3}{16}}\zeta _3\xi _L^2+{\displaystyle \frac{15}{16}}\zeta _3+{\displaystyle \frac{7}{64}}\xi _L{\displaystyle \frac{3}{8}}\xi _L^2\right]`$
$`+h^3C_A^3[{\displaystyle \frac{466373}{3888}}{\displaystyle \frac{1429}{192}}\zeta _3\xi _L+{\displaystyle \frac{93}{64}}\zeta _3\xi _L^2+{\displaystyle \frac{23}{64}}\zeta _3\xi _L^3+{\displaystyle \frac{12403}{576}}\zeta _3{\displaystyle \frac{3}{16}}\zeta _4\xi _L{\displaystyle \frac{3}{64}}\zeta _4\xi _L^2`$
$`{\displaystyle \frac{9}{64}}\zeta _4+{\displaystyle \frac{65}{32}}\zeta _5\xi _L{\displaystyle \frac{35}{32}}\zeta _5\xi _L^2+{\displaystyle \frac{5}{32}}\zeta _5\xi _L^3+{\displaystyle \frac{65}{32}}\zeta _5{\displaystyle \frac{61}{576}}\xi _L{\displaystyle \frac{327}{128}}\xi _L^2{\displaystyle \frac{361}{384}}\xi _L^3]`$
$`+h^3C_A^2Tn_f\left[+{\displaystyle \frac{150247}{1944}}+{\displaystyle \frac{13}{6}}\zeta _3\xi _L+{\displaystyle \frac{29}{9}}\zeta _3+{\displaystyle \frac{9}{2}}\zeta _4{\displaystyle \frac{199}{144}}\xi _L\right]`$
$`+h^3C_AT^2n_f^2\left[{\displaystyle \frac{5161}{486}}{\displaystyle \frac{8}{9}}\zeta _3\right]+h^3C_AC_FTn_f\left[+{\displaystyle \frac{899}{24}}22\zeta _36\zeta _4\right].`$
### C.3 The Quark Self-Energy
$`\mathrm{\Sigma }_V^{\overline{\text{ms}}}`$ $`=`$ $`hC_F\left[+\xi _L\right]+h^2C_AC_F\left[+{\displaystyle \frac{41}{4}}3\zeta _3\xi _L3\zeta _3+{\displaystyle \frac{13}{2}}\xi _L+{\displaystyle \frac{9}{8}}\xi _L^2\right]`$ (C.3)
$`+h^2C_FTn_f\left[{\displaystyle \frac{7}{2}}\right]+h^2C_F^2\left[{\displaystyle \frac{5}{8}}\right]`$
$`+h^3C_A^2C_F[+{\displaystyle \frac{159257}{648}}35\zeta _3\xi _L{\displaystyle \frac{39}{8}}\zeta _3\xi _L^2{\displaystyle \frac{1}{3}}\zeta _3\xi _L^3{\displaystyle \frac{3139}{24}}\zeta _3+{\displaystyle \frac{3}{8}}\zeta _4\xi _L+{\displaystyle \frac{3}{16}}\zeta _4\xi _L^2`$
$`{\displaystyle \frac{69}{16}}\zeta _4+{\displaystyle \frac{5}{2}}\zeta _5\xi _L+{\displaystyle \frac{5}{4}}\zeta _5\xi _L^2+{\displaystyle \frac{165}{4}}\zeta _5+{\displaystyle \frac{39799}{576}}\xi _L+{\displaystyle \frac{787}{64}}\xi _L^2+{\displaystyle \frac{55}{24}}\xi _L^3]`$
$`+h^3C_AC_F^2\left[{\displaystyle \frac{997}{24}}17\zeta _3\xi _L+\zeta _3\xi _L^3+44\zeta _3+6\zeta _4+20\zeta _5\xi _L20\zeta _5+4\xi _L+{\displaystyle \frac{3}{2}}\xi _L^2{\displaystyle \frac{1}{8}}\xi _L^3\right]`$
$`+h^3C_AC_FTn_f\left[{\displaystyle \frac{11887}{81}}+8\zeta _3\xi _L+{\displaystyle \frac{52}{3}}\zeta _3{\displaystyle \frac{1723}{72}}\xi _L\right]+h^3C_FT^2n_f^2\left[+{\displaystyle \frac{1570}{81}}\right]`$
$`+h^3C_F^2Tn_f\left[{\displaystyle \frac{79}{6}}+16\zeta _3{\displaystyle \frac{3}{2}}\xi _L\right]+h^3C_F^3\left[{\displaystyle \frac{73}{12}}{\displaystyle \frac{2}{3}}\zeta _3\xi _L^3+{\displaystyle \frac{7}{8}}\xi _L\right].`$
### C.4 The Triple Gluon Vertex
$`T_1^{\overline{\text{ms}}}`$ $`=`$ $`1+hC_A\left[{\displaystyle \frac{61}{36}}{\displaystyle \frac{1}{4}}\xi _L^2\right]+hTn_f\left[+{\displaystyle \frac{20}{9}}\right]`$ (C.4)
$`+h^2C_A^2\left[{\displaystyle \frac{9907}{576}}{\displaystyle \frac{15}{8}}\zeta _3\xi _L+{\displaystyle \frac{13}{8}}\zeta _3+{\displaystyle \frac{153}{64}}\xi _L+{\displaystyle \frac{35}{36}}\xi _L^2{\displaystyle \frac{3}{16}}\xi _L^3+{\displaystyle \frac{1}{16}}\xi _L^4\right]`$
$`+h^2C_ATn_f\left[+{\displaystyle \frac{955}{72}}+8\zeta _3{\displaystyle \frac{10}{9}}\xi _L^2\right]+h^2C_FTn_f\left[+{\displaystyle \frac{55}{3}}16\zeta _3\right]`$
$`+h^3C_A^3[{\displaystyle \frac{155555}{648}}{\displaystyle \frac{48047}{1152}}\zeta _3\xi _L{\displaystyle \frac{1481}{384}}\zeta _3\xi _L^2+{\displaystyle \frac{145}{128}}\zeta _3\xi _L^3+{\displaystyle \frac{13}{96}}\zeta _3\xi _L^4+{\displaystyle \frac{68537}{1152}}\zeta _3`$
$`+{\displaystyle \frac{9}{16}}\zeta _4\xi _L+{\displaystyle \frac{9}{64}}\zeta _4\xi _L^2+{\displaystyle \frac{27}{64}}\zeta _4+{\displaystyle \frac{2345}{128}}\zeta _5\xi _L+{\displaystyle \frac{2995}{384}}\zeta _5\xi _L^2+{\displaystyle \frac{55}{128}}\zeta _5\xi _L^3{\displaystyle \frac{35}{192}}\zeta _5\xi _L^4`$
$`+{\displaystyle \frac{2065}{128}}\zeta _5+{\displaystyle \frac{44555}{2304}}\xi _L+{\displaystyle \frac{127487}{20736}}\xi _L^2{\displaystyle \frac{271}{192}}\xi _L^3{\displaystyle \frac{149}{192}}\xi _L^4+{\displaystyle \frac{3}{32}}\xi _L^5{\displaystyle \frac{1}{64}}\xi _L^6]`$
$`+h^3C_A^2C_F[2+{\displaystyle \frac{43}{32}}\zeta _3\xi _L+{\displaystyle \frac{101}{32}}\zeta _3\xi _L^2+{\displaystyle \frac{17}{32}}\zeta _3\xi _L^3{\displaystyle \frac{1497}{32}}\zeta _3{\displaystyle \frac{325}{32}}\zeta _5\xi _L{\displaystyle \frac{215}{32}}\zeta _5\xi _L^2`$
$`{\displaystyle \frac{35}{32}}\zeta _5\xi _L^3+{\displaystyle \frac{1695}{32}}\zeta _5{\displaystyle \frac{1}{2}}\xi _L+{\displaystyle \frac{1}{4}}\xi _L^2]+h^3C_FT^2n_f^2[{\displaystyle \frac{7402}{81}}+{\displaystyle \frac{608}{9}}\zeta _3]`$
$`+h^3C_A^2Tn_f[+{\displaystyle \frac{682607}{2592}}+{\displaystyle \frac{1889}{72}}\zeta _3\xi _L{\displaystyle \frac{25}{6}}\zeta _3\xi _L^2+{\displaystyle \frac{3733}{72}}\zeta _3{\displaystyle \frac{27}{2}}\zeta _4{\displaystyle \frac{160}{3}}\zeta _5{\displaystyle \frac{521}{96}}\xi _L`$
$`{\displaystyle \frac{9511}{2592}}\xi _L^2{\displaystyle \frac{5}{6}}\xi _L^3+{\displaystyle \frac{5}{12}}\xi _L^4]+h^3C_F^2Tn_f[{\displaystyle \frac{286}{9}}{\displaystyle \frac{296}{3}}\zeta _3+160\zeta _5]`$
$`+h^3C_AC_FTn_f\left[+{\displaystyle \frac{181711}{648}}+8\zeta _3\xi _L^2{\displaystyle \frac{1438}{9}}\zeta _3+18\zeta _480\zeta _5{\displaystyle \frac{55}{6}}\xi _L^2\right]`$
$`+h^3C_AT^2n_f^2\left[{\displaystyle \frac{3415}{81}}{\displaystyle \frac{64}{9}}\zeta _3\xi _L{\displaystyle \frac{248}{9}}\zeta _3+{\displaystyle \frac{16}{9}}\xi _L{\displaystyle \frac{100}{81}}\xi _L^2\right].`$
$`T_2^{\overline{\text{ms}}}`$ $`=`$ $`hC_A\left[{\displaystyle \frac{37}{12}}+{\displaystyle \frac{3}{2}}\xi _L+{\displaystyle \frac{1}{4}}\xi _L^2\right]+hTn_f\left[+{\displaystyle \frac{8}{3}}\right]`$ (C.5)
$`+h^2C_A^2\left[{\displaystyle \frac{443}{24}}+{\displaystyle \frac{1}{2}}\zeta _3\xi _L{\displaystyle \frac{3}{2}}\zeta _3+{\displaystyle \frac{31}{48}}\xi _L+{\displaystyle \frac{59}{72}}\xi _L^2{\displaystyle \frac{11}{16}}\xi _L^3{\displaystyle \frac{1}{8}}\xi _L^4\right]`$
$`+h^2C_ATn_f\left[+{\displaystyle \frac{91}{6}}+{\displaystyle \frac{5}{2}}\xi _L{\displaystyle \frac{2}{9}}\xi _L^2\right]+h^2C_FTn_f\left[+8\right]`$
$`+h^3C_A^3[{\displaystyle \frac{2108863}{6912}}{\displaystyle \frac{787}{36}}\zeta _3\xi _L{\displaystyle \frac{1093}{96}}\zeta _3\xi _L^2{\displaystyle \frac{95}{48}}\zeta _3\xi _L^3+{\displaystyle \frac{7}{48}}\zeta _3\xi _L^4{\displaystyle \frac{19255}{96}}\zeta _3`$
$`+{\displaystyle \frac{3805}{96}}\zeta _5\xi _L+{\displaystyle \frac{845}{48}}\zeta _5\xi _L^2+{\displaystyle \frac{355}{96}}\zeta _5\xi _L^3+{\displaystyle \frac{35}{192}}\zeta _5\xi _L^4+{\displaystyle \frac{46985}{192}}\zeta _5+{\displaystyle \frac{9869}{2304}}\xi _L`$
$`+{\displaystyle \frac{178183}{20736}}\xi _L^2+{\displaystyle \frac{7039}{2304}}\xi _L^3{\displaystyle \frac{57}{128}}\xi _L^4+{\displaystyle \frac{5}{32}}\xi _L^5+{\displaystyle \frac{3}{64}}\xi _L^6]`$
$`+h^3C_A^2C_F[+{\displaystyle \frac{53}{2}}+{\displaystyle \frac{245}{8}}\zeta _3\xi _L+{\displaystyle \frac{49}{2}}\zeta _3\xi _L^2{\displaystyle \frac{5}{8}}\zeta _3\xi _L^3{\displaystyle \frac{9}{16}}\zeta _3\xi _L^4+{\displaystyle \frac{7249}{16}}\zeta _3{\displaystyle \frac{125}{2}}\zeta _5\xi _L`$
$`50\zeta _5\xi _L^25\zeta _5\xi _L^3{\displaystyle \frac{1025}{2}}\zeta _5{\displaystyle \frac{61}{4}}\xi _L+{\displaystyle \frac{3}{4}}\xi _L^3]`$
$`+h^3C_A^2Tn_f[+{\displaystyle \frac{279701}{864}}+{\displaystyle \frac{5}{18}}\zeta _3\xi _L+{\displaystyle \frac{19}{4}}\zeta _3\xi _L^2+{\displaystyle \frac{749}{12}}\zeta _3+{\displaystyle \frac{10}{3}}\zeta _5\xi _L{\displaystyle \frac{110}{3}}\zeta _5+{\displaystyle \frac{197}{12}}\xi _L`$
$`+{\displaystyle \frac{3637}{2592}}\xi _L^2{\displaystyle \frac{131}{36}}\xi _L^3{\displaystyle \frac{1}{3}}\xi _L^4]`$
$`+h^3C_AC_FTn_f\left[+{\displaystyle \frac{1105}{6}}24\zeta _3\xi _L8\zeta _3\xi _L^2208\zeta _3+{\displaystyle \frac{320}{3}}\zeta _5+25\xi _L+{\displaystyle \frac{31}{6}}\xi _L^2\right]`$
$`+h^3C_AT^2n_f^2\left[{\displaystyle \frac{1741}{27}}+{\displaystyle \frac{32}{9}}\zeta _3\xi _L16\zeta _3{\displaystyle \frac{14}{9}}\xi _L{\displaystyle \frac{140}{81}}\xi _L^2\right]`$
$`+h^3C_FT^2n_f^2\left[{\displaystyle \frac{176}{3}}+{\displaystyle \frac{128}{3}}\zeta _3\right]+h^3C_F^2Tn_f\left[4\right].`$
### C.5 The Ghost Gluon Vertex
$`\stackrel{~}{\mathrm{\Gamma }}_\mathrm{h}^{\overline{\text{ms}}}`$ $`=`$ $`1+hC_A\left[+{\displaystyle \frac{1}{2}}\xi _L\right]+h^2C_A^2\left[{\displaystyle \frac{9}{16}}\zeta _3\xi _L+{\displaystyle \frac{3}{16}}\zeta _3\xi _L^2+{\displaystyle \frac{43}{16}}\xi _L+{\displaystyle \frac{7}{16}}\xi _L^2\right]`$ (C.6)
$`+h^3C_A^3[+{\displaystyle \frac{725}{192}}\zeta _3\xi _L+{\displaystyle \frac{267}{64}}\zeta _3\xi _L^22\zeta _3\xi _L^3{\displaystyle \frac{105}{8}}\zeta _5\xi _L{\displaystyle \frac{95}{16}}\zeta _5\xi _L^2+{\displaystyle \frac{35}{16}}\zeta _5\xi _L^3`$
$`+{\displaystyle \frac{3631}{128}}\xi _L+{\displaystyle \frac{2089}{384}}\xi _L^2+{\displaystyle \frac{17}{24}}\xi _L^3]+h^3C_A^2Tn_f[+{\displaystyle \frac{29}{12}}\zeta _3\xi _L{\displaystyle \frac{493}{48}}\xi _L]`$
$`+h^3C_A^2C_F\left[27\zeta _3\xi _L6\zeta _3\xi _L^2+{\displaystyle \frac{9}{2}}\zeta _3\xi _L^3+{\displaystyle \frac{225}{8}}\zeta _5\xi _L+{\displaystyle \frac{75}{8}}\zeta _5\xi _L^2{\displaystyle \frac{15}{4}}\zeta _5\xi _L^3\right].`$
$`\stackrel{~}{\mathrm{\Gamma }}_\mathrm{g}^{\overline{\text{ms}}}`$ $`=`$ $`1+hC_A\left[+{\displaystyle \frac{3}{4}}+{\displaystyle \frac{1}{4}}\xi _L\right]`$ (C.7)
$`+h^2C_A^2\left[+{\displaystyle \frac{599}{96}}{\displaystyle \frac{9}{16}}\zeta _3\xi _L+{\displaystyle \frac{3}{16}}\zeta _3\xi _L^2+{\displaystyle \frac{97}{32}}\xi _L+{\displaystyle \frac{1}{4}}\xi _L^2\right]+h^2C_ATn_f\left[{\displaystyle \frac{29}{12}}\right]`$
$`+h^3C_A^3[+{\displaystyle \frac{43273}{432}}{\displaystyle \frac{817}{192}}\zeta _3\xi _L+{\displaystyle \frac{293}{64}}\zeta _3\xi _L^2+{\displaystyle \frac{3}{64}}\zeta _3\xi _L^3+{\displaystyle \frac{783}{64}}\zeta _3{\displaystyle \frac{85}{16}}\zeta _5\xi _L{\displaystyle \frac{325}{64}}\zeta _5\xi _L^2`$
$`+{\displaystyle \frac{5}{16}}\zeta _5\xi _L^3{\displaystyle \frac{875}{64}}\zeta _5+{\displaystyle \frac{28039}{768}}\xi _L+{\displaystyle \frac{4091}{768}}\xi _L^2+{\displaystyle \frac{31}{48}}\xi _L^3]`$
$`+h^3C_A^2C_F[+{\displaystyle \frac{27}{4}}{\displaystyle \frac{249}{16}}\zeta _3\xi _L{\displaystyle \frac{111}{16}}\zeta _3\xi _L^2+{\displaystyle \frac{3}{16}}\zeta _3\xi _L^3{\displaystyle \frac{639}{16}}\zeta _3+{\displaystyle \frac{165}{8}}\zeta _5\xi _L+{\displaystyle \frac{15}{2}}\zeta _5\xi _L^2`$
$`+{\displaystyle \frac{225}{8}}\zeta _5{\displaystyle \frac{21}{4}}\xi _L+{\displaystyle \frac{3}{4}}\xi _L^2]+h^3C_AC_FTn_f[16+12\zeta _3]`$
$`+h^3C_A^2Tn_f\left[{\displaystyle \frac{15143}{216}}+{\displaystyle \frac{55}{24}}\zeta _3\xi _L{\displaystyle \frac{49}{8}}\zeta _3{\displaystyle \frac{357}{32}}\xi _L\right]+h^3C_AT^2n_f^2\left[+{\displaystyle \frac{280}{27}}\right].`$
### C.6 The Quark Gluon Vertex
$`\mathrm{\Lambda }_\mathrm{q}^{\overline{\text{ms}}}`$ $`=`$ $`1+hC_A\left[+1+{\displaystyle \frac{1}{2}}\xi _L\right]+hC_F\left[+\xi _L\right]`$ (C.8)
$`+h^2C_A^2\left[+{\displaystyle \frac{2015}{192}}{\displaystyle \frac{3}{2}}\zeta _3+{\displaystyle \frac{181}{64}}\xi _L+{\displaystyle \frac{13}{16}}\xi _L^2\right]+h^2C_ATn_f\left[{\displaystyle \frac{95}{24}}\right]`$
$`+h^2C_AC_F\left[+{\displaystyle \frac{41}{4}}3\zeta _3\xi _L3\zeta _3+{\displaystyle \frac{15}{2}}\xi _L+{\displaystyle \frac{13}{8}}\xi _L^2\right]+h^2C_FTn_f\left[{\displaystyle \frac{7}{2}}\right]+h^2C_F^2\left[{\displaystyle \frac{5}{8}}\right]`$
$`+h^3C_A^3[+{\displaystyle \frac{2255345}{15552}}{\displaystyle \frac{119}{192}}\zeta _3\xi _L{\displaystyle \frac{19}{32}}\zeta _3\xi _L^2{\displaystyle \frac{37}{64}}\zeta _3\xi _L^3{\displaystyle \frac{1817}{72}}\zeta _3+{\displaystyle \frac{3}{16}}\zeta _4\xi _L+{\displaystyle \frac{3}{64}}\zeta _4\xi _L^2`$
$`+{\displaystyle \frac{9}{64}}\zeta _4{\displaystyle \frac{105}{32}}\zeta _5\xi _L+{\displaystyle \frac{5}{32}}\zeta _5\xi _L^2+{\displaystyle \frac{15}{32}}\zeta _5\xi _L^3{\displaystyle \frac{335}{32}}\zeta _5+{\displaystyle \frac{6511}{192}}\xi _L+{\displaystyle \frac{3295}{384}}\xi _L^2`$
$`+{\displaystyle \frac{235}{128}}\xi _L^3]+h^3C_A^2Tn_f[{\displaystyle \frac{169525}{1944}}{\displaystyle \frac{1}{6}}\zeta _3\xi _L{\displaystyle \frac{2}{9}}\zeta _3{\displaystyle \frac{9}{2}}\zeta _4{\displaystyle \frac{491}{48}}\xi _L]`$
$`+h^3C_A^2C_F[+{\displaystyle \frac{331393}{1296}}41\zeta _3\xi _L{\displaystyle \frac{45}{8}}\zeta _3\xi _L^2+{\displaystyle \frac{5}{12}}\zeta _3\xi _L^3{\displaystyle \frac{3631}{24}}\zeta _3+{\displaystyle \frac{3}{8}}\zeta _4\xi _L+{\displaystyle \frac{3}{16}}\zeta _4\xi _L^2`$
$`{\displaystyle \frac{69}{16}}\zeta _4+{\displaystyle \frac{35}{8}}\zeta _5\xi _L+{\displaystyle \frac{5}{4}}\zeta _5\xi _L^2{\displaystyle \frac{5}{8}}\zeta _5\xi _L^3+{\displaystyle \frac{125}{2}}\zeta _5+{\displaystyle \frac{13117}{144}}\xi _L+{\displaystyle \frac{311}{16}}\xi _L^2+{\displaystyle \frac{11}{3}}\xi _L^3]`$
$`+h^3C_AC_F^2[{\displaystyle \frac{253}{6}}17\zeta _3\xi _L+\zeta _3\xi _L^3+44\zeta _3+6\zeta _4+20\zeta _5\xi _L20\zeta _5+{\displaystyle \frac{59}{16}}\xi _L+{\displaystyle \frac{3}{2}}\xi _L^2`$
$`{\displaystyle \frac{1}{8}}\xi _L^3]+h^3C_AC_FTn_f[{\displaystyle \frac{121637}{648}}+8\zeta _3\xi _L+{\displaystyle \frac{118}{3}}\zeta _3+6\zeta _4{\displaystyle \frac{1067}{36}}\xi _L]`$
$`+h^3C_AT^2n_f^2\left[+{\displaystyle \frac{5161}{486}}+{\displaystyle \frac{8}{9}}\zeta _3\right]+h^3C_FT^2n_f^2\left[+{\displaystyle \frac{1570}{81}}\right]`$
$`+h^3C_F^2Tn_f\left[{\displaystyle \frac{79}{6}}+16\zeta _3{\displaystyle \frac{3}{2}}\xi _L\right]+h^3C_F^3\left[{\displaystyle \frac{73}{12}}{\displaystyle \frac{2}{3}}\zeta _3\xi _L^3+{\displaystyle \frac{7}{8}}\xi _L\right].`$
$`\mathrm{\Lambda }_\mathrm{q}^{T\overline{\text{ms}}}`$ $`=`$ $`hC_A\left[+{\displaystyle \frac{9}{4}}\xi _L{\displaystyle \frac{1}{4}}\xi _L^2\right]+hC_F\left[2\right]`$ (C.9)
$`+h^2C_A^2\left[+{\displaystyle \frac{523}{24}}{\displaystyle \frac{11}{8}}\zeta _3\xi _L{\displaystyle \frac{1}{16}}\zeta _3\xi _L^2+{\displaystyle \frac{11}{16}}\zeta _3{\displaystyle \frac{145}{48}}\xi _L{\displaystyle \frac{215}{144}}\xi _L^2{\displaystyle \frac{1}{16}}\xi _L^3+{\displaystyle \frac{1}{16}}\xi _L^4\right]`$
$`+h^2C_AC_F\left[{\displaystyle \frac{505}{18}}+{\displaystyle \frac{13}{4}}\xi _L{\displaystyle \frac{1}{2}}\xi _L^2{\displaystyle \frac{1}{4}}\xi _L^3\right]+h^2C_FTn_f\left[+{\displaystyle \frac{52}{9}}\right]`$
$`+h^2C_ATn_f\left[{\displaystyle \frac{16}{3}}4\zeta _3{\displaystyle \frac{2}{3}}\xi _L{\displaystyle \frac{5}{9}}\xi _L^2\right]+h^2C_F^2\left[+92\xi _L\right]`$
$`+h^3C_A^3[+{\displaystyle \frac{2844547}{6912}}{\displaystyle \frac{2221}{96}}\zeta _3\xi _L{\displaystyle \frac{203}{192}}\zeta _3\xi _L^2+{\displaystyle \frac{89}{48}}\zeta _3\xi _L^3+{\displaystyle \frac{29}{192}}\zeta _3\xi _L^4+{\displaystyle \frac{565}{9}}\zeta _3`$
$`+{\displaystyle \frac{655}{48}}\zeta _5\xi _L+{\displaystyle \frac{385}{96}}\zeta _5\xi _L^2{\displaystyle \frac{35}{48}}\zeta _5\xi _L^3{\displaystyle \frac{35}{192}}\zeta _5\xi _L^4{\displaystyle \frac{23575}{192}}\zeta _5{\displaystyle \frac{268897}{6912}}\xi _L`$
$`{\displaystyle \frac{40489}{2304}}\xi _L^2{\displaystyle \frac{10445}{2304}}\xi _L^3{\displaystyle \frac{155}{576}}\xi _L^4+{\displaystyle \frac{1}{16}}\xi _L^5{\displaystyle \frac{1}{64}}\xi _L^6]`$
$`+h^3C_A^2C_F[{\displaystyle \frac{1450057}{2592}}+{\displaystyle \frac{389}{48}}\zeta _3\xi _L+{\displaystyle \frac{35}{12}}\zeta _3\xi _L^2{\displaystyle \frac{1}{16}}\zeta _3\xi _L^3{\displaystyle \frac{1939}{8}}\zeta _3{\displaystyle \frac{265}{24}}\zeta _5\xi _L`$
$`+{\displaystyle \frac{5}{24}}\zeta _5\xi _L^2+{\displaystyle \frac{5}{8}}\zeta _5\xi _L^3+{\displaystyle \frac{8665}{24}}\zeta _5+{\displaystyle \frac{10469}{288}}\xi _L{\displaystyle \frac{299}{96}}\xi _L^2{\displaystyle \frac{149}{36}}\xi _L^3{\displaystyle \frac{15}{32}}\xi _L^4+{\displaystyle \frac{1}{16}}\xi _L^5]`$
$`+h^3C_A^2Tn_f[{\displaystyle \frac{171143}{864}}+{\displaystyle \frac{275}{18}}\zeta _3\xi _L{\displaystyle \frac{3}{2}}\zeta _3\xi _L^2{\displaystyle \frac{880}{9}}\zeta _3+{\displaystyle \frac{20}{3}}\zeta _5+{\displaystyle \frac{415}{72}}\xi _L{\displaystyle \frac{575}{288}}\xi _L^2`$
$`+{\displaystyle \frac{1}{36}}\xi _L^3+{\displaystyle \frac{5}{18}}\xi _L^4]+h^3C_AT^2n_f^2[+{\displaystyle \frac{700}{27}}{\displaystyle \frac{32}{9}}\zeta _3\xi _L+{\displaystyle \frac{112}{9}}\zeta _3+{\displaystyle \frac{40}{27}}\xi _L]`$
$`+h^3C_AC_FTn_f\left[+{\displaystyle \frac{78625}{648}}4\zeta _3\xi _L+4\zeta _3\xi _L^2+{\displaystyle \frac{196}{3}}\zeta _3+{\displaystyle \frac{400}{3}}\zeta _5{\displaystyle \frac{73}{18}}\xi _L{\displaystyle \frac{113}{24}}\xi _L^2{\displaystyle \frac{5}{9}}\xi _L^3\right]`$
$`+h^3C_AC_F^2\left[+{\displaystyle \frac{76339}{288}}+6\zeta _3\xi _L+316\zeta _3{\displaystyle \frac{1240}{3}}\zeta _5{\displaystyle \frac{3235}{72}}\xi _L{\displaystyle \frac{107}{32}}\xi _L^2+{\displaystyle \frac{1}{2}}\xi _L^3\right]`$
$`+h^3C_FT^2n_f^2\left[{\displaystyle \frac{2000}{81}}\right]+h^3C_F^2Tn_f\left[+{\displaystyle \frac{821}{9}}{\displaystyle \frac{160}{3}}\zeta _3160\zeta _5+{\displaystyle \frac{52}{9}}\xi _L\right]`$
$`+h^3C_F^3\left[{\displaystyle \frac{973}{12}}{\displaystyle \frac{496}{3}}\zeta _3+{\displaystyle \frac{640}{3}}\zeta _5+9\xi _L\right].`$
$`\mathrm{\Lambda }_\mathrm{g}^{\overline{\text{ms}}}`$ $`=`$ $`1+hC_A\left[+{\displaystyle \frac{1}{4}}+{\displaystyle \frac{3}{4}}\xi _L\right]+hC_F\left[\xi _L\right]`$ (C.10)
$`+h^2C_A^2\left[+{\displaystyle \frac{1075}{192}}3\zeta _3+{\displaystyle \frac{181}{64}}\xi _L+\xi _L^2\right]+h^2C_AC_F\left[+{\displaystyle \frac{3}{4}}3\zeta _3\xi _L3\zeta _3{\displaystyle \frac{3}{2}}\xi _L{\displaystyle \frac{3}{8}}\xi _L^2\right]`$
$`+h^2C_ATn_f\left[{\displaystyle \frac{55}{24}}\right]+h^2C_FTn_f\left[+{\displaystyle \frac{1}{2}}\right]+h^2C_F^2\left[+{\displaystyle \frac{19}{8}}2\xi _L^2\right]`$
$`+h^3C_A^3[+{\displaystyle \frac{59815}{1944}}{\displaystyle \frac{1235}{48}}\zeta _3\xi _L{\displaystyle \frac{73}{32}}\zeta _3\xi _L^2{\displaystyle \frac{1}{4}}\zeta _3\xi _L^3{\displaystyle \frac{10145}{288}}\zeta _3+{\displaystyle \frac{3}{16}}\zeta _4\xi _L+{\displaystyle \frac{3}{64}}\zeta _4\xi _L^2`$
$`+{\displaystyle \frac{9}{64}}\zeta _4+{\displaystyle \frac{35}{4}}\zeta _5\xi _L{\displaystyle \frac{5}{16}}\zeta _5\xi _L^2+{\displaystyle \frac{365}{16}}\zeta _5+{\displaystyle \frac{5695}{128}}\xi _L+{\displaystyle \frac{4033}{384}}\xi _L^2+{\displaystyle \frac{141}{64}}\xi _L^3]`$
$`+h^3C_A^2C_F[+{\displaystyle \frac{21971}{648}}+{\displaystyle \frac{55}{2}}\zeta _3\xi _L{\displaystyle \frac{5}{4}}\zeta _3\xi _L^2+{\displaystyle \frac{5}{12}}\zeta _3\xi _L^3{\displaystyle \frac{527}{3}}\zeta _3+{\displaystyle \frac{3}{8}}\zeta _4\xi _L+{\displaystyle \frac{3}{16}}\zeta _4\xi _L^2{\displaystyle \frac{69}{16}}\zeta _4`$
$`{\displaystyle \frac{35}{2}}\zeta _5\xi _L+{\displaystyle \frac{5}{2}}\zeta _5\xi _L^2+140\zeta _5{\displaystyle \frac{13117}{288}}\xi _L{\displaystyle \frac{159}{16}}\xi _L^2{\displaystyle \frac{173}{96}}\xi _L^3]`$
$`+h^3C_A^2Tn_f\left[{\displaystyle \frac{89777}{3888}}+{\displaystyle \frac{17}{6}}\zeta _3\xi _L+{\displaystyle \frac{359}{18}}\zeta _3{\displaystyle \frac{9}{2}}\zeta _420\zeta _514\xi _L\right]`$
$`+h^3C_AC_F^2[+{\displaystyle \frac{1259}{32}}17\zeta _3\xi _L+{\displaystyle \frac{9}{2}}\zeta _3\xi _L^2+{\displaystyle \frac{1}{2}}\zeta _3\xi _L^3+82\zeta _3+6\zeta _4+20\zeta _5\xi _L100\zeta _5`$
$`{\displaystyle \frac{747}{32}}\xi _L{\displaystyle \frac{29}{2}}\xi _L^2{\displaystyle \frac{23}{8}}\xi _L^3]+h^3C_AT^2n_f^2[+{\displaystyle \frac{1741}{486}}+{\displaystyle \frac{8}{9}}\zeta _3]`$
$`+h^3C_AC_FTn_f\left[{\displaystyle \frac{6823}{324}}2\zeta _3\xi _L+{\displaystyle \frac{58}{3}}\zeta _3+6\zeta _4+{\displaystyle \frac{761}{72}}\xi _L\right]+h^3C_FT^2n_f^2\left[{\displaystyle \frac{302}{81}}\right]`$
$`+h^3C_F^2Tn_f\left[{\displaystyle \frac{45}{2}}+16\zeta _3+{\displaystyle \frac{19}{2}}\xi _L\right]+h^3C_F^3\left[{\displaystyle \frac{109}{12}}{\displaystyle \frac{2}{3}}\zeta _3\xi _L^3+{\displaystyle \frac{41}{8}}\xi _L\right].`$
$`\mathrm{\Lambda }_\mathrm{g}^{T\overline{\text{ms}}}`$ $`=`$ $`hC_A\left[+{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{2}}\xi _L\right]+hC_F\left[+2\xi _L\right]+h^2C_A^2\left[+{\displaystyle \frac{167}{36}}{\displaystyle \frac{1}{2}}\zeta _3\xi _L{\displaystyle \frac{5}{2}}\zeta _3{\displaystyle \frac{29}{24}}\xi _L{\displaystyle \frac{5}{8}}\xi _L^2\right]`$ (C.11)
$`+h^2C_AC_F\left[+{\displaystyle \frac{43}{6}}+4\zeta _3+9\xi _L+{\displaystyle \frac{3}{2}}\xi _L^2\right]+h^2C_ATn_f\left[{\displaystyle \frac{7}{9}}\right]+h^2C_FTn_f\left[4\right]`$
$`+h^2C_F^2\left[3+2\xi _L^2\right]+h^3C_AC_FTn_f\left[{\displaystyle \frac{15827}{108}}+8\zeta _3\xi _L{\displaystyle \frac{40}{9}}\zeta _3{\displaystyle \frac{112}{3}}\xi _L\right]`$
$`+h^3C_A^3[+{\displaystyle \frac{1017553}{10368}}+{\displaystyle \frac{857}{72}}\zeta _3\xi _L+{\displaystyle \frac{7}{16}}\zeta _3\xi _L^2{\displaystyle \frac{7373}{144}}\zeta _3{\displaystyle \frac{115}{12}}\zeta _5\xi _L{\displaystyle \frac{205}{12}}\zeta _5{\displaystyle \frac{36727}{1728}}\xi _L`$
$`{\displaystyle \frac{2309}{384}}\xi _L^2{\displaystyle \frac{43}{32}}\xi _L^3]`$
$`+h^3C_A^2C_F[+{\displaystyle \frac{39011}{216}}63\zeta _3\xi _L{\displaystyle \frac{9}{4}}\zeta _3\xi _L^2+{\displaystyle \frac{3917}{36}}\zeta _3+25\zeta _5\xi _L{\displaystyle \frac{305}{3}}\zeta _5+{\displaystyle \frac{4091}{32}}\xi _L`$
$`+{\displaystyle \frac{2263}{96}}\xi _L^2+{\displaystyle \frac{63}{16}}\xi _L^3]`$
$`+h^3C_A^2Tn_f\left[{\displaystyle \frac{61085}{1296}}{\displaystyle \frac{2}{9}}\zeta _3\xi _L{\displaystyle \frac{14}{9}}\zeta _3+{\displaystyle \frac{80}{3}}\zeta _5+{\displaystyle \frac{3197}{432}}\xi _L\right]+h^3C_F^3\left[+3{\displaystyle \frac{17}{4}}\xi _L\right]`$
$`+h^3C_AC_F^2\left[{\displaystyle \frac{3515}{48}}+6\zeta _3\xi _L6\zeta _3\xi _L^2{\displaystyle \frac{184}{3}}\zeta _3+{\displaystyle \frac{280}{3}}\zeta _5+{\displaystyle \frac{361}{16}}\xi _L+17\xi _L^2+{\displaystyle \frac{11}{4}}\xi _L^3\right]`$
$`+h^3C_AT^2n_f^2\left[+{\displaystyle \frac{260}{81}}\right]+h^3C_FT^2n_f^2\left[+{\displaystyle \frac{208}{9}}\right]+h^3C_F^2Tn_f\left[+{\displaystyle \frac{28}{3}}11\xi _L\right].`$
## Appendix D Renormalization Group Coefficients
### D.1 The $`\beta `$-function
$`\beta `$ $`=`$ $`h^2\left[{\displaystyle \frac{11}{3}}C_A+{\displaystyle \frac{4}{3}}Tn_f\right]+h^3\left[+{\displaystyle \frac{20}{3}}C_ATn_f{\displaystyle \frac{34}{3}}C_A^2+4C_FTn_f\right]`$ (D.1)
$`+h^4[+{\displaystyle \frac{205}{9}}C_AC_FTn_f{\displaystyle \frac{158}{27}}C_AT^2n_f^2+{\displaystyle \frac{1415}{27}}C_A^2Tn_f{\displaystyle \frac{2857}{54}}C_A^3`$
$`{\displaystyle \frac{44}{9}}C_FT^2n_f^22C_F^2Tn_f].`$
### D.2 The Gluon Field Anomalous Dimension
$`\gamma _3`$ $`=`$ $`hC_A\left[+{\displaystyle \frac{13}{6}}{\displaystyle \frac{1}{2}}\xi _L\right]+hn_fT\left[{\displaystyle \frac{4}{3}}\right]`$ (D.2)
$`+h^2C_A^2\left[+{\displaystyle \frac{59}{8}}{\displaystyle \frac{11}{8}}\xi _L{\displaystyle \frac{1}{4}}\xi _L^2\right]+h^2C_An_fT\left[5\right]+h^2C_Fn_fT\left[4\right]`$
$`+h^3C_A^3\left[+{\displaystyle \frac{9965}{288}}{\displaystyle \frac{3}{4}}\zeta _3\xi _L{\displaystyle \frac{3}{16}}\zeta _3\xi _L^2{\displaystyle \frac{9}{16}}\zeta _3{\displaystyle \frac{167}{32}}\xi _L{\displaystyle \frac{33}{32}}\xi _L^2{\displaystyle \frac{7}{32}}\xi _L^3\right]`$
$`+h^3C_A^2n_fT\left[{\displaystyle \frac{911}{18}}+18\zeta _3+2\xi _L\right]+h^3C_AC_Fn_fT\left[{\displaystyle \frac{5}{18}}24\zeta _3\right]`$
$`+h^3C_An_f^2T^2\left[+{\displaystyle \frac{76}{9}}\right]+h^3C_Fn_f^2T^2\left[+{\displaystyle \frac{44}{9}}\right]+h^3C_F^2n_fT\left[+2\right].`$
### D.3 The Ghost Field Anomalous Dimension
$`\stackrel{~}{\gamma }_3`$ $`=`$ $`hC_A\left[+{\displaystyle \frac{3}{4}}{\displaystyle \frac{1}{4}}\xi _L\right]+h^2C_A^2\left[+{\displaystyle \frac{95}{48}}+{\displaystyle \frac{1}{16}}\xi _L\right]+h^2C_An_fT\left[{\displaystyle \frac{5}{6}}\right]`$ (D.3)
$`+h^3C_AC_Fn_fT\left[{\displaystyle \frac{45}{4}}+12\zeta _3\right]+h^3C_An_f^2T^2\left[{\displaystyle \frac{35}{27}}\right]+h^3C_A^2n_fT\left[{\displaystyle \frac{97}{108}}9\zeta _3+{\displaystyle \frac{7}{8}}\xi _L\right]`$
$`+h^3C_A^3\left[+{\displaystyle \frac{15817}{1728}}+{\displaystyle \frac{3}{8}}\zeta _3\xi _L+{\displaystyle \frac{3}{32}}\zeta _3\xi _L^2+{\displaystyle \frac{9}{32}}\zeta _3{\displaystyle \frac{17}{32}}\xi _L{\displaystyle \frac{3}{32}}\xi _L^2{\displaystyle \frac{3}{64}}\xi _L^3\right].`$
### D.4 The Quark Field Anomalous Dimension
$`\gamma _2`$ $`=`$ $`hC_F\left[\xi _L\right]+h^2C_AC_F\left[{\displaystyle \frac{25}{4}}2\xi _L{\displaystyle \frac{1}{4}}\xi _L^2\right]+h^2C_Fn_fT\left[+2\right]+h^2C_F^2\left[+{\displaystyle \frac{3}{2}}\right]`$ (D.4)
$`+h^3C_A^2C_F\left[{\displaystyle \frac{9155}{144}}{\displaystyle \frac{3}{4}}\zeta _3\xi _L{\displaystyle \frac{3}{8}}\zeta _3\xi _L^2+{\displaystyle \frac{69}{8}}\zeta _3{\displaystyle \frac{263}{32}}\xi _L{\displaystyle \frac{39}{32}}\xi _L^2{\displaystyle \frac{5}{16}}\xi _L^3\right]`$
$`+h^3C_AC_Fn_fT\left[+{\displaystyle \frac{287}{9}}+{\displaystyle \frac{17}{4}}\xi _L\right]+h^3C_AC_F^2\left[+{\displaystyle \frac{143}{4}}12\zeta _3\right]+h^3C_Fn_f^2T^2\left[{\displaystyle \frac{20}{9}}\right]`$
$`+h^3C_F^2n_fT\left[3\right]+h^3C_F^3\left[{\displaystyle \frac{3}{2}}\right].`$
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# Ordering-based Representations of Rational InferenceWork supported by Training through Research Contract No. ERBFMBICT950324 between the European Community and Università degli Studi di Roma “La Sapienza”.
## 1 Introduction
A recent breakthrough in nonmonotonic logic is the beginning of study of nonmonotonic consequence through postulates for abstract nonmonotonic consequence relations, using Gentzen-like context-sensitive sequents (, , ). The outcome of this research turns out to be valuable in at least two ways
* it provides a sufficiently general axiomatic framework for comparing and classifying nonmonotonic formalisms, and
* it gave rise to new, simpler, and better behaved systems for nonmonotonic reasoning, such as cumulative (), preferential (), and rational () inference relations.
It is unfortunate that these new inference relations enjoy only one, semantical, representation; that of preferential models (). We have that preferential, preferential transitive, and preferential modular or ranked models characterize cumulative, preferential and rational inference relations, respectively (, ). An additional second-order constraint must be imposed on these models, called stoppering or smoothness. However, this modeling is insufficient because in order to employ the above inference relations, one must be able to generate them. This is crucial when we want to design a system that reasons using the above inference relations. In such a case, one comes up with a set of rules or defaults that one wants to apply, imposes a prioritization on them, and provides a mechanism which ensures that answers are derived according to these inference relations. This is exactly the proof-theoretic approach expressed by default logic. However, no similar proof-theoretic notion is provided in the above framework.
In this paper, we offer two new, alternative representations for rational inference. The first representation is algebraic and obtained through a simple class of orderings of formulas, called rational orderings. The second representation is proof-theoretic and obtained through a class of consequence operators based on the way we handle defaults, called ranked consequence operators. Moreover, a correspondence result between these classes is established.
The first link between nonmonotonic inference relations and a class of orderings of formulas was given by Gärdenfors and Makinson in . However, the nonmonotonic system defined by an ordering of formulas is not one of the previously mentioned systems, but a translation of the well-known belief revision AGM axioms () into nonmonotonic reasoning, called expectation inference relations. Expectation inference relations are rational inference relations together with a rule called Consistency Preservation. Moreover, Gärdenfors and Makinson’s representation of expectation inference relations with orderings of formulas is not appropriate, in the sense that the correspodence is not bijective. (Two orderings of formulas can generate the same inference relation.) So, two questions remain open. Namely,
* is there a way to generate one of the independently motivated nonmonotonic inference relations (cumulative, preferential, rational) with a class of orderings of formulas?, and
* can the correspondence be bijective?
We answer affirmatively both questions for rational inference. Our approach is the following. We study the rule of Consistency Preservation and, by giving it a precise syntactic characterization, show that its role is insignificant in the context of preferential reasoning. Drawing from this intuition, we introduce new defining conditions relating the classes of orderings of formulas and nonmonotonic consequence relations and show that Gärdenfors-Makinson orderings are in bijective correspondence with rational inference. Moreover, we introduce a smaller class of orderings, which, under our translation, is in bijective correspondence with the Gärdenfors-Makinson expectation inference relations. This is how the first representation result for rational inference is obtained. This result adds to a long tradition of defining nonmonotonic logics with orderings of formulas (, , , ).
The above representation result is more “constructive” than the semantical completeness of preferential models. However, rational orderings must have a concise, constructive representation. To this end, we encode a natural way of applying defaults into a new class of consequence relations, called ranked consequence operators. Each member of this class generates a rational ordering, and conversely, hence the class of ranked consequence relations coincides with that of rational inference. Also, we show how previous default logic systems in the literature (, ) reduce to our framework.
The above results pave the way towards a study of nonmonotonicity through orderings of formulas, allow us to translate previous work in belief revision into the context of nonmonotonic reasoning, and provide a framework for designing default systems obeying rational inference.
The plan of this paper as follows. In Section 2, we briefly introduce the relations under study, explain the rule of Consistency Preservation, and provide a characterization for this rule. In Section 3, we introduce the orderings, their translations and our first representation theorem. In Section 4, we introduce ranked consequence operators and our second representation theorem. In Section 5, we show how one can generate a ranked consequence operator given a prioritized family of sets of defaults and, in Section 6, conclude. A preliminary version of the first half of this paper appeared in . Results from the second half were announced in .
## 2 Shifting underlying entailment
Before going to the main result of this section, we shall make a brief introduction to the nonmonotonic consequence relations under study. Assume a language $``$ of propositional constants closed under the boolean connectives $``$ (disjunction), $``$ (conjunction), $`\neg `$ (negation) and $``$ (implication). We shall use greek letters $`\alpha `$, $`\beta `$, $`\gamma `$, etc. for propositional variables. We shall also use $`\alpha \text{ }\beta `$ , read as “$`\alpha `$ normally entails $`\beta `$”, to denote the nonmonotonic consequence relation ($`\text{ }\times `$). Before we present the first set of rules for $`\text{ }`$, we need a symbol for a classical-like entailment. We shall use $``$. The relation $``$ need not be that of classical propositional logic. We require that $``$ includes classical propositional logic, satisfies compactness (i.e., if $`X\beta `$ then there exists a finite subset $`Y`$ of $`X`$ such that $`Y\beta `$)<sup>1</sup><sup>1</sup>1We write $`X,\alpha \beta `$ for $`X\{\alpha \}\beta `$., the deduction theorem (i.e., $`X,\alpha \beta `$ if and only if $`X\alpha \beta `$) and disjunction in premises (i.e., if $`X,\alpha \beta `$ and $`X,\gamma \beta `$ then $`X,\alpha \gamma \beta `$). The reader will notice that these are the only properties we make use of in the subsequent proofs. We shall denote the consequences of $`\alpha `$ with $`Cn(\alpha )`$ and $`C(\alpha )`$ under $``$ and $`\text{ }`$, respectively.
The rules mentioned in the following are presented in Table 1. For a motivation of these rules, see and . (The latter serves as an excellent introduction to nonmonotonic consequence relations.)
Definition 1 Following (, , ), we shall say that a relation $`\text{ }`$ on $``$ is an inference relation (based on $``$) if it satisfies Supraclassicality, Left Logical Equivalence, Right Weakening, and And. We shall call an inference relation $`\text{ }`$ preferential if it satisfies, in addition, Cut, Cautious Monotonicity, and Or. We shall call an inference relation $`\text{ }`$ rational if it is preferential and satisfies, in addition, Rational Monotonicity. Finally, we shall say that $`\text{ }`$ is an expectation inference relation (based on $``$) if it is a rational and satisfies, in addition, Consistency Preservation.
The most controversial of these rules is Rational Monotonicity, which, moreover, is non-Horn. For a plausible counterexample, see .
Expectation inference relations correspond to the so-called AGM postulates for belief revision (), as it was shown in , and only differ from rational relations in that they satisfy the following rule, called Consistency Preservation:
$$\frac{\alpha \text{ }}{\alpha }$$
where $``$ is classical entailment. Consistency Preservation says that a logically not false belief cannot render our set of beliefs inconsistent. This makes a difference between the two classes, in the following sense. Using rational inference, I can rely on an inference such as $`\beta \text{ }`$, where $`\beta `$ is the statement “I am the Queen of England”. On the other hand, expectation inference would not allow that, since, even if I am certain I am not the Queen of England, one could think a world where I could have been. This becomes more important if, instead of a belief set, one considers a conditional base. For example, consider a database for air-traffic. The statement “two airplanes are scheduled to arrive at the same time and land on the same place” should infer inconsistency on this database, although it is not a falsity. More examples can be drawn from physical laws. This means that rational inference is the logic of “hard constraints”, that is of statements (not necessarily tautologies) I cannot revise without deconstructing the whole inference mechanism. This is not admitted in expectation inference: all statements are allowed to be revised apart from tautologies (or whatever is a consequence of the empty set under the underlying entailment).
In and , it was observed that preferential entailment satisfies a weaker form of consistency preservation: there exists a consequence operation $`^{}`$ with $`^{}\text{ }`$ such that $`\text{ }`$ satisfies Consistency Preservation with respect to $`^{}`$. This was proved by semantical arguments.
In the following theorem, we make this property more precise by expressing it in syntactic terms. We show that the required underlying consequence operation retains the properties of the initial one, as it only differs on the set of assumptions. Therefore, the relation between an expectation and a rational inference relation is that of a logic with its theory.
For the proof of Theorem 3, we shall make use of the following rules (derived in any preferential inference relation).
###### Lemma 2
In any preferential inference relation, the following rules hold
1. $`{\displaystyle \frac{\alpha \text{ }\beta \beta \text{ }\alpha \alpha \text{ }\gamma }{\beta \text{ }\gamma }}`$ (Reciprocity)
2. $`{\displaystyle \frac{\alpha \beta \text{ }\gamma }{\alpha \text{ }\beta \gamma }}`$ (S)
3. $`{\displaystyle \frac{\alpha \text{ }}{\alpha \beta \text{ }}}`$
4. $`{\displaystyle \frac{\alpha \beta \text{ }}{\alpha \text{ }\neg \beta }}`$
5. $`{\displaystyle \frac{\alpha \beta \text{ }}{\alpha \text{ }}}`$
Proof. Rules 12 and 3 were introduced and shown to be derived in a preferential relation in . For 4, suppose $`\alpha \beta \text{ }`$. Applying S, we get $`\alpha \text{ }\beta `$ and, by Right weakening, we conclude 4. For 5, suppose $`\alpha \beta \text{ }`$. Then, by 3, we get $`(\alpha \beta )\alpha \text{ }`$ and, by Left Logical Equivalence, we conclude 5.
###### Theorem 3
Let $`\text{ }`$ be a preferential inference relation based on $``$. Then $`\text{ }`$ is a preferential inference relation based on $`^{}`$ that satisfies the Consistency Preservation rule, where
$$\alpha ^{}\beta \text{iff}\mathrm{\Gamma },\alpha \beta ,$$
and
$$\mathrm{\Gamma }=\{\neg \gamma :\gamma \text{ }\}.$$
Proof. We must prove that $`\text{ }`$ satisfies Supraclassicality, Left Logical Equivalence, Right Weakening and Consistency Preservation with respect to $`^{}`$. The rest of the rules are already satisfied since they do not involve an underlying consequence relation.
First notice that Consistency Preservation is immediate by definition of $`^{}`$.
For Supraclassicality, suppose $`\alpha ^{}\gamma `$ then $`\mathrm{\Gamma },\alpha \gamma `$. By compactness of $``$, there exist $`\beta _1,\mathrm{},\beta _n`$ such that $`\beta _1\text{ },\mathrm{},\beta _n\text{ }`$ and $`\neg \beta _1,\mathrm{},\neg \beta _n,\alpha \text{ }\gamma `$. By repeated applications of Or, we get $`\beta _1\mathrm{}\beta _n\text{ }`$. Let $`\beta =\beta _1\mathrm{}\beta _n`$, then $`\beta \text{ }`$ and $`\alpha \neg \beta \gamma `$. By Supraclassicality of $`\text{ }`$ on $``$, we have $`\alpha \neg \beta \text{ }\gamma `$. By Lemma 2.3, we have $`\alpha \beta \text{ }`$, so, by Lemma 2.4, we have $`\alpha \text{ }\neg \beta `$. Using Cut, we get $`\alpha \text{ }\gamma `$, as desired.
For Left Logical Equivalence, suppose $`\alpha \text{ }\gamma `$, $`\alpha ^{}\beta `$, and $`\beta ^{}\alpha `$, i.e. $`\mathrm{\Gamma },\alpha \beta `$ and $`\mathrm{\Gamma },\beta \alpha `$. By compactness, there exist $`\delta _1,\delta _2`$ such that $`\alpha \text{ }\neg \delta _1`$, $`\alpha \text{ }\neg \delta _2`$, $`\alpha \delta _1\beta `$, and $`\beta \delta _2\alpha `$. As above, we have $`\alpha \text{ }\beta `$ and $`\beta \text{ }\alpha `$. Therefore, by Lemma 2.1, we get $`\beta \text{ }\gamma `$, as desired.
Coming to Right Weakening, suppose $`\alpha \text{ }\beta `$ and $`\beta ^{}\gamma `$, i.e. there exists $`\delta `$ such that $`\alpha \text{ }\neg \delta `$ and $`\beta \neg \delta \text{ }\gamma `$. By And, we have $`\alpha \text{ }\beta \neg \delta `$, so, using Right Weakening of $`\text{ }`$ on $``$, we get $`\alpha \text{ }\gamma `$, as desired.
Notice that the result applies to rational inference relations, as well, since the latter are preferential, by definition. We interpret the above result as follows. Once we strengthen the underlying entailment, rational inference will become an expectation inference and, therefore, can be treated as such. It also implies that the logic of hard and soft constraints is basically the same, their only difference being what we consider a consequence of the underlying propositional entailment. Hard constraints are just taking a place in our belief set as “guarded” as that of, say, tautologies. Whatever remains is subject to revision, and hence a soft constraint.
## 3 Rational inference and orderings
Now that we established the correspondence between rational and expectation inference relations, we shall extend it to a particularly attractive characterization of the latter with orderings of formulas. We shall first review Gärdenfors-Makinson’s results and then present our own.
The intuition behind ordering-based formalisms is common in works on belief revision, possibilistic logic, and decision theory. We order sentences according to our expectations. A relation “$`\alpha <\beta `$” is interpreted as “$`\beta `$ is expected more than $`\alpha `$”, or “$`\alpha `$ is more surprising than $`\beta `$”, or “$`\beta `$ is more possible than $`\alpha `$”. One can treat such an ordering as a primary notion; this is the approach of this paper. However, in case of rational orderings, one can show that such an ordering induces a function from the extensions of formulas to the unit interval. This function induces a possibility measure on the extensions of formulas (see ). A possibility measure is a “weak” probability measure on these extensions. Roughly, it replaces addition with maximisation. Although the connections with probability are not clear yet (see , ), probability measures seem especially suited for modeling cases under uncertainty. Further, a possibility measure arises naturally out of a database. Zadeh’s theory for approximate reasoning () provides a method for turning available information of a certain form (“fuzzy” database) into a possibility measure and, therefore, gives rise to a rational ordering of sentences.
We find that, by a logical point of view, orderings correspond to prioritization. We prefer a proof-theoretic reading, made more explicit in Section 4, “$`\alpha `$ is more defeasible than $`\beta `$” or “$`\alpha `$ has lower priority than $`\beta `$”. A notion of proof is developed in Section 4 based on this prioritization and justifies the use of rational orderings without appealing to some probabilistic intuition.
Definition 4 A rational ordering is a relation $``$ on $``$ which satisfies the following properties:
| 1. If $`\alpha \beta `$ and $`\beta \gamma `$, then $`\alpha \gamma `$ | (Transitivity), |
| --- | --- |
| 2. If $`\alpha \beta `$, then $`\alpha \beta `$ | (Dominance), |
| 3. $`\alpha \alpha \beta `$ or $`\beta \alpha \beta `$ | (Conjunctiveness). |
The original name of these orderings was expectation orderings. However, we shall see that this name is not justifiable, since expectation inference relations correspond to a smaller class of orderings (see Definition 10).
One can easily derive from the above properties that a rational ordering satisfies
1. connectivity, i.e. $`\alpha \beta `$ or $`\beta \alpha `$, and
2. either $`\alpha \beta `$, for all $`\beta `$, or $`\neg \alpha \beta `$, for all $`\beta `$.
We should mention that the above properties of rational relations are not new. It is not easy to assign credits, but they have appeared in works in belief revision (), possibilistic logic (), ), fuzzy logic (), theory of evidence (), and economics () (see for a historical reference).
Gärdenfors and Makinson define the following maps between the class of expectation inference relations and rational orderings.
Definition 5 Given a rational ordering $``$ and an expectation inference relation $`\text{ }`$, then define a consequence relation $`\text{ }^{}`$ and an ordering $`^{}`$ as follows
($`C`$) $`\alpha \text{ }^{}\gamma `$ iff either $`\alpha \gamma `$ or there is a $`\beta `$ such that $`\alpha \beta \gamma `$ and $`\neg \alpha <\beta `$. ($`O`$) $`\alpha ^{}\beta `$ iff either $`\alpha \beta `$ or $`\neg (\alpha \beta )\text{ }\sim ̸\alpha `$.
We shall also denote $`\text{ }^{}`$ and $`^{}`$ with $`C()`$ and $`O(\text{ })`$, respectively.
Condition ($`O`$) is critical and due to Rott (). Now, one can prove the following.
###### Theorem 6
Given a rational ordering $``$ and an expectation inference relation $`\text{ }`$, then $`C()`$ is an expectation inference relation and $`O(\text{ })`$ is a rational ordering. Moreover, we have $`\text{ }=C(O(\text{ }))`$.
This theorem, although it exhibits the first connection between some class of nonmonotonic consequence relations and orderings of formulas, has two disadvantages. First, the way it achieves consistency preservation is ad hoc. If that was not the case, then the condition ($`C`$) would be inappropriate, since, in the first part, it refers explicitly to the underlying entailment<sup>2</sup><sup>2</sup>2However, the second part should remain the same since we do not mind having a few more consequences, as long as, the rules which govern the underlying entailment do not change.. Second, it fails to show an isomorphism between the class of expectation inference relations and that of rational orderings, that is $`=O(C())`$. If the second was not the case, then the condition ($`O`$) would use only the expectation inference relation to construct the ordering. Consider the following example.
Example 7 Let $`𝒟_1=\{\{\},\{,\alpha \}\}`$ and $`𝒟_2=\{\{\alpha \}\}`$. Now define orderings on $``$ as follows
$$\beta _1\gamma \text{iff}A\beta \text{implies}A\gamma ,\text{for all}A𝒟_1.$$
Similarly, for $`𝒟_2`$ and $`_2`$. We have that $`_1_2`$, and by Proposition 19, the orderings are rational. However, they generate the same expectation inference relation, using ($`O`$).
Drawing from the above intuitions and Theorem 3, we define
Definition 8 Given a rational ordering $``$ and a rational inference relation $`\text{ }`$, then define a consequence relation $`\text{ }^{}`$ and an ordering $`^{}`$ as follows
| ($`𝐂`$) | $`\alpha \text{ }^{}\gamma `$ | iff | either $`\beta \neg \alpha `$, for all $`\beta `$, |
| --- | --- | --- | --- |
| | | | or there is a $`\beta `$ such that $`\alpha \beta \gamma `$ and $`\neg \alpha <\beta `$. |
| ($`𝐎`$) | $`\alpha ^{}\beta `$ | iff | either $`\neg (\alpha \beta )\text{ }`$ or $`\neg (\alpha \beta )\text{ }\sim ̸\alpha `$. |
We shall also denote $`\text{ }^{}`$ and $`^{}`$ with $`𝐂()`$ and $`𝐎(\text{ })`$, respectively.
For Theorem 10, we need the following lemma.
###### Lemma 9
Let $``$ and $`\text{ }`$ be a rational ordering and inference relation, respectively. Then
1. If $`\alpha \text{ }`$ then $`\beta ^{}\neg \alpha `$, for all $`\beta `$, where $`^{}=𝐎(\text{ })`$.
2. If $`\beta \neg \alpha `$, for all $`\beta `$, then $`\alpha \text{ }^{}`$, where $`\text{ }^{}=𝐂()`$.
3. $`\neg \alpha <\alpha \gamma `$ iff $`\{\beta :\neg \alpha <\beta \}\alpha \gamma `$.
Now, everything falls into place.
###### Theorem 10
Given a rational ordering $``$ and a rational inference relation $`\text{ }`$, then $`𝐂()`$ is a rational inference relation and $`𝐎(\text{ })`$ is a rational ordering. Moreover, we have $`\text{ }=𝐂(𝐎(\text{ }))`$ and $`=𝐎(𝐂())`$.
Now, if rational orderings are in adjunction with rational inference relations, what is the class of orderings which corresponds to expectation inference relations? For that, observe that by Lemma 9, hard constraints are positioned on the top of rational orderings. So, it is enough to keep exclusively this place for the consequences of the empty set and add this as a condition to rational orderings.
Definition 11 An expectation ordering is a rational ordering which satisfies, in addition, the following property:
> If $`\beta \alpha `$, for all $`\beta `$, then $`\alpha `$.
Now, using the same defining conditions ($`𝐂`$) and ($`𝐎`$), we can state the improved characterization theorem for expectation inference relations.
###### Theorem 12
Given an expectation ordering $``$ and an expectation inference relation $`\text{ }`$, then $`𝐂()`$ is an expectation inference relation and $`𝐎(\text{ })`$ is an expectation ordering. Moreover, we have $`\text{ }=𝐂(𝐎(\text{ }))`$ and $`=𝐎(𝐂())`$.
Theorems 10 and 12 can now be used for giving new straightforward proofs for the characterization of rational inference with ranked preferential models (Theorem 3.12 of ) and expectation inference with nice preferential models (Theorems 3.8 and 3.9 of ). These proofs will appear elsewhere.
## 4 Ranked consequence operators
First, a word about the plan of this section. We introduce the notion of ranked consequence operation without referring to an underlying entailment (Definition 4). The reason for such a definition is that we can motivate ranked consequence operators independently of nonmonotonic reasoning. Then we define a smaller class based on an underlying entailment (Definition 4) and show that this class characterize rational inference relations. The same constraints we assumed for a language $``$ and an entailment $``$ in Section 2 continue to hold here.
Think of a reasoner whose beliefs are ordered accordingly to their defeasibility. Beliefs which are less likely to be defeated come before beliefs which are more likely to be defeated. For instance “Birds fly” will come after “Penguins do not fly” (since the former has more exceptions) and “Mary is married” might come before “Mary is married with children” (since the latter is stronger). There is a natural way to attach a consequence operator to this belief prioritization.
Definition 13 Let $`I,<`$ be a linear ordering, and $`\{A_i\}_{iI}`$ be an upward chain of sets of formulas such that $`A_iA_j`$ iff $`ij`$. Define the following consequence operators (one for each $`iI`$):
| $`\text{ }_i\beta `$ | iff | $`\beta A_i`$, |
| --- | --- | --- |
| $`\alpha \text{ }_i\beta `$ | iff | $`\neg \alpha A_i`$ and $`\alpha \beta A_i`$. |
Now let
| $`\alpha \text{ }\beta `$ | iff | either $`\alpha \text{ }_i\beta `$, for some $`iI`$, |
| --- | --- | --- |
| | | or $`\text{ }_i\neg \alpha `$, for all $`iI`$. |
The consequence operator $`\text{ }`$ will be called ranked consequence operator (induced by $`\{A_i\}_{iI}`$).
First, note that $`A_i`$’s are not necessarily deductively closed. Second, notice that, unless we add the last part of the definition of $`\text{ }`$, we do not provide for formulas $`\alpha `$, where $`\neg \alpha A_i`$, for all $`iI`$. In order to have $`\alpha \text{ }`$, there must either be an $`A_i`$ such that $`\neg \alpha A_i`$ and $`\alpha A_i`$, or $`\neg \alpha A_i`$, for all $`iI`$. This means that if our beliefs can accommodate a context where $`\alpha `$ holds, then we use the part of the ordering that remains consistent after adding $`\alpha `$. Therefore the $`A_i`$’s which contain both $`\neg \alpha `$ and $`\alpha `$ are irrelevant to the consequence operator.
Indices assign grades of relying on the set of consequences as the next example, formalizing omniscience, shows.
Example 14 Let $`\omega ,<`$ be the set of natural numbers with the usual order. Now let $``$ be the classical consequence relation and let $`A_1`$ be some set of formulas of propositional logic. Let
$$A_2=\{\varphi \text{is provable in one step from}A_1\},$$
and, inductively,
$$A_n=\{\varphi \text{is provable in less than }n1\text{ steps from}A_1\}.$$
Notice that if $`A_1`$ is consistent and $`\text{ }`$ is the ranked consequence operator defined through $`\{A_i\}_{i\omega }`$ then
$$\alpha \text{ }\beta \text{iff}A_1,\alpha \beta ,$$
where $``$ is the classical consequence operator of propositional calculus. Note that if $`A_1`$ is inconsistent, then $`\alpha `$ entails all formulas which are provable from $`A_1`$ with less steps than $`\neg \alpha `$, i.e., before we realize inconsistency. Now, if $`A_1`$ is the set of all tautologies or, better, an axiomatization of them then $`\text{ }`$ is exactly the classical consequence operator.
It is clear that the above representation is syntax-based, i.e. depends on the particular representation of $`A_i`$’s. The case where the sets of formulas $`A_i`$ are closed under consequence is the one we shall deal with in this paper. Doing that is like being logically omniscient; we do not assign any cost to derivations using $``$.
We shall now give a definition of ranked consequence operator using an underlying entailment. In case the $`A_i`$’s are closed under consequence, it coincides with the original definition (by replacing the set belonging relation $``$ with proposition entailment $``$).
Definition 15 A ranked consequence operator $`\text{ }`$ based on $``$ induced by a chain of sets $`\{B_i\}_{iI}`$ under inclusion is defined as follows:
We first define a set of consequence operators $`\text{ }_i`$ (one for each $`iI`$):
| $`\alpha \text{ }_i\beta `$ | iff | $`B_i⊬\neg \alpha `$ and $`B_i,\alpha \beta `$. |
| --- | --- | --- |
Note that we denote $`\text{ }_i\beta `$ with $`\text{ }_i\beta `$. We can now let
| $`\alpha \text{ }\beta `$ | iff | either $`\alpha \text{ }_i\beta `$, for some $`iI`$, |
| --- | --- | --- |
| | | or $`\text{ }_i\neg \alpha `$, for all $`iI`$. |
We shall use $`\{B_i\}_{iI},`$ to denote this operator.
Notice that we can have both $`⊬\alpha `$ and $`\text{ }_i\neg \alpha `$, for all $`i`$. This translates to the fact that $`\alpha `$ can be true in some possible world but it is unthinkable for us to include it in our beliefs. The above mechanism treats such a case as an instance of a hard constraint: such an $`\alpha `$ implies falsehood.
Again notice that unless $`B_i\neg \alpha `$, for all $`iI`$, we cannot derive falsity from $`\alpha `$. The reason is that, in those cases, we are able to form a context based on $`\alpha `$ (a chain of sets of formulas which prove $`\alpha `$) which is consistent. Again, the inconsistent $`B_i`$’s are irrelevant to the consequence operator. The following proposition allows us to assume that the ordering is complete, that is it has all meets and joins, and amounts essentially to Lewis’ assumption or smoothness property of preferential models.
###### Proposition 16
A ranked consequence operator $`\text{ }`$ based on $``$ is induced by a chain of sets of formulas if and only if it is induced by the closure of this chain under arbitrary unions and intersections.
This result has the following significance: it allows an assignment of a rank to an assertion of the form $`\alpha \text{ }\beta `$. Suppose that $`\alpha \text{ }\beta `$ holds. If $`\text{ }_i\neg \alpha `$ for all $`iI`$ does not hold, then the set $`I_{\alpha \text{ }\beta }=\{i:\text{ }\sim ̸_i\neg \alpha \text{ and }\alpha \text{ }_i\beta \}`$ is not empty. Moreover, it is connected. Now, it is easy to see that, in the completion of the chain, this set has a least element (because it is closed under intersection) and a greatest element (because it is closed under unions). Let $`i_1`$ and $`i_2`$ be the indices of the least and greatest elements, respectively. The rank of the assertion $`\alpha \text{ }\beta `$ is $`i_1`$ and its range $`[i_1,i_2]`$. In case $`\text{ }_i\neg \alpha `$, for all $`iI`$, then set the rank of $`\alpha \text{ }\beta `$ to $`0`$ and its range to $`[0,l]`$, where $`0`$ and $`l`$ are the indices of least and the greatest element of the linear order, respectively. In case of an assertion $`\text{ }\beta `$, observe that its range is of the form $`[i,l]`$, where $`l`$ is the index of the greatest element of the linear ordering.
Finally, notice that a ranked consequence operator is not necessarily monotonic.
Example 17 Let $`B_1=\{\alpha \}`$ and $`B_2=\{\alpha ,\neg \beta \}`$. We have $`\alpha \text{ }\neg \beta `$ because $`\neg \alpha Cn(B_2)`$ and $`\alpha \neg \beta Cn(B_2)`$. But we also have that $`\alpha \beta \text{ }\sim ̸\neg \beta `$ because $`\neg (\alpha \beta )Cn(B_1)`$.
Now, it is interesting to ask what kind of properties a ranked consequence operator satisfies. It turns out that each ranked consequence operator gives rise to a rational inference relation. Although one can show it directly, we define the rational orderings induced by such operators.
Definition 18 Given a ranked consequence operator, let
$$\alpha \beta \text{iff}B_i\alpha \text{implies}B_i\beta ,\text{for all }iI.$$
and call $``$ the ordering induced by the ranked consequence operator $`\text{ }`$.
We, now, have the following
###### Proposition 19
An ordering induced by a ranked consequence operator is rational. Moreover, $`𝐂()=\text{ }`$.
We have immediately the following.
###### Corollary 20
A ranked consequence operator is a rational inference relation.
The other direction of the above theorem holds, too. We should only show, given a rational ordering, how to generate a total order of sets of formulas. To this end, we shall define a chain of sets $`\{A_i\}_{iI}`$ which generates a ranked consequence operator $`\text{ }`$ equal to $`𝐂()`$. Let $``$ be the equivalence relation induced by $``$ (a rational ordering is a preorder). The equivalence classes will be denoted by $`\widehat{\alpha }`$ (where $`\alpha \widehat{\alpha }`$). It is also clear that the set of equivalence classes is linearly ordered. Now, for each $`\alpha `$, let
$$A_{\widehat{\alpha }}=\{\beta :\alpha \beta \}.$$
Note here that, by Dominance, the sets $`A_{\widehat{\alpha }}`$ are closed under consequence. Moreover, we have $`A_{\widehat{\alpha }}A_{\widehat{\beta }}`$ iff $`\beta \alpha `$. Now, generate a ranked consequence operator $`\text{ }`$ as in Definition 16. This ranked consequence operator turns out to be equal to the one generated by the rational order. So, we have the following.
###### Theorem 21
A rational inference relation is a ranked consequence operator.
The proof of the above theorem shows that a rational ordering can be defined by a chain of sets which induces a ranked consequence operator and conversely. However, the same rational ordering can be induced by two different ranked consequence operators. This should hardly be surprising, as ranked consequence operators play the role of axiomatizing a nonmonotonic “theory”, that is, a rational inference. Moreover,
* ranked consequence operators are proof-theoretic in their motivation, and therefore closer to what we want to describe by a rational inference relation, and
* a ranked consequence operator assigns ranks to assertions as well as to formulas therefore grading the whole process of inference.
We showed that rational and expectation inference relations are exactly the same class of consequence relations if we allow the underlying propositional entailment to “vary”. However fixing $``$, is it possible to tell if a ranked consequence operator satisfies Consistency Preservation? The answer is affirmative, for a formula $`\alpha `$ infers inconsistency ($`\alpha \text{ }`$) if and only if its negation is a consequence of the first element of the chain which induces the ranked consequence operator (as a corollary of Proposition 16, a first element always exists). To see that, suppose $`\alpha \text{ }`$, then, by definition, we must have $`\text{ }_i\neg \alpha `$, for all $`iI`$, and for that it is enough that the first element of the chain implies $`\neg \alpha `$.
## 5 Rational default systems
In this section, we shall see how one can design a ranked consequence operator. Suppose we are given a number of sets of (normal, without prerequisites) defaults in a linear well-founded prioritization. Moreover, and this is an important assumption for rational inference, we are asked to, either apply the whole set, or not apply it at all. Let the set of sets of defaults be $`𝒟=\{A_i\}_{iI}`$, where $`I,<`$ is a well-founded linear strict order, and $`A_i`$ is preferred from $`A_j`$ whenever $`i<j`$. The are two ways to read this preference.
* The first way is a strict one: if you cannot add $`A_i`$ to your set of theorems (that is, you derive inconsistency by adding $`A_i`$) then you cannot add $`A_j`$, for all $`A_j`$ less preferred from $`A_i`$.
* the other is liberal: if you cannot add $`A_i`$ to your set of theorems, then you can add $`A_j`$, where $`A_j`$ is less preferred from $`A_i`$, provided you cannot add $`A_k`$, where $`A_k`$ is more preferred than $`A_j`$.
To illustrate this, consider the following example.
Example 22 Let $`𝒟=\{A_1,A_2,A_3\}`$, where $`A_1=\{\alpha \beta \}`$, $`A_2=\{\neg \beta \}`$, and $`A_3=\{\beta \gamma \}`$. Assume $`\alpha `$. Following the strict interpretation, we can only infer $`\beta `$, from $`A_1`$. With the liberal interpretation, we can also infer $`\gamma `$, since we are allowed to add $`A_3`$, and cannot add $`A_2`$ that leads to a contradiction.
It turns out that those readings are equivalent. Not in the sense that the same set of sets of defaults generate the same consequences, but that a strict extension of a family of sets of defaults can be reduced to a liberal extension of another family of sets of defaults, and conversely. It can be easily shown that strict and liberal extensions of families of sets of defaults are instances of rational consequence operators and, therefore, rational. In particular, Proposition 19 gives us a way to construct the rational orderings of such default systems.
Given $`𝒟=\{A_i\}_{iI}`$, where $`I,<`$ is well-founded, define the following strict ordering between non-empty subsets of $`I`$:
| $`K<L`$ | iff | there exists $`iI`$ such that | a. $`iK`$ but $`iL`$, and |
| --- | --- | --- | --- |
| | | | b. for all $`j<i`$, $`jK`$ iff $`jL`$. |
It can be shown that $`<`$ is linear. Now, let $`A_K=_{KL}Cn(_{iL}A_i)`$.
Definition 23 Let $`\alpha `$ and $`𝒟=\{A_i\}_{iI}`$, where $`I,<`$ is a total strict order, and $`A_i`$, for all $`iI`$. The strict extension $`E_𝒟^s(\alpha )`$ of $`\alpha `$ with respect to $`𝒟`$ is defined as follows
$$E_𝒟^s(\alpha )=\{E_i^s:E_i^s\text{ is consistent}\},$$
where $`E_i^s=Cn(\{\alpha \}_{ji}A_j)`$. The liberal extension $`E_𝒟^l(\alpha )`$ of $`\alpha `$ with respect to $`𝒟`$ is defined as follows
$$E_𝒟^l(\alpha )=\{E_L^l:E_L^l\text{ is consistent}\},$$
where $`E_L^l=Cn(\{\alpha \}_{KL}A_K)`$.
Thus the liberal extension of $`𝒟=\{A_i\}_{iI}`$ is the strict extension of $`𝒟^{}=\{A_K\}_{K𝒫(I)^{}}`$. For the other direction, the strict extension of $`𝒟=\{A_i\}_{iI}`$ coincides with the liberal extension of $`𝒟^{\prime \prime }=\{C_i\}_{iI}`$, where $`C_i=_{ji}A_i`$.
So, it is enough to construct the ranked consequence operator for the strict extension of $`𝒟=\{A_i\}_{iI}`$. But this is easily achieved. Consider $`\{B_i\}_{iI},\text{ }`$, where $`B_i=E_i^s`$.
Thus, strict and liberal extensions of prioritized sets of set of formulas are rational. The above definition, together with Proposition 19, gives us a way to construct the rational orderings of such default systems. Given a prioritized set $`𝒟=\{A_i\}_{iI}`$ then the rational ordering of its strict extension is
$$\alpha _𝒟^s\beta \text{iff}\underset{ji}{}A_j\alpha \text{implies}\underset{ji}{}A_j\beta ,\text{for all }iI.$$
The rational ordering of its liberal extension is
$$\alpha _𝒟^l\beta \text{iff}\underset{KL}{}A_K\alpha \text{implies}\underset{KL}{}A_K\beta ,\text{for all}L𝒫(I)^{}.$$
Assuming a finite language, sets of formulas, intersections, and unions of them correspond to conjunctions, disjunctions, and conjunctions, respectively.
A study of the above default systems under the assumption of finite language, has been carried already in the context of belief revision (therefore, assuming consistency preservation, in addition to finite language), by Nebel (, ). Our strict and liberal extensions are called prioritized and linear base revision, respectively. Also, Nebel showed in that deciding if a certain formula is contained in the strict or in a liberal extension (that is, deciding $`\alpha \text{ }\beta `$) is $`\mathrm{P}^{\mathrm{NP}[O(\mathrm{log}^cn)]}`$ and $`\mathrm{P}^{\mathrm{NP}[O(n)]}`$, respectively. We expect that these results carry over to our framework.
## 6 Conclusion
We summarize our results in the following
###### Theorem 24
Let $`\text{ }`$ be a binary relation on $``$. Then the following are equivalent:
1. $`\text{ }`$ is a rational inference relation, i.e. it satisfies Supraclassicality, Left Logical Equivalence, Right Weakening, And, Cut, Cautious Monotonicity, Or, and Rational Monotonicity.
2. $`\text{ }`$ is characterized by some rational relation $``$ on $``$ using condition ($`𝐂`$).
3. $`\text{ }`$ is defined by a ranked consequence operator.
Since rational orderings are in one-to-one correspondence with rational inference, our first representation result has many ramifications. Results in belief revision can be translated in a nonmonotonic framework and vice versa. For instance, selection functions and preferential models can be used for the modeling of both. Proofs of this results are straightforward through our defining conditions for rational and expectation orderings. Work that has already been done on expectation inference relations (e.g. the study of generating expectation inference relations through incomplete rational orderings—see ) can be lifted smoothly to rational inference.
Our second representation result reveals the working mechanism of rational inference. It shows that, in order to attain rational inference, we must prioritize defaults in a particular way. We showed how default logic formalisms can fit this pattern. It enables us to assign grades to all components of the reasoning system (formulas and rules). Therefore, it is a particular attractive way to use it as an inference mechanism for nonmonotonic reasoning.
The above characterization results reveal another notion of consequence paradigm hidden behind nonmonotonicity. However, apart from Dubois and Prade’s work on possibility logic, this paradigm has been passed largely unrecognized by logicians as an appropriate method for a treatment of vagueness and uncertainty. Yet, this paradigm arose independently from various studies on different fields and appeared before nonmonotonic logic. In addition, it is applicable. Now, an important question arises: how far this paradigm extends. In other words, is it possible to reduce a nonmonotonic consequence relation to some relation expressing prioritization? The answer is positive and uniform. The important case of preferential consequence relations is treated separately in , while the general case (which includes cumulative consequence relations) appears in .
Acknowledgements: I would like to thank Gianni Amati and R. Ramanujam for their helpful comments on a preliminary version of this paper.
## Appendix A Proofs
For the proof of Theorem 3, we shall make use of the following rules (derived in any preferential inference relation).
Lemma 2 In any preferential inference relation, the following rules hold
1. $`{\displaystyle \frac{\alpha \text{ }\beta \beta \text{ }\alpha \alpha \text{ }\gamma }{\beta \text{ }\gamma }}`$ (Reciprocity)
2. $`{\displaystyle \frac{\alpha \beta \text{ }\gamma }{\alpha \text{ }\beta \gamma }}`$ (S)
3. $`{\displaystyle \frac{\alpha \text{ }}{\alpha \beta \text{ }}}`$
4. $`{\displaystyle \frac{\alpha \beta \text{ }}{\alpha \text{ }\neg \beta }}`$
5. $`{\displaystyle \frac{\alpha \beta \text{ }}{\alpha \text{ }}}`$
Proof. Rules 12 and 3 were introduced and showed to be derived in a preferential relation in . For 4, suppose $`\alpha \beta \text{ }`$. Applying S, we get $`\alpha \text{ }\beta `$, and, by Right weakening, we conclude 4. For 5, suppose $`\alpha \beta \text{ }`$. Then, by 3, we get $`(\alpha \beta )\alpha \text{ }`$, and, by Left Logical Equivalence, we conclude 5.
Theorem 3 Let $`\text{ }`$ be a preferential inference relation based on $``$. Then $`\text{ }`$ is a preferential inference relation based on $`^{}`$ that satisfies the Consistency Preservation rule, where
$$\alpha ^{}\beta \text{iff}\mathrm{\Gamma },\alpha \beta ,$$
and
$$\mathrm{\Gamma }=\{\neg \gamma :\gamma \text{ }\}.$$
Proof. We must prove that $`\text{ }`$ satisfies Supraclassicality, Left Logical Equivalence, Right Weakening and Consistency Preservation with respect to $`^{}`$. The rest of the properties are already satisfied since $`\text{ }`$ is a rational inference relation.
First notice that Consistency Preservation is immediate by definition of $`^{}`$.
For Supraclassicality, suppose $`\alpha ^{}\gamma `$ then $`\mathrm{\Gamma },\alpha \gamma `$. By compactness of $``$, there exists $`\beta _1,\mathrm{},\beta _n`$ such that $`\beta _1\text{ },\mathrm{},\beta _n\text{ }`$ and $`\neg \beta _1,\mathrm{},\neg \beta _n,\alpha \text{ }\gamma `$. By repeated applications of Or, we get $`\beta _1\mathrm{}\beta _n\text{ }`$. Let $`\beta =\beta _1\mathrm{}\beta _n`$, then $`\beta \text{ }`$ and $`\alpha \neg \beta \gamma `$. By Supraclassicality of $`\text{ }`$ on $``$ we have $`\alpha \neg \beta \text{ }\gamma `$. By Lemma 2.3, we have $`\alpha \beta \text{ }`$, so, by Lemma 2.4, we have $`\alpha \text{ }\neg \beta `$. Using Cut, we get $`\alpha \text{ }\gamma `$.
For Left Logical Equivalence, suppose $`\alpha \text{ }\gamma `$, $`\alpha ^{}\beta `$ and $`\beta ^{}\alpha `$, i.e. $`\mathrm{\Gamma },\alpha \beta `$ and $`\mathrm{\Gamma },\beta \alpha `$. By compactness there exist $`\delta _1,\delta _2`$ such that $`\alpha \text{ }\neg \delta _1`$, $`\alpha \text{ }\neg \delta _2`$, $`\alpha \delta _1\beta `$ and $`\beta \delta _2\alpha `$. As above we have $`\alpha \text{ }\beta `$ and $`\beta \text{ }\alpha `$. Therefore by Lemma 2.1 we get $`\beta \text{ }\gamma `$.
Coming to Right Weakening, suppose $`\alpha \text{ }\beta `$ and $`\beta ^{}\gamma `$, i.e. there exists $`\delta `$ such that $`\alpha \text{ }\neg \delta `$ and $`\beta \neg \delta \text{ }\gamma `$. By And we have $`\alpha \text{ }\beta \neg \delta `$, so using Right Weakening of $`\text{ }`$ on $``$ we get $`\alpha \text{ }\gamma `$.
For Theorem 10, we shall need the following lemma.
Lemma 9 Let $``$ and $`\text{ }`$ be a rational ordering and inference relation, respectively. Then
1. If $`\alpha \text{ }`$ then $`\beta ^{}\neg \alpha `$, for all $`\beta `$, where $`^{}=𝐎(\text{ })`$.
2. If $`\beta \neg \alpha `$, for all $`\beta `$, then $`\alpha \text{ }^{}`$, where $`\text{ }^{}=𝐂()`$.
3. $`\neg \alpha <\alpha \gamma `$ iff $`\{\beta :\neg \alpha <\beta \}\alpha \gamma `$.
4. If $`\alpha \text{ }^{}`$ then $`\beta \neg \alpha `$, for all $`\beta `$, where $`\text{ }^{}=𝐂()`$.
5. If $`\beta ^{}\neg \alpha `$, for all $`\beta `$, then $`\alpha \text{ }`$, where $`^{}=𝐎(\text{ })`$.
Proof. Part 2 is immediate from defining condition ($`𝐂`$). For Part 1, suppose that $`\neg (\neg \alpha \beta )\text{ }\beta `$. We must show $`\neg (\neg \alpha \beta )\text{ }`$. Since $`\alpha \text{ }`$, we have $`\alpha \text{ }\neg \beta `$, by Right Weakening. Applying Or, we get $`\alpha \neg \beta \text{ }\neg \beta `$. By hypothesis and And, we get $`\neg (\neg \alpha \beta )\text{ }`$.
The left to right direction of Part 3 is straightforward. For the right to left direction, suppose $`\{\beta :\neg \alpha <\beta \}\alpha \gamma `$. Then, by Compactness, there exist $`\beta _1,\mathrm{},\beta _n`$ such that $`\neg \alpha <\beta _i`$, for all $`i\{1,\mathrm{},n\}`$, and $`\beta _1\mathrm{}\beta _n\alpha \gamma `$. By Conjunctiveness, we have $`\neg \alpha <\beta _1\mathrm{}\beta _n`$. So, by Dominance, we have $`\beta _1\mathrm{}\beta _n\alpha \gamma `$. Hence, by Transitivity, $`\neg \alpha <\alpha \gamma `$, as desired.
For Part 4, we have, by definition of $`\text{ }^{}`$ that $`\beta \neg \alpha `$, for all $`\alpha `$ or there is $`\beta `$ such that $`\alpha \beta `$ and $`\neg \alpha <\beta `$. If the former holds then the result follows immediately. If the latter holds then $`\beta \neg \alpha `$ that contradicts $`\neg \alpha <\beta `$, by Dominance. margin: Do Part 5
###### Lemma 25
Let $``$ and $`\text{ }`$ be a rational ordering and inference relation, respectively. Then
1. $`\neg \alpha \neg \beta \text{ }\neg \alpha `$ implies $`\alpha ^{}\beta `$, where $`^{}=𝐎(\text{ })`$.
2. $`\alpha \text{ }^{}\beta `$ implies $`\neg \alpha \neg \beta \neg \alpha `$, where $`\text{ }^{}=𝐂()`$.
3. If $``$ satisfies Bounded Disjunction then $`\neg \alpha \neg \beta \neg \alpha `$ implies $`\alpha \text{ }^{}\beta `$, where $`\text{ }^{}=𝐂()`$.
Proof. For Part 1, if $`\neg \alpha \neg \beta \text{ }\sim ̸\alpha `$ we get immediately $`\alpha ^{}\beta `$, by Definition (O). If not, that is $`\neg \alpha \neg \beta \text{ }\alpha `$, then, by And and Right Monotonicity, we have $`\neg \alpha \neg \beta \text{ }`$. Again, by Definition (O), we have $`\alpha ^{}\beta `$.
For Part 2, assume $`\alpha \text{ }^{}\beta `$. If $`\alpha \beta `$. Then we have $`\neg \beta \neg \alpha `$, so $`\neg \alpha \neg \beta \neg \alpha `$, and hence, by Dominance, $`\neg \alpha \neg \beta \text{ }\neg \alpha `$, as desired. If not then there must be $`\gamma `$ such that $`\gamma \alpha \beta `$ and $`neg\alpha <\gamma `$. Therefore $`\neg \alpha <\neg \alpha \neg \beta `$. Now, suppose $`\neg \alpha \neg \beta \neg \alpha `$ towards a contradiction. By Connectivity, we have $`\neg \alpha <\neg \alpha \neg \beta `$ and so, by Conjunctiveness, $`\neg \alpha <(\alpha \beta )(\alpha \neg \beta )`$. Hence $`\neg \alpha <\neg \alpha `$, a contradiction to Reflexivity.
For Part 3, if $`\neg \alpha <\alpha \beta `$ then we immediately have $`\alpha \text{ }^{}\beta `$, by Definition (C). If not, that is $`\alpha \beta \neg \alpha `$, then applying Bounded Disjunction, we have $`\alpha \beta \neg \beta \neg \alpha `$. The latter implies $`\neg \alpha `$, so, by Dominance $`\gamma \neg \alpha `$, for all $`\gamma `$. Hence $`\alpha \text{ }^{}\beta `$, by Definition (C).
Theorem 10 Given a rational ordering $``$ and a rational inference relation $`\text{ }`$, then $`𝐂()`$ is a rational inference relation and $`𝐎(\text{ })`$ is a rational ordering. Moreover, we have $`\text{ }=𝐂(𝐎(\text{ }))`$ and $`=𝐎(𝐂())`$.
Proof. We shall try not to overlap with the proof of Gärdenfors and Makinson proof of Theorem 6 (see proof of Theorem 3.3 in ). Therefore we do not cover the case where the second half of condition (R$`\text{ }`$) applies. The list of rules we verify is Supraclassicality, Left Logical Equivalence, And, Cut, Cautious Monotony, Or and Rational Monotony. Right Weakening follows from the above list.
We shall first show that $`𝐂()`$ is a rational inference relation.
For Supraclassicality, suppose that $`\alpha \gamma `$ but not $`\beta \neg \alpha `$ for all $`\beta `$. So there exists $`\beta `$ such that $`\neg \alpha <\beta `$. But then $`\alpha \beta \gamma `$ and therefore $`\alpha \text{ }\gamma `$.
For Left Logical Equivalence, suppose that $`\alpha \text{ }\gamma `$, $`\alpha \beta `$ and $`\delta \neg \alpha `$ for all $`\delta `$. Since $`\beta \alpha `$ we have $`\neg \alpha \neg \beta `$. By Dominance we get $`\neg \alpha \neg \beta `$ and by Transitivity $`\delta \neg \beta `$ for all $`\delta `$. Therefore $`\beta \text{ }\gamma `$.
For And, suppose that $`\alpha \text{ }\beta `$ and $`\alpha \text{ }\gamma `$. In case $`\delta \neg \alpha `$ for all $`\delta `$ we have immediately $`\alpha \text{ }\beta \gamma `$.
Turning to Or, suppose that $`\alpha \text{ }\gamma `$ and $`\beta \text{ }\gamma `$. If $`\delta \neg \alpha `$ for all $`\delta `$ and $`\delta \neg \beta `$ for all $`\delta `$, then by Conjunctiveness we have either $`\neg \alpha \neg \alpha \neg \beta `$ or $`\neg \beta \neg \alpha \neg \beta `$. In either case $`\delta \neg \alpha \neg \beta `$ for all $`\delta `$ by Transitivity. Therefore $`\neg (\alpha \beta )\text{ }\gamma `$. In the mixed case, say $`\delta \neg \alpha `$ for all $`\delta `$ and there exists $`\delta _0`$ such that $`\beta \delta _0\gamma `$ and $`\neg \beta <\delta _0`$, we have $`(\alpha \beta )(\neg \alpha \delta _0)\gamma `$. Now suppose that $`\neg \alpha \delta _0<\delta _0`$. By Conjunctiveness we must have $`\delta _0\neg \alpha \neg \alpha \delta _0`$, a contradiction. Thus $`\neg (\alpha \beta )\beta <\delta _0\neg \alpha \delta _0`$. Therefore $`\alpha \beta \text{ }\gamma `$.
For Cut, suppose that $`\alpha \text{ }\beta `$ and $`\alpha \beta \text{ }\gamma `$. If $`\delta \neg \alpha `$ for all $`\delta `$ then, by definition, $`\alpha \text{ }\gamma `$. If not, there exists $`\delta _0`$ such that $`\alpha \delta _0\beta `$ and $`\neg \alpha <\delta _0`$. Now suppose that $`\delta \neg (\alpha \gamma )`$ for all $`\delta `$. Observe that $`\alpha [(\neg \alpha \neg \beta )\delta _0]\gamma `$. We moreover have that $`\neg \alpha <\delta _0(\neg \alpha \neg \beta )\delta _0`$. Therefore $`\alpha \text{ }\gamma `$.
For Rational Monotonicity, suppose that $`\alpha \text{ }\gamma `$ and $`\alpha \text{ }\sim ̸\neg \beta `$. If $`\delta \neg \alpha `$ for all $`\delta `$, then we get a contradiction because $`\alpha \text{ }\neg \beta `$.
For Cautious Monotony, suppose that $`\alpha \text{ }\beta `$ and $`\alpha \text{ }\gamma `$. Observe that in case $`\alpha \text{ }\sim ̸\neg \beta `$ then the result follows by an application of Rational Monotony. If not, i.e. $`\alpha \text{ }\neg \beta `$, then by applying And we have $`\alpha \text{ }`$. If $`\delta \neg \alpha `$ for all $`\delta `$, then, since $`\neg \alpha \neg \alpha \neg \beta `$, we have $`\delta \neg \alpha \neg \alpha \neg \beta `$ but $`\neg \alpha \neg \beta \neg (\alpha \beta )`$ therefore $`\alpha \beta \text{ }\gamma `$. Otherwise there exists $`\delta `$ such that $`\alpha \delta `$ and $`\neg \alpha <\delta `$. But then we have that $`\delta \neg \alpha `$ and therefore $`\delta \neg \alpha `$ which is a contradiction to our hypothesis.
Definition ($`𝐎`$) is identical to Gärdenfors and Makinson’s one in the second disjunct. Therefore we shall only treat the first disjunct.
For Dominance, suppose $`\alpha \beta `$ and $`\neg (\alpha \beta )\text{ }\alpha `$. We have $`\neg \beta \neg \alpha `$ and $`\neg \alpha \neg \alpha `$. By Or, we get $`\neg \beta \neg \alpha \neg \alpha `$. By Supraclassicality, we have $`\neg (\alpha \beta )\text{ }\neg \alpha `$. Applying And, we get $`\neg (\alpha \beta )\text{ }`$.
For Conjunctiveness, suppose $`\neg (\alpha (\alpha \beta ))\text{ }\alpha `$ and $`\neg (\beta (\alpha \beta ))\text{ }\beta `$. These imply $`\neg (\alpha \beta )\text{ }\alpha `$ and $`\neg (\alpha \beta )\text{ }\beta `$, by Left Logical Equivalence. Applying And, we get $`\neg (\alpha \beta )\text{ }\alpha \beta `$. By reflexivity of $``$ and And, we have $`\neg (\alpha \beta )\text{ }`$. By Left Logical Equivalence again, we have $`\neg (\alpha (\alpha \beta ))\text{ }`$, and so $`\neg (\beta (\alpha \beta ))\text{ }`$, as desired.
For Transitivity, let $`\alpha \beta `$ and $`\beta \gamma `$. Suppose $`\neg (\alpha \beta )\text{ }`$. By Lemma 2.5, we have $`\neg \beta \text{ }`$. Lemma 2.3 gives $`\neg (\beta \gamma )\neg \beta \text{ }`$. By S, we have $`\neg (\beta \gamma )\text{ }\beta `$. So, we have $`\neg (\beta \gamma )\text{ }`$. Using the initial hypothesis and Or, we get $`\neg \alpha \neg \beta \neg \gamma \text{ }`$. By Lemma 2.5, we have $`\neg \alpha \neg \gamma \text{ }`$, i.e. $`\neg (\alpha \gamma )\text{ }`$. Now, suppose $`\neg (\beta \gamma )\text{ }`$ and $`\neg (\alpha \gamma \text{ }\sim ̸\alpha `$. Then, by Lemma 2.5, we have $`\neg \gamma \text{ }`$. Since $`\neg \alpha \text{ }\neg \alpha `$, we have $`\neg \alpha \neg \gamma \text{ }\neg \alpha `$. Therefore if $`\neg (\alpha \gamma )\text{ }\alpha `$, then And gives $`\neg (\alpha \gamma )\text{ }`$.
We shall now show that the initial rational inference relation $`\text{ }`$ and the induced one $`\text{ }_{}`$ by the expectation ordering with ($`𝐂`$) are the same.
We show first that $`\text{ }\text{ }_{}`$. Let $`\alpha \text{ }\gamma `$. We must show that $`\alpha \text{ }_{}\gamma `$. If $`\delta \neg \alpha `$, for all $`\delta `$, then it clearly holds. If not, let $`\beta \neg \alpha \gamma `$ then $`\alpha \beta \alpha \gamma `$. So $`\alpha \beta \gamma `$. Also, $`\alpha \neg \beta \alpha `$. If $`\alpha \text{ }`$ then $`\delta \neg \alpha `$, for all $`\beta `$ (using Lemma 9.1), so, by our hypothesis, $`\alpha \neg \beta \text{ }\sim ̸`$. Observe that $`\alpha \neg \beta \text{ }\gamma `$, and $`\gamma \neg \alpha \gamma \beta `$. Right Weakening gives $`\alpha \neg \beta \text{ }\beta `$. So $`\beta \neg \alpha `$ and therefore $`\neg \alpha <\beta `$. Hence $`\alpha \text{ }_{}\gamma `$.
For the other direction, i.e. $`\text{ }_{}\text{ }`$, let $`\alpha \text{ }_{}\gamma `$. Suppose first that $`\beta \neg \alpha `$ for all $`\beta `$. Therefore $`\neg \alpha `$. This gives either $`\neg (\neg \alpha )\text{ }`$ or $`\neg (\neg \alpha )\text{ }\sim ̸`$. Since obviously the latter does not hold we must have that $`\alpha \text{ }`$ and hence, by Right Weakening, $`\alpha \text{ }\gamma `$. Now suppose that there exists $`\beta `$ with $`\alpha \beta \gamma `$ and $`\neg \alpha <\beta `$, i.e. $`\beta \neg \alpha `$. The latter implies that $`\alpha \neg \beta \text{ }\sim ̸`$ and $`\alpha \neg \beta \text{ }\beta `$. Observe that $`\alpha \beta (\alpha \neg \beta )\beta `$ hence, by Supraclassicality, $`(\alpha \neg \beta )\beta \text{ }\gamma `$. Applying Cut on the latter and $`\alpha \neg \beta \text{ }\beta `$ we get $`\alpha \neg \beta \text{ }\gamma `$. Applying And on $`\alpha \neg \beta \text{ }\beta `$ and $`\alpha \neg \beta \text{ }\alpha \neg \beta `$ we get $`\alpha \neg \beta \text{ }\alpha `$. Since $`\alpha \neg \beta \text{ }\sim ̸`$ we have that $`\alpha \neg \beta \text{ }\sim ̸\neg \alpha `$. Applying Rational Monotonicity on the latter and $`\alpha \neg \beta \text{ }\gamma `$ we get $`(\alpha \neg \beta )\alpha \text{ }\gamma `$, i.e. $`\alpha \text{ }\gamma `$.
It remains to show $`=𝐎(𝐂())`$. Let $`^{}`$ be $`𝐎(𝐂())`$ and assume $`\alpha ^{}\beta `$. By definition of $``$, we have that $`\neg (\alpha \beta )\text{ }`$ or $`\neg (\alpha \beta )\text{ }\sim ̸\alpha `$, where $`\text{ }=𝐂()`$. The former implies $`\gamma \alpha \beta `$, for all $`\gamma `$, by Lemma 9. By Dominance and Transitivity, we have $`\alpha \alpha \beta \beta `$, as desired. The latter implies that $`\neg (\alpha \beta )\alpha \alpha \beta `$. Since $`\neg (\alpha \beta )\alpha `$ is (classically) equivalent to $`\alpha `$, we get $`\alpha \alpha \beta `$ and so $`\alpha \beta `$, by Dominance and Transitivity.
For the other direction, assume $`\alpha \beta `$. If $`\gamma \alpha \beta `$, for all $`\gamma `$ then $`\neg (\alpha \beta )\text{ }`$, by Lemma 9.2, and therefore $`\alpha ^{}\beta `$, by Definition (O). If not then, by Conjunctiveness on the hypothesis, we have $`\alpha \alpha \beta \beta `$. Since $`\neg (\alpha \beta )\alpha `$ is (classically) equivalent to $`\alpha `$, we get $`\neg (\alpha \beta )\alpha \alpha \beta `$. So $`\alpha \beta \neg (\alpha \beta )\alpha `$ and so, by Definition (C) and Lemma 9.3, we have $`\neg (\alpha \beta )\text{ }\sim ̸\alpha `$. The latter implies, by Definition (O), $`\alpha ^{}\beta `$, as desired.
Proposition 16 A ranked consequence operator $`\text{ }`$ based on $``$ is induced by a chain of sets of formulas if and only if it is induced by the closure of this chain under arbitrary unions and intersections.
Proof. Let $`\text{ }`$ and $`\text{ }^{}`$ be the ranked consequence operators induced by $`\{A_i\}_{iI}`$ and $`\{A_j\}_{jJ}`$ respectively, where the latter is the closure of the former under arbitrary unions and intersections. Without loss of generality we can assume that the sets belonging in $`\{A_i\}_{iI}`$ carry the same indices in $`\{A_j\}_{jJ}`$. We must prove that $`\text{ }`$ is equal to $`\text{ }^{}`$.
From left to right, suppose $`\text{ }\beta `$ and there exists $`iI`$ such that $`\text{ }_i\beta `$. Since $`iJ`$ we also have $`\text{ }^{}\beta `$. Suppose $`\alpha \text{ }\beta `$. If $`\alpha \text{ }_i\beta `$ for some $`ii`$ then as above $`\alpha \text{ }^{}\beta `$. If $`\text{ }_i\neg \alpha `$ for all $`i`$, i.e. $`\neg \alpha A_i`$ for all $`i`$, then $`\neg \alpha A_j`$ for all $`jJ`$. For either $`A_j=_kA_{i_k}`$ or $`A_j=_kA_{i_k}`$, where $`i_kI`$, and $`\neg \alpha A_{i_k}`$ for all $`k`$.
From right to left, suppose $`\text{ }^{}\beta `$, then there is $`jJ`$ such that $`\text{ }_j\beta `$. Either $`A_j=_kA_{i_k}`$ or $`A_j=_kA_{i_k}`$, where $`i_kI`$. In both cases there is some $`k_0`$ such that $`\beta A_{i_{k_0}}`$. Therefore $`\text{ }\beta `$. Suppose now that $`\alpha \text{ }^{}\beta `$. If there exists $`jJ`$ such that $`\neg \alpha A_j`$ and $`\alpha \beta A_j`$ then either $`A_j=_kA_{i_k}`$ or $`A_j=_kA_{i_k}`$, where $`i_kI`$. In the first case there exists $`k_0`$ such that $`\alpha \beta A_{i_{k_0}}`$. We also have that $`\neg \alpha A_{i_k}`$ for all $`k`$. So for the same $`k_0`$ we have that $`\neg \alpha A_{i_{k_0}}`$. Therefore $`\alpha \text{ }\beta `$. In the second case we have that $`\alpha \beta A_{i_k}`$ for all $`k`$ while there exists $`k_0`$ such that $`\neg \alpha A_{i_{k_0}}`$. If $`\neg \alpha A_j`$ for all $`jJ`$, then we immediately have that $`\neg \alpha A_i`$ for all $`iI`$ and hence $`\alpha \text{ }\beta `$.
Proposition 19 An ordering induced by a ranked consequence operator is rational.
Proof. Let $`\{B_i\}_{iI},\text{ }`$ be a ranked consequence operator. Denote $`𝐎(\text{ })`$ with $``$, and $`𝐂()`$ with $`\text{ }_{}`$.
We should verify that $``$ satisfies Supraclassicality, Transitivity, and Conjunctiveness.
For Supraclassicality, suppose $`\alpha \beta `$. We have $`\alpha \beta `$, so if $`B_i\alpha `$ than $`B_i\beta `$, for all $`iI`$. Hence $`\alpha \beta `$.
For Transitivity, suppose $`\alpha \beta `$ and $`\beta \gamma `$. Pick an $`iI`$ such that $`B_i\alpha `$. We have $`B_i\beta `$, since $`\alpha \beta `$. Hence $`B_i\gamma `$, since $`\beta \gamma `$, as desired.
For Conjunctiveness, suppose $`\alpha \alpha \beta `$ and $`\beta \alpha \beta `$, towards a contradiction. By our assumptions, there exist $`B_i`$ and $`B_j`$ with $`i,jJ`$ such that $`B_i\alpha `$ and $`B_i⊬\alpha \beta `$ and $`B_j\beta `$ and $`B_j⊬\alpha \beta `$. Now, $`B_i`$’s form a chain under inclusion, so either $`B_jB_i`$, or $`B_iB_j`$. If $`B_jB_i`$ then $`B_i\beta `$, a contradiction, since $`B_i\alpha `$ and $`B_i⊬\alpha \beta `$. Similarly, for $`B_iB_j`$.
We must now show that $`\alpha \text{ }\beta `$ iff $`\alpha \text{ }_{}\beta `$. Assume $`\alpha \text{ }\beta `$. We have either $`\text{ }_i\neg \alpha `$, for all $`iI`$, or there exists $`iI`$ such that $`B_i⊬\neg \alpha `$ and $`B_i\alpha \beta `$. Assume the former. We have immediately $`\gamma \neg \alpha `$, for all $`\gamma `$. Hence $`\alpha \text{ }_{}`$, by Lemma 9.2. Assume the latter, then $`\alpha \beta \neg \alpha `$, that is, $`\neg \alpha <\alpha \beta `$. Hence $`\alpha \text{ }\beta `$, by Lemma 9.3. The other direction is similar.
Corollary 20 A ranked consequence operator is a rational inference relation.
Proof. We shall give an alternative proof with a straightforward verification of the rules of rational inference. We shall show that a ranked consequence operator $`\text{ }`$ based on $``$ induced by a chain of sets $`\{B_i\}_{iI}`$ satisfies Supraclassicality, Left Logical Equivalence, Right Weakening, And, Cut, Cautious Monotonicity, Or and Rational Monotonicity.
For Supraclassicality, suppose $`\alpha \beta `$. We have either $`\text{ }_i\neg \alpha `$ for all $`iI`$ or there exists some $`iI`$ such that $`\text{ }\sim ̸_i\neg \alpha `$. In the first case we have immediately $`\alpha \text{ }\beta `$. In the second case we have that $`B_i,\alpha \beta `$, by our hypothesis, and therefore $`\alpha \text{ }\beta `$.
For Left Logical Equivalence, suppose that $`\alpha \beta `$ and $`\alpha \text{ }\gamma `$. We have either $`\text{ }_i\neg \alpha `$, i.e. $`B_i\neg \alpha `$, for all $`iI`$ or there exists some $`iI`$ such that $`\text{ }\sim ̸_i\neg \alpha `$ and $`B_i,\alpha \gamma `$. Since $`\alpha `$ and $`\beta `$ are equivalent under $``$, we have in the first case that $`B_i\neg \beta `$ for all $`iI`$ and In the second case that there exists some $`iI`$ such that $`\text{ }\sim ̸_i\neg \beta `$ and $`B_i,\beta \gamma `$. In both cases we have $`\beta \text{ }\gamma `$.
For Right weakening, suppose $`\alpha \text{ }\beta `$ and $`\beta \gamma `$. If $`\text{ }_i\neg \alpha `$ for all $`iI`$ then we immediately get $`\alpha \text{ }\gamma `$. If there exists $`iI`$ such that $`\text{ }\sim ̸_i\neg \alpha `$ and $`B_i,\alpha \beta `$ then we also have that $`B_i,\alpha \beta \gamma `$ by hypothesis. Therefore $`B_i,\alpha \gamma `$ and $`\alpha \text{ }\gamma `$.
For And, suppose $`\alpha \text{ }\beta `$ and $`\alpha \text{ }\gamma `$. If $`\text{ }_i\neg \alpha `$ for all $`iI`$ then we immediately have $`\alpha \text{ }\beta \gamma `$. If not then there exists $`i,kI`$ such that $`\text{ }\sim ̸_i\neg \alpha `$, $`\text{ }\sim ̸_k\neg \alpha `$, $`B_i,\alpha \beta `$, and $`B_k,\alpha \gamma `$. Since $``$ is linear let $`ik`$. Then $`B_iB_k`$ and therefore $`B_k,\alpha \beta `$. So $`B_k,\alpha \beta \gamma `$ and $`\alpha \text{ }\beta \gamma `$.
For Cut, suppose $`\alpha \text{ }\beta `$ and $`\alpha \beta \text{ }\gamma `$. If $`\text{ }_i\neg \alpha `$ for all $`iI`$ then we immediately have $`\alpha \text{ }\gamma `$. If not, then there exists $`iI`$ such that $`\text{ }\sim ̸_i\neg \alpha `$ and $`B_i,\alpha \beta `$. If $`\text{ }_j\neg (\alpha \beta )`$( $`\neg \alpha \neg \beta `$) for all $`jI`$, then for $`i`$, $`B_i\neg \alpha \neg \beta `$, i.e. $`B_i,\alpha \neg \beta `$. Combining with our hypothesis we get $`B_i\neg \alpha `$, i.e. $`\text{ }_i\neg \alpha `$, a contradiction. Therefore there exists $`kI`$ such that $`\text{ }\sim ̸_k\neg \alpha \neg \beta `$ and $`B_k,\alpha \beta \gamma `$. There are two cases: either $`ki`$ or $`ik`$. In the first case we have $`B_i,\alpha \beta \gamma `$ as well, so by (regular) cut on $``$ and our hypothesis we get $`B_i,\alpha \gamma `$. Therefore $`\alpha \text{ }\gamma `$. In the second case, observe that $`B_k,\alpha \beta `$ so as above $`B_k,\alpha \gamma `$. Since $`B_k⊬\neg \alpha \neg \beta `$ then $`B_k⊬\neg \alpha `$ so again $`\alpha \text{ }\gamma `$.
For Cautious Monotony, suppose $`\alpha \text{ }\beta `$ and $`\alpha \text{ }\gamma `$. If $`\text{ }_i\neg \alpha `$ for all $`iI`$ then we also have $`\text{ }_i\neg \alpha \neg \beta `$ hence $`\text{ }_i\neg (\alpha \beta )`$ for all $`iI`$. Therefore $`\alpha \beta \text{ }\gamma `$. If not, there exists $`i,kI`$ such that $`\text{ }\sim ̸_i\neg \alpha `$, $`B_i,\alpha \beta `$, $`\text{ }\sim ̸_j\neg \alpha `$ and $`B_j,\alpha \gamma `$. Let $`l=\mathrm{max}(i,k)`$ then $`B_k⊬\neg \alpha `$ and both $`B_l,\alpha \beta `$ and $`B_l,\alpha \gamma `$. From the latter we get $`B_l,\alpha \beta \gamma `$. Suppose that $`B_k\neg \alpha \neg \beta `$ then $`B_l,\alpha \neg \beta `$. Combining it with above we get $`B_l\neg \alpha `$ which is a contradiction to our hypothesis.
For Or, suppose that $`\alpha \text{ }\gamma `$ and $`\beta \text{ }\gamma `$. If both $`\text{ }_i\neg \alpha `$ and $`\text{ }_i\neg \beta `$ for all $`iI`$ then we also have $`\text{ }_i\neg \alpha \neg \beta `$ for all $`i`$ and therefore $`\alpha \beta \text{ }\gamma `$. If this is true for only one of them, say $`\text{ }_i\neg \alpha `$ for all $`iI`$ but there exists $`kI`$ such that $`\text{ }\sim ̸_k\neg \beta `$ and $`B_k,\beta \gamma `$, then we also have $`\text{ }\sim ̸_k\neg \beta \neg \alpha `$ and $`B_k`$ which implies $`B_k,\alpha \gamma `$. So $`B_k,\alpha \beta \gamma `$ and therefore $`\alpha \beta \text{ }\gamma `$. If neither of them holds then there exist $`i,kI`$ such that $`\text{ }\sim ̸_i\neg \alpha `$, $`B_i,\alpha \gamma `$, $`\text{ }\sim ̸_k\neg \beta `$ and $`B_k,\beta \gamma `$. Let $`l=\mathrm{max}(i,k)`$ then we have $`\text{ }\sim ̸_l\neg \alpha \neg \beta `$, $`B_l,\alpha \gamma `$ and $`B_l,\beta \gamma `$. So $`B_l,\alpha \beta \gamma `$ and therefore $`\alpha \beta \text{ }\gamma `$.
For Rational Monotonicity, suppose that $`\alpha \text{ }\sim ̸\neg \beta `$ and $`\alpha \text{ }\gamma `$. We can’t have $`\text{ }_i\neg \alpha `$ for all $`iI`$ because we have that $`\alpha \text{ }\sim ̸\neg \beta `$. So there exists $`iI`$ such that $`\text{ }\sim ̸_i\neg \alpha `$ and $`B_i,\alpha \gamma `$. We have by monotonicity of $``$ that $`B_i,\alpha \beta \gamma `$. Now observe that $`\alpha \text{ }\sim ̸\neg \beta `$ implies that if $`B_i,\alpha \neg \beta `$ then $`\text{ }_i\neg \alpha `$. However since $`\text{ }\sim ̸_i\neg \alpha `$ we have that $`B_i,\alpha ⊬\neg \beta `$ so $`B_i⊬\neg \alpha \neg \beta `$. Therefore $`\alpha \beta \text{ }\gamma `$ by definition.
Theorem 21 A rational inference relation is a ranked consequence operator.
Proof. Denote the comparative rational inference relation by $`\text{ }_{}`$. We shall define a chain of sets $`\{A_i\}_{iI}`$ which generates a ranked consequence operator $`\text{ }`$ equal to $`\text{ }_{}`$. Let $``$ be the equivalence relation induced by $``$ (an expectation ordering is clearly a preorder). The equivalence classes will be denoted by $`\widehat{\alpha }`$ (where $`\alpha \widehat{\alpha }`$). It is also clear that the set of equivalence classes is linearly ordered. Now, for each $`\alpha `$, let
$$A_{\widehat{\alpha }}=\{\beta :\alpha \beta \}.$$
Note here that, by Dominance, the sets $`A_{\widehat{\alpha }}`$ are closed under consequence. Moreover, we have $`A_{\widehat{\alpha }}A_{\widehat{\beta }}`$ iff $`\beta \alpha `$. Now, generate a ranked consequence operator $`\text{ }`$ as in Definition 4.
We must show that $`\text{ }`$ and $`\text{ }_{}`$ are identical.
Let $`\alpha \text{ }_{}\beta `$, i.e. either $`\beta \neg \alpha `$ for all $`\beta `$, or there exists $`\delta `$ such that $`\neg \alpha <\delta `$ and $`\delta \alpha \gamma `$. In the first case $`\neg \alpha A_{\widehat{\beta }}`$ for all $`\beta `$. So $`\text{ }_{\widehat{\beta }}\neg \alpha `$ for all $`\beta `$ and therefore $`\alpha \text{ }\gamma `$. in the second case consider $`A_{\widehat{\delta }}`$. We have $`\delta A_{\widehat{\delta }}`$ so $`A_{\widehat{\delta }},\alpha \gamma `$. Suppose that $`A_{\widehat{\delta }}\neg \alpha `$ then by compactness there exists $`\delta ^{}A_{\widehat{\delta }}`$ such that $`\delta ^{}\neg \alpha `$. By Dominance we have $`\delta ^{}\neg \alpha `$ and by definition of $`A_{\widehat{\delta }}`$, $`\delta \delta ^{}`$, i.e. $`\delta \neg \alpha `$, a contradiction.
Now let $`\alpha \text{ }\gamma `$. If $`\text{ }_{\widehat{\beta }}\neg \alpha `$ for all $`\beta `$ then $`\beta \neg \alpha `$ for all $`\beta `$ and we are done. If not, then there exists $`\widehat{\delta }`$ such that $`A_{\widehat{\delta }}⊬\neg \alpha `$ and $`A_{\widehat{\delta }},\alpha \gamma `$. Since $`A_{\widehat{\delta }}⊬\alpha `$ we must have $`\neg \alpha <\delta `$. By compactness there exists $`\delta ^{}A_{\widehat{\delta }}`$ such that $`\delta ^{},\alpha \gamma `$, i.e. $`\delta ^{}\alpha \gamma `$ and $`\neg \alpha <\delta \delta ^{}`$. Therefore $`\alpha \text{ }\gamma `$.
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# The impact of atomic precision measurements in high energy physics 1footnote 11footnote 1 Talk given at the XVII International Conference on Atomic Physics, ICAP 2000, Florence, June 4-9, 2000
## Introduction
The aim of this talk is to review some of the atomic precision measurements in atomic physics leading to precious informations in the realm of high-energy physics. The idea of atomic physics bringing light on the high-energy physics world requires some qualification due to the very different scales of energy involved in the two cases. In fact, typically one has a separation of about six or seven order of magnitude between the two scales and one expects the two physics being almost decoupled. In fact, if we look at some observable, $`A`$, at a scale $`\mathrm{\Lambda }_1\mathrm{\Lambda }_2`$, we expect that the observable can be represented in the form
$$A(\mathrm{\Lambda }_1,\mathrm{\Lambda }_2)=A(\mathrm{\Lambda }_1)+𝒪\left(\left(\frac{\mathrm{\Lambda }_1}{\mathrm{\Lambda }_2}\right)^n\right)$$
(1)
In order to be able to derive informations about the physics at the scale $`\mathrm{\Lambda }_2`$, being at the scale $`\mathrm{\Lambda }_1`$, one starts considering a combination of observables corresponding to the corrections coming from the higher scale
$$B=c\left(\frac{\mathrm{\Lambda }_1}{\mathrm{\Lambda }_2}\right)^n$$
(2)
In order to measure $`B`$ one needs either the coefficient $`c`$ being very large in such a way to partially compensate the scale factor, or having an extremely good experimental sensitivity. In this talk I will consider two particular examples of situations where atomic physics can be relevant to high-energy physics, namely Atomic Violation of Parity (APV) and possible violations of the discrete symmetry CPT, that is the product of charge-conjugation, parity and time-reversal. In fact, using heavy atoms like cesium in APV measurements one can get a good enhancement factor. On the other hand, CPT symmetry can be tested by using the extraordinary opportunities offered by the atomic traps in order to obtain a very accurate determination of frequencies.
## Atomic Parity Violation in Atoms
In this Section I will discuss mainly the latest determination of the weak charge in atomic cesium and some of its implications in models of physics beyond the Standard Model (SM). The SM has been tested very precisely at machines such as LEP and SLC, where, working at an energy around the $`Z`$ mass, one is mainly testing the property of the $`Z`$ itself. Therefore, the physics beyond the SM that can be looked for at these machines is the one giving corrections to the $`Z`$-propagator and/or to the couplings of the $`Z`$ with fermion-antifermion pairs. Namely, new massive vector bosons, $`Z^{}`$, which mix to the $`Z`$, or new particles running in loops and contributing to the $`Z`$ self-energy or to vertex corrections. But consider, for instance, the case of a massive vector boson which does not mix to the $`Z`$, and therefore invisible at LEP (except for tiny radiative corrections). If the $`Z^{}`$ is coupled to fermions, in the low-energy limit it gives rise to an effective four-fermi interaction. Therefore, low-energy experiments are complementary to the high-energy ones, and furthermore they are able to measure directly the couplings of the $`Z`$ to light quarks; something that at LEP and SLC can be done only in an indirect way. Among the low-energy experiments a particular role is played by the APV experiments, due to the precision almost at the level of the one reached at LEP/SLC.
Let us now recall some feature of APV in atoms. First of all, within the SM the four-fermi parity violating hamiltonian density for nucleons is given by
$$^{PV}=\frac{G_F}{\sqrt{2}}\left[(\overline{e}\gamma _\mu \gamma _5e)\underset{N=p,n}{}c_{1N}\overline{N}\gamma ^\mu N+(\overline{e}\gamma _\mu e)\underset{N=p,n}{}c_{2N}\overline{N}\gamma ^\mu \gamma _5N\right]$$
(3)
where
$$c_{ip}=2c_{iu}c_{id},c_{in}=c_{iu}2c_{id},i=1,2$$
(4)
and
$$c_{1q}=8a_ev_q=(T_3^q2s_\theta ^2Q^q),c_{2q}=8v_ea_q=T_3^q(14s_\theta ^2),q=u,d$$
(5)
Here $`v_e`$, $`v_q`$, $`a_e`$ and $`a_q`$ are the vector and vector-axial couplings of the $`Z`$ to the electrons and quarks. For a point-like nucleus with $`Z`$ protons and $`N`$ neutrons, the hamiltonian density, in the non-relativistic limit, is given by
$`_{PV}`$ $`=`$ $`{\displaystyle \frac{G_F}{4\sqrt{2}m_e}}[Q_W(Z,N)\stackrel{}{\sigma }_{\mathrm{}}[\stackrel{}{p},\delta ^3(\stackrel{}{r})]_++2(c_{2p}\stackrel{}{S}_p+c_{2n}\stackrel{}{S}_n)[\stackrel{}{p},\delta ^3(\stackrel{}{r})]_+`$ (6)
$``$ $`2i\stackrel{}{\sigma }_{\mathrm{}}(c_{2p}\stackrel{}{S}_p+c_{2n}\stackrel{}{S}_n)[\stackrel{}{p},\delta ^3(\stackrel{}{r})]_+]`$ (7)
where $`\stackrel{}{p}`$ is the momentum of the electron, $`\stackrel{}{S}_{p(n)}`$ the total spin of the protons (neutrons) and $`m_e`$ the electron mass. I have also defined the weak charge of the atom as
$$Q_W(Z,N)=2\left[c_{1p}Z+c_{1n}N\right]$$
(8)
Notice that for a heavy atom (large values of $`Z`$) the matrix element of the first term in $`_{PV}`$ is roughly proportional to $`Z^3`$, one factor coming from $`Q_W`$, one from the momentum of the electron and the third one from the wave function evaluated at the origin. This coherence effect was noticed by Bouchiat and Bouchiat bouchiat2 and it provides, in the case of cesium ($`Z=55`$) an enhancement factor of about $`10^5`$, more or less what is necessary in order to compensate for the decoupling factor from the scales mentioned in the Introduction.
In order to get a rough idea of the bounds on new physics that can be obtained by a measurement of $`Q_W`$ with a given sensitivity, we parametrize the new physics contribution to $`Q_W`$ by a four-fermi effective interaction musolf
$$_{NP}^{PV}=\frac{g_{NP}^2}{\mathrm{\Lambda }^2}\overline{e}\gamma _\mu \gamma _5e\underset{q=u,d}{}h_{1q}\overline{q}\gamma ^\mu q$$
(9)
If we assume $`h_{1q}c_{1q}`$, for a sensitivity $`\mathrm{\Delta }Q_W/Q_W1\%`$ one gets a bound
$$\mathrm{\Lambda }(5g_{NP})TeV$$
(10)
If new physics is strongly interacting ($`g_{NP}^24\pi `$), then $`\mathrm{\Lambda }17TeV`$, whereas in the weakly interacting case ($`g_{NP}^24\pi \alpha `$) we get $`\mathrm{\Lambda }1.5TeV`$. In any case we see that at 1% level of sensitivity, $`Q_W`$ is able to test new physics for scales greater than 1 $`TeV`$.
In APV measurements one looks at optical transitions between a pair of states $`|\psi _\pm `$ mixed by $`_{PV}`$ and a state $`|\psi _0`$ of the same nominal parity as $`|\psi _+`$. The mixing of the two eigenstates of parity is given by
$$\eta =\frac{\psi _{}|H_{PV}|\psi _+}{\mathrm{\Delta }E}$$
(11)
where $`\mathrm{\Delta }E`$ is the splitting between the two levels. If I denote by $`M_1`$ and $`E_1^{PV}`$ the amplitudes for the two unperturbed transitions $`|\psi _+|\psi _0`$ and $`|\psi _{}|\psi _0`$, the transition probability, after the mixing, is given by
$$W=M_1^2+|E_1^{PV}|^2\pm 2Im(E_1^{PV})M_1$$
(12)
The choice of the sign depends on the helicity of the photon which is emitted or absorbed in the transition. In the actual experiment on cesium one measures the circular dichroism, that is the asymmetry for the absorption cross-section
$$\delta =\frac{\sigma _+\sigma _{}}{\sigma _++\sigma _{}}2\frac{Im(E_1^{PV})}{M_1}$$
(13)
Of course, the PV amplitude $`E_1^{PV}`$ is proportional to the mixing parameter $`\eta `$ and therefore measuring $`\delta `$ one can get the matrix element of the PV hamiltonian. These ideas have been applied in particular to the transition $`6S7S`$ in atomic cesium $`{}_{55}{}^{133}Cs`$ bouchiat ; boulder1 ; boulder2 , but also to other atoms as thallium thallium . The typical value of $`\delta `$ is $`10^4÷10^5`$, but there is a strong background which can be overcomed by letting the PV amplitude to interfere with a large electro-induced (Stark) transition. Eventually one extracts from the experiment the matrix element of $`H^{PV}`$ which is proportional to $`Q_W`$ times an atomic form factor $`\kappa _{PV}`$ which must be evaluated theoretically in order to extract the value of the weak charge. Therefore the measurement must be coupled with theoretical calculations of similar accuracy in order to get a precise determination of $`Q_W`$. In the case of atomic cesium the calculation of $`\kappa _{PV}`$ was performed independently by two groups dzuba ; blundell . This calculation is not an easy task, as one has to use many-body perturbation theory coupled with Hartree-Fock techniques. The theoretical errors are quite difficult to estimate. The authors of Refs. dzuba ; blundell did their estimate by looking at the differences between the theoretical and the experimental values of parity conserving quantities as dipole matrix elements and hyperfine splittings for the $`6S_{1/2}`$, $`7S_{1/2}`$, $`6P_{1/2}`$ and $`7P_{1/2}`$ states. In this way the error $`\mathrm{\Delta }\kappa _{PV}/\kappa _{PV}1\%`$ was obtained. After the new measurement of the weak charge of the cesium by the Boulder group boulder2 , which improved the accuracy of the previous experiment boulder1 by more than a factor five, Bennett and Wieman bennett re-examined the theoretical errors on $`\kappa _{PV}`$. In fact, since the time of the previous estimate there have been a number of new and more precise measurements of the quantities of interest. The result is that now the agreement is much better than before, and as a consequence Bennett and Wieman got the estimate $`\mathrm{\Delta }\kappa _{PV}/\kappa _{PV}0.4\%`$. It should be noticed that there is a third element which contributes to the extraction of $`Q_W`$ from the data. This is the Stark mixing-induced electric dipole moment amplitude, $`\beta `$. The experiments in Refs. bouchiat ; boulder1 were using a theoretical determination of $`\beta `$. In boulder2 the ratio $`M_{hf}/\beta `$ has been measured. The off-diagonal magnetic dipole moment induced by the hyperfine interaction is well known empirically and it is possible to extract a precise value for $`\beta `$. However, in a contribution to this Conference dzuba2 , the matrix element $`M_{hf}`$ has been accurately calculated with the result that the empirical formula for it should be corrected by a factor of 0.24% increasing the discrepancy with the SM (see later). I would like also to comment about some possible neglected contribution in the evaluation of the atomic form factor. It has been pointed out in ref. pollock that there could be a contribution arising from the difference of neutron and proton spatial distributions inside the nucleus. This contribution turns out to be very difficult to estimate, in fact it is quite model dependent. Most probably it could introduce a further error on $`Q_W(Cs)`$ of about 0.3. This would not change the conclusions in a very significant way. Another point has been raised recently in ref. derevianko . This author argues that the contribution from the Breit interaction (exchange of a transverse photon between two electrons) could have been underestimated. The Breit interaction contribution to the atomic form factor was estimated in dzuba and it was found to be very small. However in ref. derevianko it is found that the total effect, taking into account also second and third order contributions, is about twice the first order effect. As a consequence, if ”all” the higher order contributions could be shown to be negligible, the experimental measure would reconcile with the SM expectation for $`Q_W(Cs)`$. However, see also ref. kozlov .
To conclude this analysis I think that an evaluation of the atomic form factor by taking into account the next order in the many-body perturbative theory is highly desirable in order to settle the question. In any case I find of some interest to assume that the theoretical error is indeed at the level of $`0.4\%`$ in order to see which are the possible implications of the APV in high-energy physics.
I can start now to discuss the experimental results on $`Q_W(Cs)`$. It is interesting to recall the value obtained in boulder1 combined with the theoretical determination of $`\kappa _{PV}`$ dzuba ; blundell
$$Q_W(Cs)=71.04\pm (1.58)_{\mathrm{exp}}\pm (0.88)_{\mathrm{th}}$$
(14)
The total error of these measurement on $`Q_W(Cs)`$ is at 2.5% level of accuracy that, at that time, was comparable with the sensitivity obtained at LEP1. In fact, this determination of $`Q_W(Cs)`$ lead to the first indication that technicolor models, in their most simple version obtained from scaling of QCD, could not possibly fit the data. The new experimental result on $`Q_W(Cs)`$ boulder2 combined with the new determination of the theoretical error bennett gives
$$Q_W(Cs)^{\mathrm{exp}}=72.06\pm (0.28)_{\mathrm{exp}}\pm (0.34)_{\mathrm{th}}$$
(15)
A result at 0.6% level of accuracy. On the theoretical side, $`Q_W`$ can be expressed as marciano
$$Q_W(Cs)^{\mathrm{th}}=72.72\pm 0.13102ϵ_3^{\mathrm{rad}}+\delta _NQ_W$$
(16)
including hadronic-loop uncertainty. I use here the variables $`ϵ_i`$ (i=1,2,3) of ref. altarelli , which include the radiative corrections, in place of the set of variables $`S`$, $`T`$ and $`U`$ originally introduced in ref. peskin . In the above definition of $`Q_W^{\mathrm{th}}(Cs)`$ I have explicitly included only the Standard Model (SM) contribution to the radiative corrections. New physics (that is physics beyond the SM) contributions are represented by the term $`\delta _NQ_W`$. Also, I have neglected a correction proportional to $`ϵ_1^{\mathrm{rad}}`$. In fact, as well known marciano , due to the particular values of the number of neutrons ($`N=78`$) and protons ($`Z=55`$) in cesium, the dependence on $`ϵ_1`$ almost cancels out. For a top mass of 175 $`GeV`$ and $`m_H=100(300)GeV`$ the value of $`ϵ_3^{\mathrm{rad}}`$ is given by altarelli2
$$ϵ_3^{\mathrm{rad}}=5.110(6.115)\times 10^3$$
(17)
For $`m_H=100GeV`$, corresponding roughly to the lower experimental bound from direct search at LEP2 higgs , one gets
$$Q_W(Cs)^{\mathrm{exp}}Q_W(Cs)^{SM}=1.18\pm 0.46$$
(18)
giving rise to a deviation of about $`2.57SD`$. Furthermore, for increasing mass of the Higgs the discrepancy increases. Therefore, if we assume as being correct the experimental result, the theoretical evaluation of $`\kappa _{PV}`$ and the evaluation of the theoretical errors, we are forced to conclude that the SM is disfavored at 99% CL.
We can draw another conclusion, that is, that in order to explain the data on $`Q_W(Cs)`$ we need new physics not constrained by the LEP and SLC data. In fact, as an example let me consider a type of new physics visible at LEP as, for instance, contributing to the self-energy of the $`Z`$, the so called oblique corrections. In such a case one can write $`\delta _NQ_W(\mathrm{oblique})=102ϵ_{3N}`$, and in order to compensate for the discrepancy on $`Q_W(Cs)`$ one needs
$$ϵ_{3N}=(11.6\pm 4.5)\times 10^3$$
(19)
whereas from LEP and SLC data one can determine the sum
$$ϵ_3^{\mathrm{exp}}=ϵ_3^{\mathrm{rad}}+ϵ_{3N}=(4.19\pm 1)\times 10^3$$
(20)
Therefore one gets $`ϵ_{3N}10^3`$, one order of magnitude too small to explain the data on $`Q_W(Cs)`$.
I would like also recall the experimental result of APV on Thallium thallium
$$Q_W(Tl)^{\mathrm{exp}}=114.8\pm (1.2)_{\mathrm{exp}}\pm (3.4)_{\mathrm{th}}$$
(21)
This result is not as precise as the one on $`Cs`$, and in fact the total error is about 3%. At this level it is perfectly compatible with the SM prediction
$$Q_W(Tl)^{\mathrm{SM}}=116.7\pm 0.1$$
(22)
A new experiment on cesium is being planned in Paris but the experimental sensitivity is going to be lower than the one obtained in Boulder.
In Berkeley and Seattle there are plans for isotope ratio measurements. In this case the dependence on the atomic form factor would go away eliminating the theoretical error. However these ratios depend on the variation of the neutron density along the isotope chain. This would introduce errors at least twice as big as the experimental ones chen .
We are now in the position of discussing the implications of eq. (18) on new physics. Assuming that the contribution of new physics, $`\delta _NQ_W`$, is such to reproduce the experimental results, we can make use of eqs. (15) and (16) to write casalbuoni
$$Q_W(Cs)^{\mathrm{exp}}Q_W(Cs)^{\mathrm{th}}(m_H)=0.66+102ϵ_3^{\mathrm{rad}}(m_H)\delta _NQ_W\pm 0.46$$
(23)
For $`m_H=100GeV`$ at a 95% CL we find
$$0.28\delta _NQ_W2.08$$
(24)
Notice that the lower positive bound arises since the SM (corresponding to $`\delta _NQ_W=0`$) does not fit the experimental value of $`Q_W(Cs)`$ at this CL value. This is quite important since it implies an upper bound on the scale of new physics. For the same reason new physics with a contribution $`\delta _NQ_W<0`$ is not allowed. Also notice that lower and upper bounds both increase for increasing Higgs mass.
### Contact interactions from compositness.
A typical four-fermi operator in composite models contributing to the PV lagrangian is langacker ; casalbuoni
$$\pm \frac{g^2}{\mathrm{\Lambda }^2}\overline{e}\gamma _\mu \frac{1\gamma _5}{2}e\overline{q}\gamma ^\mu \frac{1\gamma _5}{2}q$$
(25)
The effect of this interaction is to modify the coefficients $`c_{1u,1d}`$
$$c_{1u,1d}c_{1u,1d}\frac{\sqrt{2}\pi }{G_F\mathrm{\Lambda }^2}$$
(26)
where, since composite models correspond to strongly interacting new physics, we have assumed $`g^2=4\pi `$. From
$$Q_W=2[(2Z+N)c_{1u}+(Z+2N)c_{1d}]$$
(27)
we see that the negative sign for the operator (25) is excluded. For the positive sign we get the bounds
$$12.1\mathrm{\Lambda }(TeV)32.9$$
(28)
The typical lower bound from high energy physics is about 3.5 $`TeV`$ PDG .
### Extra-dimension models.
In ref. pomarol a minimal extension to higher dimensions of the SM, with extra dimensions compactified, was considered. In this model the fermions live in a 4-dimensional subspace, the wall, whereas the gauge bosons live in the full D-dimensional space, the bulk. In general, there might be two Higgs fields, one living in the bulk, $`\varphi _1`$, and the other living on the wall, $`\varphi _2`$. The propagation of the gauge fields in the bulk is equivalent to the exchange of an infinite tower of Kaluza-Klein (KK) excitations with increasing mass. For example, for $`D=5`$, $`M=n/R`$, $`n=1,\mathrm{},\mathrm{}`$, with $`R`$ the compactification radius. If only the Higgs field $`\varphi _2`$ is present, the ordinary gauge bosons do not mix with the KK resonances and it is easy to see that the contribution of these modes to $`Q_W`$ is negative casalbuoni2 . Therefore the model does not fit the data on $`Q_W(Cs)`$. For the more general case of both Higgs fields present it has been shown casalbuoni2 that the LEP/SLC and $`Q_W(Cs)`$ experimental data are not compatible among them at 95% CL.
### Extra $`Z^{}`$ models.
The implications of models with an extra neutral vector boson $`Z^{}`$ for APV have been considered in the literature for quite a long time qwz ; casalbuoni2 ; erler . The $`Z^{}`$ has couplings comparable to the ones of the $`Z`$ in the SM and therefore this is an example of weakly interacting new physics. There is a continuum of such models characterized by an angle $`0^0\theta _690^0`$. To any value of $`\theta _6`$ it corresponds a different model. The 95% CL regions allowed by $`Q_W`$, in the plane $`(\theta _6,M_Z^{})`$, for different values of the Higgs mass, are shown in Fig. 1. In deriving these Figures the assumption of zero mixing between $`Z`$ and $`Z^{}`$ has been made. In the Figure are also shown three popular models: $`\eta `$ ($`\theta _652^0`$), $`\chi `$ ($`\theta _6=0^0`$), $`\psi `$ ($`\theta _6=90^0`$). We see that the $`\eta `$ and the $`\psi `$ models are not allowed by the data. The direct search at the Tevatron for a $`Z^{}`$ within the $`\chi `$ model gives a direct lower bound at 95% CL, $`M_Z^{}590GeV`$ (a similar bound holds for all these models) . Therefore this model is compatible with the data. A recent best fit to all the data (including APV) gives for the $`\chi `$ model the following results erler , $`M_Z^{}=812_{152}^{+339}GeV`$ and a mixing angle compatible with zero, $`\theta _M=(1.12\pm 0.80)\times 10^3`$.
## Atomic physics and CPT violation
The CPT theorem is one of the fundamental results in local relativistic field theories. Therefore the idea of possible violations of this theorem implies that some of the axioms of these theories should be reviewed. Let me recall here the exact statement of the theorem cpt : In a field theory satisfying
1. Locality
2. Lorentz invariance
3. Analiticity of the Lorentz group representations in the boost parameters
the CPT transformation is a symmetry of the theory itself.
The first two conditions say that one is dealing with a local relativistic field theory, whereas the third one is satisfied in any finite-dimensional representation of the Lorentz group. It is interesting to notice that unitary representations fail to be analytic and as a consequence the CPT theorem can be violated in this case. The first example of this situation dates back to Majorana majorana when he formulated a first order wave equation without negative-energy solutions. He was able to do that by making use of a unitary infinite-dimensional representation of the Lorentz group. Since this theory does not contain antiparticles the CPT symmetry is broken. However, the quarks and leptons described by the SM belong to finite-dimensional representation of the Lorentz group and therefore this does not seem a possible way to break the theorem. It seems also very hard to give up locality, since it guarantees the microcausality of the theory. Therefore, the only sensible way to avoid the consequences of the CPT theorem in a local field theory seems to break Lorentz invariance. A situation of this type could arise at a more fundamental level as in string theory, where it is possible that Lorentz invariance is spontaneously broken around the Planck mass, $`M_P`$ kostelecky1 . One can take into account these effects by writing down a local effective lagrangian with Lorentz and CPT breaking terms. These terms can be written as an expansion in derivatives over the Planck mass. For instance, considering a single fermion, the violating term can be written as
$$_v=\underset{n}{}\frac{g_n}{M_P^n}T\overline{\psi }\mathrm{\Gamma }(i)^n\psi $$
(29)
I have used a somewhat symbolic notation where $`\mathrm{\Gamma }`$ stays for a generic combination of Dirac matrices and $`T`$ is a constant tensor and I take the mass dimensions of $`g_n`$ as $`[g_n]=1`$. Furthermore I will assume the same internal symmetries as in the SM, that is $`SU(3)SU(2)U(1)`$ kostelecky2 ; kostelecky3 . Since the breaking terms should vanish in the limit $`M_P\mathrm{}`$ also for $`n=0`$, I will require
$$g_0=c_o\frac{m^2}{M_P}$$
(30)
where $`m`$ is some low-energy mass scale parameter. We see that the relevant terms are the ones with $`n=0`$ and $`n=1`$, and therefore the resulting theory preserves the renormalizability property.
Let me now consider a single fermion interacting with the electromagnetic field. One adds to the standard QED lagrangian the following two terms
$$_v^{(n=0)}=\overline{\psi }[a_\mu \gamma ^\mu b_\mu \gamma _5\gamma ^\mu \frac{1}{2}H_{\mu \nu }\sigma ^{\mu \nu }]\psi $$
(31)
and
$$_v^{(n=1)}=\overline{\psi }[ic_{\mu \nu }\gamma ^\mu D^\nu +id_{\mu \nu }\gamma _5\gamma ^\mu D^\nu ]\psi $$
(32)
where $`D_\mu =_\mu iqA_\mu `$, with $`q`$ the electric charge of the fermion. There are other possible terms with $`n=1`$, but they are not compatible with the symmetries of the SM and therefore they should be suppressed. The following orders of magnitude are expected
$$a_\mu ,b_\mu ,H_{\mu \nu }𝒪(m^2/M_P),c_{\mu \nu },d_{\mu \nu }𝒪(m/M_P)$$
(33)
The terms in $`_v^{(n=0,1)}`$ violate Lorentz invariance, since all the tensors in eq. (33) are constant ones. However only the terms proportional to $`a_\mu `$ and $`b_\mu `$ violate CPT symmetry since $`\gamma _\mu `$, $`\gamma _\mu \gamma _5`$ and $`D_\mu `$ are CPT odd, whereas the other covariant terms are CPT even. Therefore, in the following I will take into consideration only $`_v^{(n=0)}`$. Notice also that when dealing with a single fermion the term in $`a_\mu `$ does not have physical meaning since we can write $`a_\mu =_\mu (ax)`$, showing that $`a_\mu `$ is a trivial gauge background field. Of course, the situation changes when dealing with different fermions having different $`a_\mu `$’s. From eq. (33) we expect that the order of magnitude of the CPT and Lorentz breaking terms is given by $`m/M_P10^{22}÷10^{17}`$ for $`m=m_e÷v`$, where $`m_e`$ is the electron mass and $`v250GeV`$ is the electroweak symmetry breaking scale. Lorentz and CPT breaking terms could appear also in the photon part of the total lagrangian. This instance is discussed thoroughly in the second paper of ref. kostelecky2 , but I will not consider it in this talk.
Here I want to illustrate some atomic physics experiment about CPT violation. But before doing that let me just give a list of other existing or planned experiments about the violation of this fundamental symmetry
* $`K\overline{K}`$ mass difference. This experiment gives the best high-energy result PDG
$$\frac{|m_Km_{\overline{K}}|}{m_K}10^{18}$$
(34)
* Experiments on neutral meson oscillations to be done at meson factories kostelecky4 .
* Experiments on muons kostelecky5 .
* Experiments with spin-polarized solids bluhm1 .
* Experiments from clock-comparison kostelecky6 .
CPT violation may have also some relevance for baryogenesis and this subject has been discussed in ref. bertolami .
Let me now consider atomic physics experiments for testing CPT using atomic traps. Several of these experiments have been performed by confining single particles or antiparticles in a Penning trap for a long time. These experiments have a very high precision, of order $`10^9`$ or better, whereas the precision in experiments about mesons (see eq. (34)) is much lower, of order $`10^3`$. I recall here the comparison of the electron and positron gyromagnetic ratios, $`g_{}`$, obtained measuring their cyclotron and anomaly frequencies (see later), which gives the figure of merit dyck
$$\left|\frac{g_{}g_+}{g_{\mathrm{av}}}\right|2\times 10^{12}$$
(35)
Measuring the proton and antiproton cyclotron frequencies, one can get their charge-to-mass ratios. $`r_{p.\overline{p}}`$ gabrielse
$$\left|\frac{r_pr_{\overline{p}}}{r_{\mathrm{av}}}\right|9\times 10^{11}$$
(36)
Analogously, from the charge-to-mass ratio for electron and positron schwinberg
$$\left|\frac{r_e^{}r_{e^+}}{r_{\mathrm{av}}}\right|1.3\times 10^7$$
(37)
As we see the relevant figures of merit are much bigger than the one for the mass difference $`K\overline{K}`$, although, as noticed, these measurements are about six order of magnitude more sensitivity than the one leading to (34). In ref. bluhm2 it has been argued that these figures of merit could not be the relevant ones in testing CPT breaking. In fact, within the approach presented here, at the lowest order in the CPT violating parameters, one has $`g_{}=g_+`$, and similarly the charge-to-mass ratios do not depend on these parameters bluhm2 . To review this point, let me start by the Dirac equation for an electron or a proton including the breaking terms contained in $`_v^{(n=0)}`$ (of course, the breaking parameters may depend on the type of particle one is considering)
$$\left(i\gamma ^\mu D_\mu mb_\mu \gamma _5\gamma ^\mu \frac{1}{2}H_{\mu \nu }\sigma ^{\mu \nu }\right)\psi =0$$
(38)
In a Penning trap the radial confinement is obtained through a strong axial magnetic field, whereas the axial confinement is obtained by a quadrupole electric field. The main corrections due to the CPT and Lorentz breaking parameters are obtained by taking $`A_\mu `$ as the four-potential for a constant magnetic field. Then, to obtain the energy shifts generated by the breaking parameters one makes use of the relativistic Landau levels wave functions and the expressions containing the full QED corrections for the unperturbed levels bluhm2 ; bluhm3 . However, the underlying physics can be understood quite simply recalling the expression for the non-relativistic Landau levels
$$E_{n,\sigma }=\left(n+\frac{1}{2}+\frac{g}{2}\right)\frac{Be}{m},\sigma =\pm \frac{1}{2}$$
(39)
The cyclotron and anomalous frequencies are obtained comparing two Landau levels with different quantum number $`n`$ and with the same and opposite spin configurations respectively
$`\omega _c`$ $`=`$ $`E_{1,1/2}E_{0,1/2}={\displaystyle \frac{Be}{m}}`$
$`\omega _a`$ $`=`$ $`E_{0,+1/2}E_{1,1/2}={\displaystyle \frac{g2}{2}}{\displaystyle \frac{Be}{m}}`$ (40)
The relevant CPT and Lorentz breaking corrections to the energy levels are given by bluhm3
$$\delta E_{n,\pm 1/2}^e^{}=b_3\pm H_{12},\delta E_{n,\pm 1/2}^{e^+}=b_3H_{12}$$
(41)
where we have taken the third axis along the magnetic field of the trap. The frequencies for the antiparticles that we need according to the CPT theorem are the ones with inverted spin, therefore
$$\omega _c^e^{}=\omega _c^{e^+}=\omega _c,Toconcludethisomega_a^e^{}=\omega _a2b_3+2H_{12}$$
(42)
We get
$$\mathrm{\Delta }\omega _c\omega _c^e^{}\omega _c^{e^+}=0,\mathrm{\Delta }\omega _a\omega _a^e^{}\omega _a^{e^+}=4b_3$$
(43)
We recall that these equations hold only at the first order in the breaking parameters and also that the usual relation $`(g2)/2=\omega _a/\omega _c`$ does not hold here since, as noted before, the gyromagnetic ratios do not change at the lowest order.
Since the observables that are measured in a Penning trap are the anomalous and cyclotron frequencies, it seems natural to introduce figures of merit related to these observables. A such figure of merit for CPT violation is bluhm2
$$r_{\omega _a}^e=\frac{|_{n,\sigma }^e^{}_{n,\sigma }^{e^+}|}{_{n,\sigma }^e^{}}=\frac{|\delta E_{n,\sigma }^e^{}\delta E_{n,\sigma }^{e^+}|}{_{n,\sigma }^e^{}}$$
(44)
where $`=E+\delta E`$. For a weak magnetic field one gets
$$r_{\omega _a}^e=\frac{|\mathrm{\Delta }\omega _a|}{2m}=2\frac{|b_3|}{m}$$
(45)
A new analysis of the 1987 experiment by Dehmelt et al. dyck has been done recently in ref. dehmelt obtaining the following bound
$$r_{\omega _a}^e1.2\times 10^{21}$$
(46)
However, the vector $`b_\mu `$ is absolutely constant and as such it rotates with a diurnal period of 23 h and 56 m, when seen in the laboratory frame wich is fixed with respect to the earth. This effect might have given rise to non favorable situations during the observation, and therefore the bound has been a bit relaxed dehmelt
$$r_{\omega _a}^e3\times 10^{21}÷2\times 10^{20}$$
(47)
In the case of proton and atiproton there is no experiment at the moment. Assuming an experimental sensitivity analogous to the electron positron case (meaning $`\delta \omega _a2Hz`$) one gets bluhm3
$$r_{\omega _a}^p=2\frac{|b_3^p|}{m}_p10^{23}$$
(48)
The last case I consider is the spectroscopy of free or magnetically trapped hydrogen ($`H`$) and antihydrogen ($`\overline{H}`$). This is interesting since the two-photon $`1S2S`$ transition has been measured with a precision of $`3.4\times 10^{14}`$ udem in a cold atomic beam of $`H`$ and with a precision of $`10^{12}`$ in trapped $`H`$ cesar . However for the free case the dependence of the $`1S2S`$ transition on the CPT and Lorentz breaking parameters is suppressed by a factor $`\alpha ^2/8\pi `$, since the $`1S`$ and $`2S`$ levels shift by the same amount at the leading order in the breaking bluhm4 . Consider now the spectroscopy of $`H`$ and $`\overline{H}`$ in a magnetic field $`B`$. In the basis $`|m_J,m_I`$ the four $`1S`$ and $`2S`$ hyperfine Zeeman levels are, for $`n=1,2`$
$`|b_n=|1/2,1/2,|d_n=|1/2,1/2`$
$`|a_n=\mathrm{cos}\theta _n|1/2,1/2\mathrm{sin}\theta _n|1/2,1/2`$
$`|c_n=\mathrm{sin}\theta _n|1/2,1/2+\mathrm{cos}\theta _n|1/2,1/2`$ (49)
with $`\mathrm{tan}2\theta _n=(51\mathrm{mT})/n^3B`$. Transitions of the type $`|c_1|c_2`$ have leading-order sensitivity to Lorentz and CPT violation, but they are field-dependent. As a consequence there is a problem connected with the broadening of the lines due to trapping field inhomogeneities.
Consider now hyperfine transitions in the ground state. Again there is the problem of the Zeeman broadening. However one can try to eliminate the frequency dependence on $`B`$ (at lowest order) by choosing a field independent transition point bluhm4 . For $`B0.65T`$ the state $`|c_1`$ is highly polarized ($`|1/2,1/2`$). Then the effect on the transition $`|c_1|d_1`$ of the CPT and Lorentz violating parameters is $`\delta \omega _{cd}^{H,\overline{H}}=2(b_3^p+H_{12}^p)`$. Therefore by putting $`\mathrm{\Delta }\omega _{cd}=\omega _{cd}^H\omega _{cd}^{\overline{H}}`$ the corresponding figure of merit can be defined as
$$r_{cd}^H=\frac{|\mathrm{\Delta }\omega _{cd}|}{m_H}=4\frac{|b_3^p|}{m_H}$$
(50)
Attaining a resolution of $`1mHz`$, one would get bluhm4
$$r_{cd}^H5\times 10^{27}$$
(51)
## Conclusions
In this talk I have reviewed some important consequences of atomic physics measurements in the domain of high-energy physics. In particular APV in cesium could be the first real indication of new physics beyond the SM. The atomic physics tests of the CPT symmetry are already at a spectacular level of sensitivity, and the future experiments on $`H`$ and $`\overline{H}`$ could give bounds well below the one expected from string theory.
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# Indications of Spin-Charge Separation at Short Distance and Stripe Formation in the Extended t-J Model on Ladders and Planes
## I Introduction
The study of high temperature superconductors continues attracting the attention of the Condensed Matter community. In recent years, much effort has been devoted to the understanding of the spin incommensurability that appears in neutron scattering experiments for some of these compounds. Tranquada’s interpretation of the experiments is based on “stripes” where charge is confined to one-dimensional (1D) paths in the crystal, with an average hole charge n<sub>h</sub>=0.5, namely one hole for every two sites along the stripe. The stripe interpretation appears robust in the one-layer material $`\mathrm{La}_{2\mathrm{x}}\mathrm{Sr}_\mathrm{x}\mathrm{CuO}_4`$, although it is still controversial in the bilayer $`\mathrm{YBa}_2\mathrm{Cu}_3\mathrm{O}_{6+\delta }`$. Experimental results compatible with $`metallic`$ stripes have also been reported using other techniques. On the theory front, early studies discussed the presence of stripes in Hubbard and t-J models, although typically with density n<sub>h</sub>=1.0. Recently, the existence of n$`{}_{h}{}^{}<`$ 1 stripes was discussed using computational techniques directly in the t-J model at intermediate coupling J/t, without the need of long-range Coulomb interactions. Sometimes these stripes are described as a condensation of d-wave pairs. Stripes appear also in Monte Carlo studies of the spin-fermion model. However, the origin of such complex spin-charge arrangement is still under much discussion, and even its stabilization remains controversial in the standard t-J model.
In parallel with the developments in neutron scattering experiments, in recent years photoemission techniques have provided a plethora of information about the cuprates. In particular, the experimental study of the one-hole spectral function of the parent insulator compound gave us useful information to judge the quality of existing models for these materials. Based on the experimental one-hole dispersion it has been convincingly shown that the t-J model is $`not`$ enough to address the cuprates but extra terms must be added in order to reproduce the photoemission results for the insulator. These terms appear in the form of extra hole hoppings, regulated in intensity by hopping amplitudes usually denoted by t and t<sup>′′</sup>. The bare value of these extra amplitudes is small compared with the nearest-neighbor hopping t, but its influence is substantial since they link sites belonging to the same sublattice and they are not so heavily renormalized to smaller values as it occurs with t. The origin of these extra hoppings was discussed before and they appear naturally in mappings from the three-band model for cuprates to a one-band Hamiltonian, and also in band-structure calculations. In fact, what is unnatural is to assume that the standard t-J model, with t=t<sup>′′</sup>=0.0, is valid for the cuprates without corrections. Unlike in gauge theories, there are no renormalizability or symmetry arguments at work preventing the introduction of nonzero extra hoppings in models for copper oxides. For these reasons it is important to study the extended t-J model in detail and, in particular, whether stripes occur in its ground state.
Recent studies of the extended t-J model have provided useful and sometimes surprising information. For instance, it was observed that the one-hole QP weight Z is considerably reduced by the addition of t and t<sup>′′</sup> at fixed J/t. This result is particularly dramatic at momenta $`(0,\pi )`$ and $`(\pi ,0)`$, and its origin lies in the generation of across-the-hole antiferromagnetic (AF) correlations, namely in the reference frame of the mobile hole the two spins immediately next to it along one axis are aligned $`antiferromagnetically`$. This tendency is the opposite as expected from a vacancy in an antiferromagnet where those two spins should be ferromagnetically aligned, since they are in the same sublattice. Such a curious result was also observed in previous DMRG studies of ladders and planes using the standard t-J model by White and Scalapino. Its importance and origin was recently discussed by Martins et al. where its presence was conjectured to be caused by tendencies in the system to spin-charge separation at short distances. In other words, the kinetic energy of the hole with the extra mobility induced by t and t<sup>′′</sup> dominates most of the physics of the one-hole problem and it tries to generate in its vicinity an environment that allows for the hole to move without fighting against the spin background. In one-dimension this is easy to set up, and indeed an antiferromagnetic across-the-hole environment is stabilized on chains leading to spin-charge separation. In two-dimensions the hole tries to create such an environment in its vicinity but it cannot produce a total spin-charge separation since frustration is induced when antiferromagnetic across-the-hole interactions are generated along both axes. It is expected that for realistic values of the couplings, the spin-charge separation will only be local (i.e. at short distances). This result, obtained using small square clusters, is clear in the presence of nonzero t and t<sup>′′</sup> of the proper sign and magnitude, but it also occurs in the standard t-J model in the regime of very small J/t, which remains mostly unexplored. This interesting observation redefines the relevant value of the parameter J/t from the number usually used in most studies (0.3-0.4) to smaller values close to 0.1.
Also motivated by the results in Ref., recently the analysis of many holes at couplings that are expected to lead to robust across-the-hole correlations was reported by Martins et al.. Indications of stripes were obtained in this case, similarly as those previously found in Ref., rationalized as the natural tendency of locally spin-charge separated states to avoid the spin frustration caused by the nontrivial spin environment that each hole generates. The hole density in these stripes is in good agreement with experiments, and qualitatively the results resemble those in the “holons in a row” picture of Zaanen. Tight hole bound-states are not needed to generate stripes.
These observations suggest that either the extended t-J model or the standard t-J model at very small coupling must be analyzed in detail if the goal is to reproduce the physics of the striped regime of the cuprates. In addition, it is important to confirm the generality of previous results by studying other systems with similar physics but where computational studies can be carried out on larger clusters. For this purpose, here $`ladders`$ are investigated using the extended and standard t-J models in the coupling regime where in Ref. the robust across-the-hole feature was observed on small square clusters. Our main result is that these robust correlations are also very clearly identified on ladders, even with only two legs, and, thus, its presence is more general than naively anticipated. Based on our results it is clear that holes on ladders tend to form one-dimensional (1D) like environments in their vicinity to improve their kinetic energy leading to an effective reduction in dimensionality. Spin incommensurability is generated by this procedure, a novel result in two-leg ladder systems to the best of our knowledge, although its presence in two dimensions has been known for some time.
The present paper has a second goal which is the detailed analysis of stripe formation upon doping of ladders with many holes. For this purpose here the recent results of Ref. where stripes were observed are extended to couplings and parameters not reported before, to illustrate and confirm the tendencies toward stripe formation in this context. The proposed rationalization for this stripe stabilization was already discussed and it is based on the sharing of the locally spin-charge separated spin environment created by each individual hole, thus avoiding spin frustrating effects.
The model used here is the extended t-J model defined as
$$\mathrm{H}=\mathrm{J}\underset{\mathrm{𝐢𝐣}}{}(𝐒_𝐢𝐒_𝐣\frac{1}{4}\mathrm{n}_𝐢\mathrm{n}_𝐣)\underset{\mathrm{𝐢𝐦}}{}\mathrm{t}_{\mathrm{𝐢𝐦}}(\mathrm{c}_𝐢^{}\mathrm{c}_𝐦+\mathrm{h}.\mathrm{c}.),$$
$`(1)`$
where $`\mathrm{t}_{\mathrm{𝐢𝐦}}`$ is t for nearest-neighbors (NN), t for next NN, and t<sup>′′</sup> for next next NN sites, and zero otherwise. The scale of energy will be t=1 unless otherwise specified. The rest of the notation is standard. Comparison with PES experiments showed that t=-0.35 and t<sup>′′</sup>=0.25 are relevant to explain PES data. To simplify our studies the ratio t/t<sup>′′</sup> will be fixed to -1.4 in most of our analysis, although other ratios will be used in some cases. As numerical techniques, the Density Matrix Renormalization Group (DMRG), Exact Diagonalization, and an algorithm using a small fraction of the ladder rung-basis (optimized reduced-basis approximation, or ORBA) are here used. The paper is organized by the number of legs of the ladders considered, from two to six. Analytical results are discussed after the numerical methods, illustrating the appearance of tendencies toward across-the-hole antiferromagnetism in two dimensional systems. In the last section it is concluded that the mostly unexplored extended t-J model contains interesting physics related to the cuprates, that deserves further work.
## II Two-Leg Ladders:
### A One Hole:
The present computational analysis of ladder systems starts with the two-leg ladder case. Previous literature has shown that in the absence of t and t<sup>′′</sup> hoppings, and at intermediate values of J, such as 0.4, one hole behaves as expected from a carrier in an antiferromagnet even if the ladder spin background has only short-range magnetic order. In other words, a hole creates a spin distortion in its vicinity (spin polaron) and the one-hole ground state has a finite QP weight Z. Two of these spin polarons bind into a hole pair that leads to superconducting correlations upon further doping. However, in this paper our goal is to analyze a region of parameter space not studied before, to the best of our knowledge, where the hole is expected to be substantially more mobile. For this purpose J is reduced to values between 0.1 and 0.2, still within the domain considered realistic in studies of the cuprates, and t, t<sup>′′</sup> are made nonzero and of the sign and magnitude as suggested by photoemission experiments. At such “small” values of J/t, the spin background no longer is expected to fully control the behavior of the hole but the spins have to arrange in such a way that the hole kinetic energy is optimized (but still without leading to a ferromagnetic state that may be the best configuration at J/t$``$0.0).
Consider first Exact Diagonalization results in the one-hole sector. In Fig.1a, the spin-spin correlations are shown for the case where the hole is projected on the (arbitrary) site indicated, from the full one-hole lowest-energy state with momentum $`𝐤`$=$`(\pi ,0)`$, working at J=0.2, t=-0.35 and t<sup>′′</sup>=0.25. It can be observed that at a distance of about three lattice spacings from the hole the spin correlations are similar to the ones expected for an undoped ladder, with robust antiferromagnetism along both the rungs and legs. However, in the vicinity of the hole the spin correlations are substantially altered. In this case the spins belonging to the same rungs are no longer strongly antiferromagnetically correlated and the system in the vicinity of the hole resembles a pair of weakly-coupled portions of a chain, one per leg. In this respect it appears that an “effective” transition from a two-leg ladder to 1D chains has occurred near the hole, a surprising result. Holes seem to optimize their energy by creating 1D environments in its vicinity. Note the presence of robust across-the-hole spin correlation on the upper leg of Fig.1a, which also appears in the exact solution of the 1D Hubbard model due to spin-charge separation. Fig.1a suggests that the spin associated to the hole is spread uniformly in the 1D-like segment dynamically generated, and in this respect spin-charge separation occurs locally, as anticipated in the Introduction and as it was reported in Ref. using small square clusters. Note in Fig.1a that the spin belonging to the same rung as the hole appears as “free” and there is a strong AF bond at distance of two lattice spacings across it, as it occurs across-the-hole.
The comparison of Figs.1a and 1b, the latter obtained without extra hoppings, shows that t and t<sup>′′</sup> are important at J=0.2 to regulate the size of the 1D environment around the hole: as t, t<sup>′′</sup> grow in magnitude, the size of the 1D-like region also grows. This same effect is produced reducing J at fixed t,t<sup>′′</sup> as shown in Fig.1c: here the distortion is larger than at J=0.2 and even weak ferromagnetic links are generated between the chains, an unexpected result quite different from the well-established physics of undoped ladders at intermediate J. The results here are also weakly dependent on the momentum: in Fig.1d, the case of $`𝐤`$=$`(\pi /2,0)`$ (momentum of the overall ground-state of one hole) presents a distortion similar in size as at $`𝐤`$=$`(\pi ,0)`$, although the AF long-bond in the leg opposite to where the hole is projected (Fig.1a) is no longer present. However, the long-bond across- the-hole is still robust. Fig.1e is another case illustrating the momentum dependence in the problem, corresponding this time to $`𝐤`$=$`(\pi ,\pi )`$. Overall, it can be concluded that the lowest one-hole energy states for all momenta studied here, and in large regions of parameter space, present very similar spin arrangements. The across-the-hole feature emphasized in Ref. is very robust in all cases and in this sense there is a clear $`\pi `$-shift across- the-hole even in the simple case of a two-leg ladder, similarly as it occurs in 2D systems. The dynamically induced 1D-like regions near the carrier are also clear in our studies.
The generation of a quasi-1D structure in the vicinity of the hole and the expected spread of the hole’s spin-1/2 over several lattice spacings (spin-charge separation at short distances) causes a drastic reduction in the QP weight Z. In Fig.2, the one-hole spectral function A($`𝐤`$,$`\omega `$) is presented at three different momenta. The results are shown as it is usual in experimental photoemission literature with the lower energy states appearing near $`\omega `$=0, reference energy which is located at an arbitrary position in this half-filled case, and the rest of the states running to the left. Data for J=0.2 and several values of t, at fixed t/t<sup>′′</sup>=-1.4, are presented, together with the weight Z and position of the first state (arrows). Fig.2a shows that Z at this value of J and $`𝐤`$=$`(\pi ,0)`$ is already small even in the absence of extra hoppings, result similar to those observed in previous investigations of the t-J model on square lattices. However, it is remarkable the rapid reduction of Z with increasing $`|`$t$`|`$, to values that in practice are basically negligible by the time these extra hoppings reach its realistic values. This is correlated with the appearance of the complex quasi-1D structure discussed in Fig.1 A similar phenomenon but even more pronounced exists for $`𝐤`$=$`(\pi ,\pi )`$ (Fig.2c), where Z is negligible even without extra hoppings. On the other hand, the results at $`𝐤`$=$`(\pi /2,0)`$ (Fig.2b), while still corresponding to small values of Z of the order of 10% of the maximum, do not show such a dramatic reduction with increasing t. In this case the spin appears still located in the vicinity of the hole, although it is spread over several sites. Z becomes negligible at this momentum only at values of $`|`$t$`|`$ larger than shown in the figure.
The results for the QP weight Z are more explicitly illustrated in Fig.3a where its $`|`$t$`|`$ dependence is shown for several momenta at fixed J=0.2. Clearly for all momenta the presence of t eventually leads to a drastic reduction of the Z value, although for $`𝐤`$=$`(\pi ,0)`$ it happens rapidly with increasing t while at $`𝐤`$=$`(\pi /2,0)`$ it needs a t larger than the nearest-neighbors hopping t. In Fig.3b the lines of constant weight Z corresponding to $`𝐤`$=$`(\pi ,0)`$ are shown. It is clear that Z at this momentum rapidly reduces its value, both at fixed J increasing $`|`$t$`|`$ or for fixed hoppings decreasing J. Although a careful finite size scaling is difficult, previous experience with two-dimensional clusters suggest that Z tends to decrease as the lattice size grows. As a consequence it is expected that the values reported in Figs.2 and 3 will actually be even smaller for a very long two-leg ladder. Note also that the lines of constant Z suggest that the physics of, say, intermediate J and finite t may be smoothly connected to that of small J and zero t, quite similarly as it happens on square clusters. Then, it may occur that analyzing the regime of abnormally small J ($``$0.1) of the t-J model may effectively account for the presence of extra hoppings, as was conjectured in Ref.. This assumption is interesting due to the anticipated unusual properties of the small J/t region of the t-J model, and it avoids the somewhat aesthetically unpleasant use of extra hoppings in the present theoretical studies.
### B Several holes:
The previous subsection showed that the two-leg ladder system presents exotic behavior in the one-hole sector when the regime of small J is investigated or alternatively when at a fixed J the hoppings t-t<sup>′′</sup> are switched on. It is important to explore what occurs for more holes in the problem, and Fig.4 contains exact results for the case of two holes on an 2$`\times `$8 ladder. Data is shown for the case when the two holes are projected to its most probable location in the ground-state, which in the cases studied coincide with the maximum distance among them allowed in the cluster. The spin structures in Fig.4 near the holes clearly resemble those found around individual holes in the previous subsection. In particular the across-the-hole AF correlation is very prominent in all cases shown in Fig.4. This structure naturally leads to spin incommensurability, as observed in the spin structure factor presented in Fig.5a. It is interesting to observe that the coupling between chains can be weakly $`ferromagnetic`$ in some cases (Fig.4a), weakly AF in others (Fig.4b), or even virtually negligible (Fig.4c) depending on the actual couplings used. Overall it is clear that the legs are approximately decorrelated upon hole doping at the couplings analyzed here, and each leg behaves as a 1D-chain described by the same original t-J model. Each individual leg resembles the expected result for the 1D Hubbard model at large U, as discussed before for 2D clusters. It is remarkable that very recently nuclear spin relaxation results for two-leg ladder materials have also reported the decoupling of the two legs at high temperature. The analysis of the relation between our results and those of Ref. deserves further work.
In the two-hole ground-state the expectation value of the number operator for a particular momentum is shown in Fig.5b. Clearly the momenta $`\mathrm{k}_\mathrm{x}`$ dominating in this state is not only $`\pi /2`$, the momentum of the lowest energy state of the one-hole sector, but actually $`3\pi /4`$ and $`\pi `$ are also important as well as the $`\mathrm{k}_\mathrm{y}`$=$`\pi `$ sector. Then, there is no indication of “hole-pockets” in the doped system, and momenta over a wide range contribute to the two-hole ground-state. This suggests that the one-hole states with the largest weight in the two-hole ground-state have the across-the-hole structure and local spin-charge separation.
The charge structure factor was also calculated in these investigations. In Figs.6a-b exact results on a 2$`\times `$8 cluster and DMRG results on 2$`\times `$16 and 2$`\times `$32 clusters, all at hole density x=0.125, are shown. They are very similar, illustrating the lack of strong size effects systematically observed in our studies of two-leg ladders. The results obtained in the more traditionally studied case of J=0.4, t=t<sup>′′</sup>=0.0 are also shown in Figs.6c-d. They are qualitatively different from those in Figs.6a-b, difference likely related with the formation of the quasi-1D regions mentioned before.
## III Three-Leg Ladders:
### A One hole:
Results for three-leg ladders have also been gathered in this work. In Fig.7 the hole is projected at the location shown, from the lowest-energy state at the momenta indicated in the caption. In Fig.7a results are presented at J=0.1 and nonzero t-t<sup>′′</sup>. The formation of the across-the-hole structure is very clear, in both directions. Along the vertical one, a strong AF bond is formed between the two spins of the rung where the hole is located. Its strength is close to that of a perfect singlet and, thus, the coupling of those spins with the rest is small. This structure is present at all the couplings investigated as shown in Figs.7a-c, even intermediate and large J, and it is also clearly present in the exact solution of a three-site t-J model at many couplings. In the small J/t regime of the standard t-J model, and in the extended t-J model as well, this strong AF bond denotes a precursor of the $`\pi `$-shift across the stripe that will form (along the PBC direction) as the hole-density increases. Regarding the horizontal AF bond (along the legs), its strength depends on the particular value of J and t-t<sup>′′</sup>. In the intermediate coupling region often studied, exemplified by J=0.5 and no extra hoppings as in Fig.7c, there is no across-the-hole AF correlation along the legs. On the other hand, for the couplings of Fig.7a, the clear 1D-like AF segment along the central leg induced in the vicinity of the hole suggests a larger mobility of the carrier and generation of 1D-like segments, as discussed in the Introduction and for two-leg ladders.
The tendency toward local spin-charge separation can be studied further if the mean-value of the spin in the z-direction is calculated when the hole is projected from its ground-state to a given site, as done before in Figs.7a-c. Since removing one spin from a cluster with an even number of sites creates states with spin-1/2, together with the hole in Fig.8 there must be a spin-1/2 spread over the cluster. The results are shown in Fig.8a for the case of small J and nonzero t-t<sup>′′</sup>. In this situation it is clear that the z-component of the spin 1/2 is distributed approximately uniformly along the central leg where the hole is projected (the spin in the outer legs is negligible). In agreement with the introductory discussion, the mobile hole creates a 1D environment (along the central leg in this case) to help in its propagation, and in this context the spin and charge are separated. On the other hand, the results corresponding to an intermediate J of value 0.5 and no extra hoppings corresponds to a staggered spin background surrounding the hole (Fig.8b), similar to results for a vacancy in a Néel background, with the exception of the two spins in the same rung as the hole, which try to form a strong singlet and thus its mean z-component spin tends to vanish. There is a clear qualitative difference between Figs.8a and b, caused by the high mobility of the hole in the presence of a small J and nonzero hoppings beyond nearest-neighbors. The QP weights Z (not shown) associated to the structure of Fig.8a tend to vanish at $`𝐤`$=$`(0,\pi )`$ and others, as found before for two-leg clusters and small two-dimensional systems.
### B Many holes:
When two holes are considered on the three-leg ladder the situation is similar as in the case already described of two-legs, namely each individual hole carries a spin arrangement in its vicinity similar to that found in the one hole case. Results are shown in Fig.9a at small J and Fig.9b at intermediate J, in both cases without t-t<sup>′′</sup>. In the first case, AF bonds across-the-hole are found in both directions, while in the second the bond along the legs turns ferromagnetic as for vacancies in a Néel background and as in Fig.7c for one mobile hole.
If three holes are considered at small J (Fig.9c) the most likely hole configuration corresponds to a n<sub>h</sub>=0.5 stripe along the central leg, with clear and robust AF bonds across it joining the two outer legs (in spite of having open boundary conditions along the rungs). Within the central leg or stripe, AF bonds are formed across-the-hole, as in a 1D system with the same Hamiltonian. Fig.9c qualitatively indicates the dynamical separation of a chain subsystem carrying the charge from the rest of the spins which are correlated as if that chain would not exist. As emphasized in other parts of this paper and in Refs., this appears to be the way in which the system achieves a partial separation of spin and charge in two dimensions.
## IV Four- and Six-leg Ladders:
### A One Hole:
Results similar to those found in the case of two- and three-leg ladders are also observed with four legs. Results for the case of one hole with momentum $`𝐤`$=$`(\pi ,0)`$ are shown in Fig.10 on a 4$`\times `$6 cluster solved exactly using the Lanczos method. In Fig.10a results for small J and nonzero t-t<sup>′′</sup> are presented: here the AF bonds across-the-hole can be observed and 1D-like segments are created near the hole in both directions. The effect is amplified if at fixed J the values of the hoppings t-t<sup>′′</sup> are increased. Fig.10b contains the results using abnormally large extra hoppings: now the AF bonds across-the-hole are quite robust, with a strength that grows as the extra hoppings grow in magnitude. There is a smooth connection between the physically acceptable values of t-t<sup>′′</sup> and those used in Fig.10b to amplify the effect. Such a connection suggests a common origin to the structure.
The one-hole spectral function $`A(𝐤,\omega )`$ is shown in Figs.11a-b for the momenta indicated and as a function of $`|`$t$`|`$. At $`(0,\pi )`$ the QP weight Z rapidly decreases with an increasing t hopping amplitude, while at a momenta closer to the expected ground-state momentum of one hole it remains more robust, but it eventually tends to vanish at large $`|`$t$`|`$ (see also Fig.12). These results are very similar to those observed for two-leg ladders, for planes, and also for three-leg ladders (not shown).
### B Many Holes:
To clarify the behavior of two holes on four-leg ladders, consider the use of cylindrical boundary conditions (OBC in one direction and PBC in the other). Suppose the two-hole ground-state is considered and one hole is projected to a site belonging to the central legs, where the density of holes is the largest. In this situation it is possible to study the probability of finding the other hole, with results shown in Fig.13a: the second hole clearly prefers to be along the PBC direction and still within the two central legs. Actually the largest probability is found at a distance of two lattice spacings along the same leg. The results can be interpreted as the formation of a loop of charge which wraps around the direction with PBC, as discussed recently in Ref. and also in Ref.. The state appears to correspond to a n<sub>h</sub>=0.5 stripe, not rigid but fluctuating in the direction perpendicular to it.
In Fig.13b the spin correlations are shown for the holes projected into the configuration with the largest weight in the two-hole ground-state. In agreement with the discussion associated with Fig.13a, the holes appear to be forming a n<sub>h</sub>=0.5 stripe along the direction with PBC. Moreover, the spin correlations across-the-stripe are clearly antiferromagnetic as in the three-leg ladder (Fig.9c) and in experiments, and the two spins of the four-site stripe considered here are correlated also antiferromagnetically. This result is representative of a large number of similar data gathered in our present analysis, namely it appears that in general holes tend to form loops of charge around the closed direction if a system has cylindrical boundary conditions. This is also in excellent agreement with the conclusions of earlier work by White and Scalapino, although their description of stripes is sometimes based on the condensation of d-wave pairs while ours is based on a kinetic energy optimization (namely a one-hole problem).
The appearance of stripes based on the 4$`\times `$4 results of Fig.13 needs to be confirmed increasing the cluster size. This analysis was done in part in Ref. as discussed in the Introduction, but here those results are expanded and more details are provided. Note that our present analysis using the DMRG method is restricted to the t-J-t model, namely the t<sup>′′</sup> hopping will be considered to be zero. The reason is that in the implementation of the DMRG technique sites are aligned along a one-dimensional pattern even for ladder geometries, and a t<sup>′′</sup> hopping would link sites along this equivalent chain which are located several lattice spacings from each other, reducing the accuracy of the method. Nevertheless even without t<sup>′′</sup> the stripes reported in Fig.13 and Ref. can be clearly observed on larger systems. For example, consider in Fig.14 two holes on a 4$`\times `$12 cluster with cylindrical boundary conditions (PBC along the rungs), J=0.5, and t=-0.3 (t<sup>′′</sup>=0.0). One of the holes is projected to an arbitrary site of one of the central rungs, where the hole density is the largest. The distribution of the second hole around the projected one is quite similar to the result observed on the 4$`\times `$4 cluster, namely the largest chances are at distance of two lattice spacings along the rung. This state resembles a bound-state in the sense that the two holes are close to each other, but its shape is better described as a stripe configuration or a loop of charge that wraps around the short direction.
If more holes are added to the four-leg ladder, it appears that stripes similar to those observed in Fig.14 are formed. For example, consider the case of four holes on a 4$`\times `$8 cluster at small J and without t and t<sup>′′</sup>. In Fig.15 the density of holes is shown for the case where one hole is projected into one of the rungs with the largest density. It is clear from the figure that in the vicinity of the projected hole there is a structure similar to that observed for the two-hole ground state, namely a local maximum in the hole density is observed at two lattice spacings from the projected hole. In addition, clearly a large accumulation of charge appears on the other side of the cluster from where the projected hole is located. This other sector populated by holes involves two rungs in width, and it is framed in Fig.15. Note that the density in that second stripe is very uniform showing that there are no charge correlations between the two stripes (the two charged loops are independent of each other).
In Fig.16a the charge structure factor is shown for the case of an 4$`\times `$8 cluster with four holes, studied with the DMRG method at small J. The result resembles those obtained for two-leg ladders in Fig.6. In Fig.16b the spin structure factor is shown for the same cluster and density. Spin incommensurability clearly appears in this system, as remarked before in Ref. working at other couplings, and as presented for the two-leg ladder (Fig.5a) as well. It is clear that the two-, three- and four-leg ladder systems share very similar physics in the charge and spin sectors.
### C Six-leg Ladders:
Accurate results for six-leg ladders are difficult to obtain in the regime of parameters studied in this paper, namely small J and nonzero t-t<sup>′′</sup>. However, some nontrivial results can still be gathered. For instance, consider in Fig.17 the case of three-holes on a 6$`\times `$4 cluster with PBC along the direction with six sites, and OBC along the short one. In this case a good approximation to the ground-state can be obtained using the ORBA method with about one million 4-site rung-basis states. The procedure to obtain the results shown in Fig.17 are the following: first it was noticed that the central long legs are the ones the most populated by holes. Then, one hole was projected at an arbitrary site belonging to those central legs. In the next step, the hole density is obtained with that hole projected, and at the position with the largest density a second hole is now projected. The distance between the first and second projected holes is two lattice spacings. With those holes projected as shown in Fig.17, the hole density is recalculated. It is clear that at two lattice spacings from the projected holes the density has a maximum. This result is once again compatible with a n<sub>h</sub>=0.5 stripe formed this time by three-holes in a closed loop around the direction with PBC. The stripe is not rigid but fluctuating perpendicular to its main direction.
The density of holes with a given momentum $`𝐤`$ for the 6$`\times `$4 cluster with three holes is shown in Fig.18. Clearly the momenta the most important are those around $`(\pi ,\pi )`$, result similar to those reported for the two-legs cluster in Fig.5b and for a 4$`\times `$6 two-holes cluster with PBC in both directions in Ref..
## V Analytical calculations: spin-polaron approach
In previous sections, it has been shown among other items that the across-the-hole correlations identified before numerically on small square clusters for the one-hole problem exist also on ladders. For completeness, in this section an analytical approach to the one-hole problem in two dimensions is described. In agreement with the above mentioned results, once again the existence of across-the-hole correlations is confirmed this time using a non-numerical method, highlighting the robustness of this feature.
The analytic approach used here is based on the picture of the spin-bag QP or magnetic polaron. The spin deviations from a perfect antiferromagnet in the vicinity of the hole are described by a set of fluctuation operators $`A_n`$ , which create strings of spin defects attached to the hole. The one-hole Green’s function is evaluated using a cumulant version of Mori-Zwanzig projection technique.
The undoped ground state is modeled by $`|\psi =\mathrm{exp}(_\nu \alpha _\nu S_\nu )|\varphi _{\mathrm{N}\stackrel{´}{\mathrm{e}}\mathrm{el}}`$ where the operators $`S_\nu `$ create clusters of spin fluctuations in the classical Néel state. Following Ref. it is possible to obtain a non-linear set of equations for the coefficients $`\alpha _\nu `$, which is solved self-consistently. The fluctuation operators $`A_n`$ for describing the spin deviations caused by hole motion include the usual string operators up to a certain length $`l_{\mathrm{max}}`$ and additional operators for local fluctuation configurations. The Green’s function is calculated by projection technique in the subspace spanned by the operators $`A_n`$, which amounts to the diagonalization of the dynamical matrix $`\psi |A_m^{}HA_n|\psi `$. Processes outside the subspace formed by $`A_n|\psi `$ (i.e., self-energies) are neglected, therefore a discrete set of poles for each spectrum is obtained. The eigenvectors of the dynamical matrix represent the wavefunctions corresponding to the poles and they allow for the calculation of expectation values such as spin correlation functions. In the following the pole next to the Fermi surface representing the QP is the only one considered, the remaining part of the spectrum forms an incoherent background well separated from the QP peak. The approach sketched here has been successfully applied to the one-hole problem in several contexts ; details are described in Ref. . In the present calculations strings with $`l_{\mathrm{max}}=5`$ have been used (total number of fluctuation operators $`A_n`$ around 1200). Note that the accuracy of the approximation decreases with increasing $`t`$ since the employed expansion parameter is basically the amplitude of the deviations from the undoped state. Therefore it is better to restrict ourselves to the parameter range with $`\mathrm{J}/\mathrm{t}`$$`>`$0.05 where the comparison with numerical results reveals the good accuracy of the analytical approach.
First, let us discuss the weight of the QP pole which is known to decrease with decreasing J/t since the size of the spin polaron increases. The full parameter dependence of the QP weight Z at momentum $`(\pi ,0)`$ is shown in Fig.19a. The introduction of t$`<`$0 rapidly suppresses the value of $`\mathrm{Z}_{\pi ,0}`$ even for intermediate or large J, such that the entire region of moderate t is characterized by small $`\mathrm{Z}_{\pi ,0}`$, in good agreement with the computational results shown in this paper and in Ref.. On the other hand, the t dependence of Z at $`(\pi /2,\pi /2)`$ is weaker (not shown), also as found in Ref.. Nevertheless, its associated QP weight Z is also strongly suppressed with the reduction of J.
In order to study the spin configuration near the mobile hole, static spin correlations in real space have been calculated relative to the hole position. It is known that the limit of small t (static vacancy) leads to an increase of the antiferromagnetic correlations on bonds near the hole, whereas hole hopping tends to “scramble” the spins in the hole environment which in general weakens the AF tendencies. This general behavior is well reproduced by our calculations. However, the introduction of t leads additionally to antiferromagnetic correlations $`S_{AH}`$ across-the-hole, see Figs.19b and 20b, quite consistent with those found numerically here and in previous studies. In Fig.20b, the strength of the antiferromagnetic bonds at J=0.05, t=-0.35, t<sup>′′</sup>=0.25, and momentum $`(\pi ,0)`$ are shown as illustration. Across the mobile hole strong antiferromagnetic correlations develop which are supported by further strong AF bonds forming a chain segment. The other nearest-neighbor bonds are weaker by a factor of two than the chain bonds. The values of $`\mathrm{S}_{\mathrm{AH}}`$, obtained with the analytic approximation used in this section are shown in Fig.19b. Although they are not as strong as found numerically, the qualitative trends agree quite well with computational studies.
The stronger tendency for local spin-charge separation, at small J/t or in the presence of t and t<sup>′′</sup>, can be understood within the string picture. The string of defect spins attached to the hole basically connects the spin and charge parts of the excitation, confining them at long distances. Therefore longer average strings correspond to less tightly bound spin and charge, i. e., larger spin polarons, immediately implying a smaller quasiparticle weight. The energy cost per unit length of string is proportional to J, hence smaller J/t allows for longer strings. Interestingly, the presence of t and t<sup>′′</sup> leads to a related effect, namely the hole can easily ‘get rid’ of its attached string by intra-sublattice hopping, the string being then absorbed in the background spin fluctuations. This mechanism makes long strings effectively less costly, leading to the described large polarons with small QP weight Z.
For completeness, in Fig.21 the calculated QP dispersions along high-symmetry lines in the Brillouin zone are shown. The points correspond to (a) J=0.40, t=0.0, (b) J=0.40, t=-0.35 and (c) J=0.05, t=-0.35 with t/t<sup>′′</sup>=-1.4. The change in the curvature near $`(\pm \pi /2,\pm \pi /2)`$ is clearly visible when adding t and decreasing J. Whereas the dispersion is strongly anisotropic for large J and small t, there is an entire region of near isotropy for moderate values of t and a large range of J/t, in agreement with PES experiments. This is also illustrated in Fig.20a where the ratio of the effective QP masses at $`𝐊_0`$=$`(\pi /2,\pi /2)`$ for the entire parameter space are presented. The masses are defined as usual as eigenvalues of the tensor $`𝐦`$ given by $`ϵ_𝐤=(2𝐦)_{\alpha \beta }^1(𝐤𝐊_\mathrm{𝟎})_\alpha (𝐤𝐊_\mathrm{𝟎})_\beta `$ where $`ϵ_𝐤`$ is the QP energy near $`𝐊_0`$.
Summarizing, the introduction of t leads to a suppression of the QP weight, especially at momentum $`(\pi ,0)`$, a delocalization of the spin carried by the spin polaron, and the formation of antiferromagnetic correlations across the mobile hole. All the results are in good qualitative agreement with the numerical calculations reported here.
## VI Conclusions
In this paper, tendencies toward spin-charge separation at short distances have been discussed in the context of the extended t-J model using ladder geometries. This effort generalized previous calculations carried out on small square clusters. Analytic approximations have also been used, with results in good agreement with the computational ones. Overall it is concluded that in regimes where the hole kinetic part of the Hamiltonian dominates, holes tend to arrange the spin environment in such a way that across-the-hole robust antiferromagnetic correlations are generated, both on ladders and planes. This arrangement helps the hole move easily among the spins, and it resembles the structure found in one-dimensional spin-charge separated systems. For a variety of reasons described here, it is believed that at least at short distances similar tendencies toward spin-charge separation are at work in ladders and two-dimensional systems, at very small J in the standard t-J model, or in the extended t-J model. At finite hole density, holes share their nontrivial spin environment, forming half-doped stripes as recently discussed by some of the authors in Ref.. Here more evidence substantiating this previous result has been provided. The stabilization of stripe tendencies discussed here is based upon a $`small`$ J/t picture, and it has no obvious relation with those emerging in the opposite limit of large J/t based upon the frustration of phase separated tendencies. Although more work is certainly still needed to confirm the picture described here, it appears that a gas of spinons and holons, as envisioned in two-dimensional theories for cuprates based upon spin-charge separation, may not be operative at low hole densities, but instead stripes of holons seem to be forming.
## VII Acknowledgments
E. D. thanks NSF (DMR-9814350) and the Center for Materials Research and Technology (MARTECH) for support, J. C. X. thanks FAPESP-Brazil for support and C. G. thanks Fundación Antorchas for partial support. M. V. acknowledges support by the DFG (VO 794/1-1) and by US NSF Grant DMR 96-23181.
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# Experimental realization of a highly structured search algorithm
## Abstract
The highly structured search algorithm proposed by Hogg\[Phys.Rev.Lett. 80,2473(1998)\] is implemented experimentally for the 1-SAT problem in a single search step by using nuclear magnetic resonance technique with two-qubit sample. It is the first demonstration of the Hogg’s algorithm, and can be readily extended to solving 1-SAT problem for more qubits in one step if the appropriate samples possessing more qubits are experimentally feasible.
Quantum computers\[1-5\] can outperform classical ones owing to utilization of quantum mechanical effects. The potential for quantum parallel computing mainly resulted from superpositions and entanglements of quantum states has made quantum computers capable of solving classically intractable problems (such as factoring large integers) and finding tractable solutions more rapidly(e. g., searching an unsorted database). However, carrying out an actual quantum computing is much difficult although the rudiment of the quantum computing is easily understood and several quantum algorithm have been discovered. Among the approaches meeting the requirement for performing quantum computing, the nuclear magnetic resonance(NMR) technique is considered to be one of current promising method, with which the Deutsch-Jozsa algorithm and Grover’s search one has been experimentally implemented with two-qubit samples\[10-13\]. Recently, it was reported that an algorithmic benchmark has been experimentally realized using NMR with seven qubits. It was easily found from those works that NMR can act as a considerably satisfactory small-scaled quantum computer to test the simple cases of the various quantum algorithm. In this contribution, also using NMR scheme with two-qubit sample, we report the first experimental demonstration of another quantum search algorithm. This algorithm, called a highly structured quantum search algorithm, was claimed to be much more efficient than Grover’s algorithm that ignores the problem structure, and can be exploited in many practical examples instead of for academic study alone as an artificial problem.
The Hogg’s algorithm is associated with the combinatorial searching known as nondeterministic polynomial-time problems(NP), an important class of intractable problems. Its prototype is the satisfiability problems(SAT). A SAT consists of a logic formula in $`n`$ boolean variables, V<sub>1</sub>,…,V<sub>n</sub>, and the requirement to find an assignment, specifying a value for each binary variable, that makes the formula true. The logic formula can be expressed as a conjunction of $`m`$ clauses and each clauses is a disjunction of some variables. When all the clauses have exactly $`k`$ variables, the problem is called k-SAT, with 1-SAT being the simplest case. Combinatorial search aims at seeking the solution of a SAT from the total of 2<sup>n</sup> assignments. In general, the computational cost of solving a SAT grows exponentially with $`n`$ in the worst case, making the SAT one of the most difficult NP problems. For a few simple cases like 1-SAT and 2-SAT, however, solutions can be quickly and accurately determined without exponential cost using the regular structure of the problem. As each false clause for a given assignment is counted as a conflict, solutions are assignments with no conflicts. The classical search cost for such cases is O($`n`$) due to each clause eliminating one value for a single variable, but the quantum counterpart with Hogg’s algorithm for 1-SAT and maximally constrained k-SAT a single step, with the help of quantum parallelism and interference combined with the problem structure, which shows very high efficiency for large $`n`$.
Actually, Hogg’s algorithm consists of three stages. First, one makes an equal superposition of bases, $`|s>`$, of the $`n`$ two-state quantum system as the initial state $`|\psi _i>=2^{\frac{n}{2}}_s|s>`$, with bit strings $`s`$ being all 2<sup>n</sup> assignments of $`n`$ variables. Secondly, a transformation R is applied on $`|\psi _i>`$. The elements R<sub>ss</sub> of the diagonal matrix R are determined by the conflicts c<sub>s</sub> of $`s`$, which are governed by the constraints expressed in clauses of the logic formula.It is explicitly expressed as
$`R(s)=\{\begin{array}{cc}\sqrt{2}\mathrm{cos}(2c1)\frac{\pi }{4}\hfill & \mathrm{for}\mathrm{even}\mathrm{m}\hfill \\ i^c\hfill & \mathrm{for}\mathrm{odd}\mathrm{m}.(1)\hfill \end{array}`$
Finally, an operation U with the form
$`U_{rs}=\{\begin{array}{cc}2^{\frac{n1}{2}}\mathrm{cos}[(nm+12d)\frac{\pi }{4}]\hfill & \mathrm{for}\mathrm{even}\mathrm{m}\hfill \\ 2^{\frac{n}{2}}e^{i(nm)\frac{\pi }{4}}(i)^d\hfill & \mathrm{for}\mathrm{odd}\mathrm{m},\hfill \end{array}`$
is acted. It can be seen that the matrix elements U<sub>rs</sub> depend only on the Hamming distance d<sub>r,s</sub> between $`r`$ and $`s`$. As it describes the inconsistency of all the corresponding bits of bit strings $`r`$ and $`s`$, d<sub>r,s</sub> has some correlation with conflicts c<sub>s</sub>. By exploiting the intrinsic relation between d<sub>r,s</sub> and c<sub>s</sub> and ingeniously constructing R<sub>ss</sub> and U<sub>rs</sub>with d<sub>r,s</sub> and c<sub>s</sub>, making the combined transformation UR on $`|\psi _i>`$ once is able to find the expected state, i. e., UR $`|\psi _i>=`$ $`|\psi _0>`$, which has equal amplitudes among solutions and no amplitudes among nonsolutions.
In the view of quantum physics, Hogg’s algorithm subtly makes use of the properties of quantum superposition and interference. When the calculating system is prepared on the equal superposition state, $`|\psi _i>=2^{\frac{n}{2}}_s|s>`$, all the possible values, which can be assigned to the logical formula, regardless correct or incorrect, are embodied in the state.Then the state is subjected to the transform $`R`$, which gives a varied phase to each component state of the superposition state according to its conflicts $`c_s`$. The final operation $`U`$ destructively annihilates every component state that its conflict is not zero, construtively enhances the correct component state that its conflicts is zero, finds out the exact component state that satisfies the logical formula. In one word, the essence of searching achieved in a single step is to examine all the assignments simultaneously and to make use of quantum interference embedded in UR. By a single step we mean the number of search steps or sequentially examined assignments, not the number of elementary computational operations required.
Our implementation of Hogg’s algorithm was performed using solution NMR, applied to an ensemble of carbon-13 labelled chloroform molecules $`(^{13}`$CHCl$`{}_{3}{}^{})`$. Two heteronuclear spins <sup>1</sup>H and <sup>13</sup>C in <sup>13</sup>CHCl<sub>3</sub> work as two qubits in our quantum computing. Similar to other experiments\[9-12,14,17-19\] of quantum information processing using NMR, all the logic operations were realized by a specified sequence of radio-frequency (RF) pulses and the spin-spin coupling between the nuclei, and should be finished in the period of time much shorter than the relaxation times T<sub>1</sub> and T<sub>2</sub> of the nuclei in order to minimize the decoherence effect.
A quantum circuit for achieving this highly structured search in a two-qubit system is shown in Fig. 1. The algorithm starts in the state $`|\psi _s>=|00>`$, labelling the states of <sup>1</sup>H and <sup>13</sup>C spins from left to right respectively. As the state of a bulk sample at room temperature is described by a density matrix $`\rho _{th}`$ for a thermally equilibrated system,we prepared an effective pure state $`\rho =|00><00|`$ by temporal averaging, with which a deviation density matrix with diagonal elements of diag($`\rho _s`$)= could be extracted from the summation over cyclic permutations of the populations of the $`|01>`$, $`|10>`$, and $`|11>`$ states three times on the thermal matrix $`\rho _{th}`$. The Hadamard gates H in Fig. 1 were implemented by using RF pulse sequence $`\left(\frac{\pi }{2}\right)_y\left(\pi \right)_x`$ or its equivalents and thus a uniform superposition state denoted by the density matrix $`\rho _i=|\psi _i><\psi _i|`$ with diag($`\rho _i)`$= was obtained. Then two crucial operations R and U were applied, which are related to the problem structure and expressions for the logic formula and lead to the enhancement of solution states and the cancel of nonsolution states.
For our two-qubit system(i. e., $`n=2`$), $`m`$, the number of clauses in the 1-SAT formula with satisfiable solutions can be assumed 2 or 1. When $`m=3`$ and 4, no satisfying assignment exists due to over-constrained formula being unsatisfiable. Four satisfiable formulas in the $`m=2`$ case are $`V_1V_2`$, $`V_1\overline{V}_2`$, $`\overline{V}_1V_2`$ and$`\overline{V}_1\overline{V}_2`$, where $`\overline{V}_i`$ ($`i=1,2`$) are the negation of $`V_i`$. The corresponding solutions are easily found to be $`|11>`$, $`|10>`$, $`|01>`$ and $`|00>`$ respectively, where the logic variable with subscript 1 corresponds to the high qubit and 2 the low qubit. Similarly, the formulas for $`m`$=1 are $`V_2`$, $`V_1`$, $`\overline{V}_2`$ and $`\overline{V}_1`$, corresponding to the unnormalized answer states $`|01>+|11>`$, $`|10>+|11>`$, $`|00>+|10>`$ and $`|00>+|01>`$, respectively. According to the number of states in the solutions, they can be called as single- and multiple- item searching. The algorithm on two instances of 1-SAT was experimentally demonstrated: one for $`m=2`$ with clauses $`V_1`$ and $`V_2`$, and the other for $`m=1`$ with clause $`V_2`$.
First of all, we have to derived the expressions for R and U from Ref.. With Eq.(1), after calculating the conflicts $`c`$ of each string assigned to the logical formula respectively, the elements of the diagonal matrices R$`_{V_1V_2}`$ and R$`_{V_2}`$ were calculated to be $`diag(R_{V_1V_2})=[1,1,1,1]`$ and $`diag(R_{V_2})=[i,1,i,1]`$. U could be represented by $`W\mathrm{\Gamma }W`$, where $`W=H_1H_2`$ with $`H_i`$ being the Hadamard gates, and $`\mathrm{\Gamma }`$ is a diagonal matrix with the diagonal elements
$`\mathrm{\Gamma }_{rr}=\{\begin{array}{cc}\sqrt{2}\mathrm{cos}[(m2|r|1)\frac{\pi }{4}\hfill & \mathrm{for}\mathrm{even}\mathrm{m}\hfill \\ i^{|r|}e^{im\frac{\pi }{4}}\hfill & \mathrm{for}\mathrm{odd}\mathrm{m}.\hfill \end{array}`$
Elements of the diagonal matrices $`\mathrm{\Gamma }`$ for $`m=2`$ and $`m=1`$ were evaluated respectively as $`diag(\mathrm{\Gamma }_{m=2})=[1,1,1,1]`$ and $`diag(\mathrm{\Gamma }_{m=1})=[1,i,i,1]`$.
In order to realize the transformations with NMR, we designed the following RF pulse sequences respectively:
$`R_{V_1V_2}`$: $`\left(\frac{\pi }{2}\right)_{y_1}\left(\frac{\pi }{2}\right)_{x_1}\left(\frac{\pi }{2}\right)_{y_1}\frac{1}{2J}\left(\frac{\pi }{2}\right)_{y_2}\left(\frac{\pi }{2}\right)_{x_2}\left(\frac{\pi }{2}\right)_{y_2}`$
$`R_{V_2}`$: $`\left(\frac{\pi }{2}\right)_{y_1}\left(\frac{\pi }{2}\right)_{x_1}\left(\frac{\pi }{2}\right)_{y_1}`$
$`\mathrm{\Gamma }_{m=2}`$:$`\left(\frac{\pi }{2}\right)_{y_1}\left(\frac{\pi }{2}\right)_{x_1}\left(\frac{\pi }{2}\right)_{y_1}\frac{1}{2J}\left(\frac{\pi }{2}\right)_{y_2}\left(\frac{\pi }{2}\right)_{x_2}\left(\frac{\pi }{2}\right)_{y_2}`$
$`\mathrm{\Gamma }_{m=1}`$: $`\left(\frac{\pi }{2}\right)_{y_1}\left(\frac{\pi }{2}\right)_{x_1}\left(\frac{\pi }{2}\right)_{y_1}\left(\frac{\pi }{2}\right)_{y_2}\left(\frac{\pi }{2}\right)_{x_2}\left(\frac{\pi }{2}\right)_{y_2}`$
To make the best use of the available coherence time and to diminish errors due to the increased number of RF pulses, we finally optimized or reduced the pulse sequence constructed from the executing pulses of H, $`\mathrm{\Gamma }`$ and R mentioned above with the help of NMR principle. The reduced pulse sequence for (UR)$`_{V_1V_2}`$ and (UR)$`_{V_2}`$ shown in Fig. 2 were applied upon $`\rho _i`$, resulting the output matrix $`\rho _0|\psi _0><\psi _0|=(UR)\rho _i(UR)^{}`$.
In order to read out the results of Hogg’s algorithm accurately, a method of quantum state tomography was adopted to reconstruct all the elements of the output density matrix $`\rho _0`$. For this end, we applied different read-out pulses immediately after the searching pulses UR, one for each run. Explicitly, nine pulses E<sub>1</sub>E<sub>2</sub>, E$`{}_{1}{}^{}\left(\frac{\pi }{2}\right)_{x_2}^{}`$, E$`{}_{1}{}^{}\left(\frac{\pi }{2}\right)_{y_2}^{}`$, $`\left(\frac{\pi }{2}\right)_{x_1}`$E<sub>2</sub>, $`\left(\frac{\pi }{2}\right)_{x_1}\left(\frac{\pi }{2}\right)_{x_2}`$, $`\left(\frac{\pi }{2}\right)_{x_1}\left(\frac{\pi }{2}\right)_{y_2}`$, $`\left(\frac{\pi }{2}\right)_{y_1}`$E<sub>2</sub>, $`\left(\frac{\pi }{2}\right)_{y_1}\left(\frac{\pi }{2}\right)_{x_2}`$ and $`\left(\frac{\pi }{2}\right)_{y_1}\left(\frac{\pi }{2}\right)_{y_2}`$were exploited, where E<sub>i</sub> $`(i=1,2)`$ means no pulse acted on the i-th spin. Then we calculated the signal intensities by integrating the proton and carbon spectral lines acquired in each run and reproduced the whole matrix $`\rho _0`$ using the least-square fitting to the coupled equations connecting the signal intensities with the matrix elements. The recovered matrices $`\rho _0`$ from the experimental data were depicted in Fig. 3, along with the theoretical ones for comparison. It is clearly seen from Fig. 3 that only the elements of $`<11|\rho _0(V_1V_2)|11>`$, $`<01|\rho _0(V_2)|01>`$, $`<01|\rho _0(V_2)|11>`$, $`<11|\rho _0(V_2)|01>`$ and $`<11|\rho _0(V_2)|11>`$ are significant, which justifies the searched items to be $`|11>`$ for $`V_1V_2`$ and unnormalized $`|01>+|11>`$ for $`V_2`$ formula. Due to experimental errors, however, several elements that should be zero theoretically still have small amounts of modulus, with the maxima being 8% for single-item and 19% for multiple-item searching.
Errors in the experiments result from several sources. The effective pure state $`|00>`$ prepared by temporal averaging is generally not ideal, and small residual populations in the states $`|01>`$, $`|10>`$ and $`|11>`$ of the deviation density matrix will surely cause errors in the later stage of experiments. The inhomogeneity of RF fields and static magnetic fields and imperfections of the pulse-length calibration are also important sources of errors. Although we have tried to minimize these effects by careful adjustments of and measurements with the apparatus, some factors were still not well under control. Lastly but not insignificantly, losses of the coherence and populations of the density matrices in the whole process of experiments will lead to deviations of measured data from those expected by theory, especially for the off-diagonal elements. Perhaps that is why the maximum error in the multiple-item searching is larger than that in the single-item case. In all senses, minimizing experimental errors in all available ways is essential to accurately operating NMR quantum computers.
We have experimentally demonstrated a quantum algorithm for the highly structured combinatorial searching and found solutions to the 1-SAT problem in one algorithmic step almost surely using an NMR quantum computer with 2 qubits. By contrasting with other search methods that ignore the problem structure requiring O(2<sup>n</sup>) steps classically and O(2$`^{\frac{n}{2}}`$) steps on quantum computers, the searching with Hogg’s algorithm can be accomplished in principle more efficiently. However, there is no difference between Hogg’s algorithm and Grover’s one when only two qubits are considered. The potential of the high efficiency owned by Hogg’s algorithm will be displayed in solving various 1-SAT problems for $`n>2`$ qubits in one step with current NMR technique, which can be readily extended from the present experiment, if the appropriate samples possessing more qubits can be experimentally feasible. Furthermore, Hogg’s algorithm for 1-SAT could be generalized to the maximally constrained k-SAT problems, which have practical examples such as scheduling, finding low energy states of spin glasses and proteins, and automatic theorem proving. While scaling up an NMR computer to much large systems poses daunting challenges, building such a device with a few qubits by some creative approaches for demonstrating algorithms stated above and the others is promising.
We thank Xijia Miao for help in the early stage of experiments. This work was supported by the Chinese Academy of Sciences and National Natural Science Foundation of China.
Captions of the figures
Figure 1 A quantum circuit for implementing a highly structure search algorithm on a two-qubit computer. Two Hadamard gates H$``$H transform an effective pure states $`|\psi _s>=|00>`$ into a uniform superposition state $`|\psi _i>`$, which is then transformed to the answer state $`|\psi _0>`$ after the action of gates R and U. For the definition of R and U, see text.
Figure 2 The reduced NMR pulse sequences used to execute a) (UR)$`_{V_1V_2}`$and b) (UR)$`_{V_2}`$ operations. Narrow and wide boxes correspond to $`\frac{\pi }{2}`$ and $`\pi `$ pulses respectively. X and Y denote the pulses along the x- and y-axis, $`\overline{X}`$ and$`\overline{Y},`$ opposite to the x- and y-axis. The time period $`\tau `$ is set equal to $`\frac{1}{4J}`$, where J is the size of the spin-spin coupling between nuclei $`{}_{}{}^{1}H`$ and $`{}_{}{}^{13}C`$. Experiments were performed using a Bruker ARX500 spectrometer. $`{}_{}{}^{13}C`$-labelled CHCl<sub>3</sub> were obtained from Cambridge Isotope Laboratories.
Figure 3 Experimentally recovered and theoretically expected deviation density matrices, $`\rho _0,`$ after completion of the combinatorial search with the aim of satisfying logic formulas of both V$`{}_{1}{}^{}V_2`$ and V<sub>2</sub>. The ordinate represents the modulus of matrix elements of $`\rho _0`$ (not normalized). The numbers 0, 1, 2 and 3 in the horizontal plane denote the subscripts $`|00>`$, $`|01>`$, $`|10>`$ and $`|11>`$ of the elements, respectively. a) Experiment for $`V_1V_2`$, b) Theory for $`V_1V_2`$, c) Experiment for $`V_2`$ and d) Theory for $`V_2`$.
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# A series representation of the nonlinear equation for axisymmetrical fluid membrane shape
\[
## Abstract
Whatever the fluid lipid vesicle is modeled as the spontaneous-curvature, bilayer-coupling, or the area-difference elasticity, and no matter whether a pulling axial force applied at the vesicle poles or not, a universal shape equation presents when the shape has both axisymmetry and up-down symmetry. This equation is a second order nonlinear ordinary differential equation about the sine $`sin\psi (r)`$ of the angle $`\psi (r)`$ between the tangent of the contour and the radial axis $`r`$. However, analytically there is not a generally applicable method to solve it, while numerically the angle $`\psi (0)`$ can not be obtained unless by tricky extrapolation for $`r=0`$ is a singular point of the equation. We report an infinite series representation of the equation, in which the known solutions are some special cases, and a new family of shapes related to the membrane microtubule formation, in which $`sin\psi (0)`$ takes values from $`0`$ to $`\pi /2`$, is given.
preprint: Submit to PRE
\]
Vesicles are closed surfaces of amphiphilic molecules dissolved in water, which form flexible bilayer membranes in order to minimize contact between the hydrocarbon chains of the lipid and water. Recently, the vesicle shapes have attracted wide interest from different communities such as physics, mathematics, chemistry, and biology. The naive model was given by Canham, in which only the surface bending elasticity is considered. An important progress was the introduction of the spontaneous curvature into the Canham’s theory, and it was made by Helfrich with analogue to the spontaneous splay in the liquid crystal molecular layer. Since the membrane consists in two monolayers, if assuming that two layers depart from each other at a fixed distance, a so-called bilayer-coupling model was explored. However, the individual monolayers can in fact expand elastically under tensile stress, the so-called area-difference elasticity model was introduced and investigated. The latter three models are more realistic than the naive one, and have been intensively studied in recent years. A surprising common property of these three models is that they give exactly the same shapes, while differing from each other in accounting for the shape transitions, such as budding and vesiculation transition, transition from prolates to oblates etc. . Furthermore, pulling or pushing the vesicle axially making it up-down and axis symmetrical, as it used in the experimental investigation of the tether formation, does not complicate the form of the mathematical equation for the general axisymmetrical vesicle in the spontaneous curvature model without force included. But the meanings of parameters need to re-specify. Except some numerical approach using the powerful software Surface Evolver, usual numerical as well as analytical method is employed to solve the shape equation representing axisymmetrical vesicle. However, in analytical side, there is not a general applicable method to study the equation. In numerical side, $`r=0`$ (cf. the following Eq.(1) ) is a singular point of the equation, so the shape around the point cannot be obtained unless by tricky extrapolation. In this paper, we are going to reveal a novel property of the equation: the equation has an equivalent series representation, in which the two above problems do not present anymore. As a meaningful demonstration, a new family of shape related to the tether formation will be given.
Since the detailed or key steps of the derivation of the same equation from different models is available in many papers, e.g., we can start our further discussion from the equation itself with clearly specifying the physics meaning of each quantity in it. This may of course makes this paper more concise. And the equation we start from reads:
$`cos^2\psi {\displaystyle \frac{d^2\psi }{dr^2}}{\displaystyle \frac{sin(2\psi )}{4}}({\displaystyle \frac{d\psi }{dr}})^2+{\displaystyle \frac{cos^2\psi }{r}}{\displaystyle \frac{d\psi }{dr}}`$ (1)
$`{\displaystyle \frac{sin(2\psi )}{2r^2}}{\displaystyle \frac{\delta pr}{2kcos\psi }}`$ (2)
$`{\displaystyle \frac{sin\psi }{2cos\psi }}({\displaystyle \frac{sin\psi }{r}}c_0)^2{\displaystyle \frac{\lambda sin\psi }{kcos\psi }}={\displaystyle \frac{C}{rcos\psi }}.`$ (3)
The quantities in this equation need some explanations. For an axisymmetrical shape, we can choose the symmetrical axis to be the $`z`$-axis. The contour of the shape can be plotted in $`rz`$ plane, with $`r`$ being the radial coordinate denoting the distance from the symmetric $`z`$ axis. Then the tangent angle $`\psi (r)`$ of the contour is measured clockwise from $`r`$ axis. Using $`s`$ to denote the arclength along the contour measured from the north pole of the shape and the $`\varphi `$ the azimuthal angle, we have the following geometrical relations:
$`\{\begin{array}{c}\frac{dr}{ds}=cos\psi (r)\hfill \\ \frac{dz}{ds}=sin\psi (r)\hfill \\ z(r)z(0)=_0^rtan\psi (r)𝑑r\hfill \end{array}`$ (7)
$`𝐧=(sin\psi cos\varphi ,sin\psi sin\varphi ,cos\psi )`$ (8)
in which $`𝐧`$ denotes the normal of the surface. As to constants $`k`$, $`C`$, $`c_0`$, $`\delta p`$, and $`\lambda `$, they have different meaning in different models. Taking the spontaneous curvature theory as example, $`k`$ the elastic modulus, $`C`$ the integral constant, $`c_0`$ the spontaneous curvature, $`\delta p`$ and $`\lambda `$ the two Lagrangian multipliers taking account for the constraints of constant volume and area, which may be physically understood as the osmotic pressure between the ambient and the internal environments, and the surface tension coefficient, respectively. When $`C=0`$, Eq. (3) reduces to be the original Helfrich shape equation. When $`c_0N`$, $`(2\lambda +c_0^2)/4L`$, $`\delta p6M`$, $`C4F`$, Eq.(1) becomes the equation representing the shapes with the same symmetry in area-difference elasticity model with the inclusion of the axially external force $`F`$ (Eq. (19) in ).
Although no one attempts to analytically solve it in a systematical way, we know that the following simple functions.
$$sin\psi (r)=ar+b/r+c,\text{ and},ar+c_0rlnr.$$
(9)
solve the equation. Adjusting the constants $`a`$, $`b`$, and $`c`$ in permissible ranges, the former form of the solution gives following shapes: 1) cylinder $`r=r_0`$, $`sin\psi (r_0)=1`$; 2) the axisymmetrical constant mean curvature surfaces ( the so-called the Delaunay surfaces) with $`c=0`$ and an extension with $`c0`$; and 3) Clifford torus: $`sin\psi (r)=\sqrt{2}+r`$. The latter form of solution gives shapes from the self-intersecting oblates, oblate, sphere, prolate to very long capped prolate. These analytical solutions not only have explained the existing but puzzling shapes, but also had predicted something new. Clifford torus, for instance, was first theoretically predicted, then experimentally confirmed.
We will show that Eq.(3) can be analytically solved for the physically interesting case. For convenience, we make a transformation for Eq. (3)
$$\psi (r)=arcsinf(r)$$
(10)
and rescaling the unit such that $`k=1`$. Then we get from Eq.(3) an equation for $`f(r)`$:
$`2Cr\delta pr^32f(r)c_0^2r^2f(r)2\lambda r^2f(r)`$ (11)
$`+2c_0rf^2(r)+f^3(r)+2krf^{}(r)2rf^2(r)f^{}(r)`$ (12)
$`+r^2f(r)f_{}^{}{}_{}{}^{2}(r)+2r^2f^{\prime \prime }(r)2r^2f^2(r)f^{\prime \prime }(r)=0`$ (13)
This is a second order nonlinear ordinary differential equation (ODE) not belonging to any well-studied mathematical type. For our purpose to attack it, let us recall the Taylor series expansion of a function. As well-known, for any analytical function $`f(r)`$ defining in a closed interval $`r[r_1,r_2]`$, global Taylor expansion of $`f(r)`$ around a point $`r_0`$ is
$`f(r)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{f^{(k)}(r_0)(rr_0)^k}{k!}}+R_n(rr_0)`$ (14)
$`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{f^{(k)}(r_0)(rr_0)^k}{k!}},`$ (15)
where $`R_n(rr_0)=(rr_0)^{(n+1)}f^{(n+1)}(\xi )`$ with $`r_1<\xi <r_2`$. Both the above series are exact and equivalent to each other. The first finite series is easy to serve the digit purpose, while the second infinite one is useful to obtain a closed expression once the general form is determined. To note that every axisymmetrical shape must intersect the equatorial plane at right angle $`\psi =\pi /2`$. This means that there exits a point $`r_0`$, such that $`f(r_0)=1`$. At this point $`r_0`$, the two terms containing the second derivative in Eq.(13) cancel. We have then value of its first derivative $`f^{}(r_0)`$ of $`f(r)`$ at point $`r_0`$ as
$$f^{}(r_0)=\pm \frac{\sqrt{1+2Cr_02c_0r_0+c_0^2r_0^2+2\lambda r_0^2+\delta pr_0^3}}{r_0}.$$
(16)
Since $`f(r)f(r_0)=1`$, we have $`f^{}(r_0)=\underset{rr_0}{lim}(f(r)1)/(rr_0)0`$ if and only if $`r_0`$ is the maximum radius of vesicle. Since we treat $`f(r)`$ a single valued function of $`r`$ and $`r_0`$ is really the maximum radius of vesicle, only the positive $`f^{}(r_0)`$ is relevant. To note that the physically interesting vesicles are either free or axially forced, and there is no force or torque acting on the waist. It implies that the high rank derivatives of $`f(r)`$ at $`r_0`$ as $`f^{(k)}(r_0),k=2,3\mathrm{}`$ exist, and they can be obtained by the following method. Taking derivative with respect to variable $`r`$ in both side, Eq.(13) becomes a third order differential equation. Putting $`r=r_0`$ and substituting the value $`f(r_0)`$ $`(=1)`$ and $`f^{}(r_0)`$ (given by Eq. (16) ) into the third order equation, the third order terms cancel, and we can obtain the value of its second derivative $`f^{\prime \prime }(r_0)`$. Similarly, taking derivative with respect to variable $`r`$ in both side of the third order derivative equation, we can obtain $`f^{(3)}(r_0)`$, and so forth, all its higher rank derivatives $`f^{(k)}(r_0),k=4,5\mathrm{}`$ can be determined. It is a remarkable result: the nonlinear ODE (13) can be transformed into an infinite series with the recurrence relation between coefficients uniquely determined by the equation itself. In the mathematical point of view, all its solutions with physical significance are obtained.
In this paper, we only treat the simplest aspect of the solution without considering the congruence of the segments of shapes. Since the parameterization of the axisymmetrical surface is $`\psi (r)`$ rather than $`\psi (s)`$, from Eq.(1) $`f(r)`$ is a single value function of $`r`$ in interval $`\psi (r)[0,\pi /2]`$. We can therefore expand $`f(r)`$ in Taylor series directly. Conversely, once all Taylor coefficients $`f^{(k)}(r_0)/k!,(k=1,2,3,\mathrm{})`$, are known, the function $`f(r)`$ is determined from the relation (15). In principle, we can give such a series for each shape. As checks of our method, we like to give two examples. 1) Cylinder with arbitrary radius $`r=r_0`$. It is the case $`f(r_0)1`$ provided all parameters $`k,\delta p,\lambda ,c_0,C`$ satisfying $`\delta pr_0^3+2\lambda r_0^2+k(c_0r_01)^2=0`$ and $`C=0`$. It is the result discussed previously. 2) The Delaunay surfaces $`f(r)=ar+b/r`$ (a, b are two constants). These constant mean curvature surfaces appear when $`\lambda =\delta p=C=0`$. Now we have simply $`f^{(1)}(r_0)=\pm (c_01/r_0)`$ at position $`r_0`$ satisfying $`f(r_0)=1`$. We study the case with positive sign. We can easily find $`f^{(k)}(r_0)=(1)^{(k+1)}k!(c_0r_02)/(2r_0^k)`$, $`(k=2,3,4\mathrm{})`$. Then the solution can be exactly written as $`f(r)=ar+b/r`$ with $`a=c_0/2`$, $`b=(r_0c_0r_0^2/2)`$. It is the result discussed also previously. One can verify this method using another analytical or numerical solution as exercises, and can easily find that they are some special cases of the infinite series with special relation between the parameters.
Now we wish to give new family of shapes. These shapes can be inferred from the works of three research groups in studying the meaning of the integral constant $`C`$. First by Zheng and Liu, they noted that the constant $`C`$ roughly meant that the constant pressure difference $`\delta p`$ in Helfrich shape equation was replaced by $`\delta p+2kC/r^2`$. Since that the all Helfrich shapes with axisymmetry are smooth at the poles $`r=0`$, the constant $`C`$ means a stress singularity at two poles. This singularity may lead to two horns at both poles. Secondly, Jülicher and Seifert claimed that the nonvanishing $`C`$ connected to shapes with torus topology. Thirdly, Podgornik, Svetina, and Z̆eks̆ in and Boz̆ic̆, Svetina, and Z̆eks̆ in showed that a point axial force can lead to such a constant $`C`$. We have known that such a force can lead to microtubule formation of vesicle , and in real experiment a finite force is enough to produce such structure. However, the expected shapes with two horns at both poles as exerting finite force have not been obtained theoretically.
For our purpose, we treat a case in which the quantity in the square of Eq.(16) to form the square of a quantity with nonzero $`C`$. For our purpose, we choose the simplest parameters as $`C=c_0`$, $`\delta p=0,\lambda =c^2/2`$, and the positive sign of $`f^{}(r_0)`$. The first seven derivatives $`f^{(k)}(r_0)`$, $`(k=0,1,2,\mathrm{..6})`$ are
$`\{1,{\displaystyle \frac{1}{r_0}},{\displaystyle \frac{2C}{r_0}},{\displaystyle \frac{4C}{r_0^2}},{\displaystyle \frac{2\left(35C+4C^2r_0\right)}{5r_0^3}},`$ (17)
$`{\displaystyle \frac{2C\left(1155+268Cr_0+12C^2r_0^2\right)}{35r_0^4}},`$ (18)
$`{\displaystyle \frac{8C\left(5040+1689Cr_0+99C^2r_0^2+32C^3r_0^3\right)}{105r_0^5}}\}.`$ (19)
Since the conformal invariance of the shape, the magnitude of $`C`$ does not matter. We can therefore choose $`C=1`$. Then the Taylor series Eq.(15) can not give real and positive root of $`r_0`$ for equation $`f(0)`$= a negative number. It means that in this case the equation is not self-consistent. In contrast, it can give a unique and positive root for equation $`f(0)`$= a positive constant. To note that $`arcsinf(0)`$ is the angle between the tangent of the contour at pole and the $`r`$ axis. For a sequence of numbers of $`f(0)`$ as
$$N=\{0,0.1,0.2,0.3,0.4,0.5,0.6,0.7,0.8,0.9,0.972,1\},$$
(20)
we have form the equation $`f(0)=N`$ a corresponding sequence of unique and (semi-)positive roots for $`r_0`$ as
$`r_0`$ $`=`$ $`\{0,0.0095195,0.0189412,0.0282679,0.0375022,`$ (23)
$`0.0466463,0.0557024,0.0646727,0.0735589,`$
$`0.0823627,0.0886836,0.0910858\}.`$
Do not care the magnitude of these values. The energy of the vesicle depends on the shape due to the conformal invariance. For clearly specifying each shape, a conformal invariant, the so called relative volume $`v=V/(4\pi R^3/3)`$ with $`R=\sqrt{A/4\pi }`$ is usually used, where $`V`$ and $`A`$ are the volume and area of the vesicle respectively. The one-to-one corresponding relative volumes are:
$`v`$ $`=`$ $`\{1,0.999971,0.999879,0.999725,0.999505,`$ (26)
$`0.999213,0.998844,0.998386,0.997819,`$
$`0.997107,0.996449,0.996136\}`$
In FIG. 1, the two limit shapes, sphere (thick solid line) and the prolate with the sharpest horn (thin line), are plotted. Other shapes with the horns of angles $`\psi (0)`$ satisfying $`0<\psi (0)<\pi /2`$ will occupy the domain between these two.
For studying the relation between the “external force” $`C`$ and the relative volume, we define a dimensionless external force as $`C(z_0r_0)`$ with $`z_0`$ being half of the length of the prolate, and rescale the relative volume as $`v^{}=4.725(1v)`$. We can find that both the dimensionless external force and the rescaled relative volume have nearly the same relationship to the zenith angle, as shown in FIG. 2. Results plotted in FIG.1 and FIG.2 clearly show that the zenith angle becomes sharper as the external force becomes larger, or equivalent, the less the relative volume, the sharper the horn.
In summary, we have found that the second order nonlinear different equation for axisymmetrical fluid vesicle shapes with up-down symmetry has an equivalent Taylor series representation, in which the recurrence relation of Taylor coefficients are uniquely and completely determined be the equation. Using this representation, we can not only reproduce the known solutions, but also new ones. The new solutions discussed in this article give the shapes with horns at two poles. These horns can never form simultaneously if no axially external force exerted. So, our study directly answers a question whether the constant $`C`$ is related to the presence of the axially external force. This result is compatible with that obtained previously. The obtained shapes are related to the formation of the membrane microtubule on the axially strained vesicles, and finite force is enough to lead to such structure. Results also show that the shapes with less relative volume is easier to form such microtubule.
Finally, we give two comments. First, even the Taylor series method is useful in obtaining and analyzing the solution to the equation, one should be careful to check the convergence of the series. In fact, we explore many cases with different parameter set of $`\delta p,\lambda ,c_0,C`$, the convergence is sometime uncertain. Secondly, the presence of the constant $`C`$ may not be the sufficient condition to relate the axial force, especially in the torus topology situation.
Acknowledgments
One of the authors (Liu) is indebted to Profs. Peng Huan-Wu, Vipin Srivastava and Dr. ZhouHaijun for enlightening discussions. This subject is supported in part by grants from the Hong Kong Research Grants Council (RGC), the Hong Kong Baptist University Faculty Research Grant (FRG), and in part by National Natural Science Foundation of China.
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# Power Spectrum Analysis of the ESP Galaxy Redshift Survey
## 1 Introduction
The quantitative characterisation of the galaxy distribution is a major aim in the study of the large-scale structure of the Universe. During the last 20 years, several surveys of galaxy redshifts have shown that galaxies are grouped in clusters and superclusters, drawing structures surrounding large voids (see e.g. Guzzo 1999 for a review). The power spectrum of the galaxy distribution provides a concise statistical description of the observed clustering that, under some assumptions on its relation to the mass distribution, represents an important test for different structure formation scenarios (e.g. Peacock 1997 and references therein). Indeed, under the assumption of Gaussian fluctuations, the power spectrum totally describes the statistical properties of the matter density field (e.g. Peebles 1980).
In recent years, several estimates of the galaxy power spectrum have been obtained using galaxy samples selected at different wavelenghts: radio (Peacock & Nicholson 1991), infrared (Feldman, Kaiser & Peacock 1994, Fisher et al. 1993, Sutherland et al. 1999) and optical (Park et al. 1994, da Costa et al. 1994, Tadros & Efstathiou 1996, Lin et al. 1996, Hoyle et al. 1999), to mention the most recent ones.
The ESO Slice Project redshift survey (ESP, Vettolani et al. 1997, 1998) is one of the two deepest wide–angle surveys currently available, inferior only to the larger Las Campanas Redshift Survey (LCRS, Shectman et al. 1996). During the last few years, it has produced a number of statistical results on the properties of optically–selected galaxies, as e.g. the luminosity function (Zucca et al. 1997) or the correlation function (Guzzo et al. 2000). The geometry of the survey (a thin row of circular fields, resulting in an essentially 2D slice in space) is such that an estimate of the power spectrum represents a true challenge. In this paper we present the results of a detailed analysis that overcomes these difficulties, producing a reliable measure of the power spectrum from the ESP redshift data. The technique developed here to cope with the specific geometry of the survey is potentially interesting also for application to other surveys consisting of separate patches on the sky, as could be the case, for example, of preliminary sub–samples of the ongoing SDSS (Margon 1998) and 2dF (Colless 1998) surveys.
The outline of the paper is as follows. We shall first recall the main features of the ESP survey and the sample selection (section 2), then discuss the power spectrum estimator adopted for the analysis (section 3) and the numerical tests performed in order to check its validity range (section 4). We shall then present the estimated power spectrum (section 5) and its consistency with the correlation function (section 6), and then discuss it in comparison to the results from other surveys (section 7). Section 8 summarises the results obtained, drawing some conclusions.
## 2 The ESO Slice Project
The ESO Slice Project galaxy redshift survey (ESP, Vettolani et al. 1997, 1998) was constructed between 1993 and 1996 to fill the gap that existed at the time between shallow, wide angle surveys as the CfA2, and very deep, one-dimensional pencil beams. The survey was designed in order to allow the sampling of volumes larger than the maximum sizes of known structures and an unbiased estimate of the luminosity function of field galaxies to very faint absolute magnitudes. The survey and the data catalogue are described in detail in Vettolani et al. (1997, 1998). Here we limit ourselves to a summary of the main features relevant for the present analysis.
The ESP survey (see Figures 1 and 2) extends over a strip of $`\alpha \times \delta =22\mathrm{°}\times 1^\mathrm{o}`$, plus a nearby area of $`5\mathrm{°}\times 1\mathrm{°}`$, five degrees west of the main
strip, in the South Galactic Pole region ($`22^h30^m\alpha 01^h20^m`$, at a mean declination of $`40^\mathrm{o}15^{}`$ (1950)). This region was covered with a regular grid of adjacent circular fields, with a diameter of 32 arcmin each, corresponding to the field of view of the multifibre spectrograph OPTOPUS (Avila et al. 1989) at the ESO 3.6 m telescope. The total solid angle covered by the survey is 23.2 square degrees and its position on the sky was chosen in order to minimize galactic absorption ($`75^\mathrm{o}b^{II}60^\mathrm{o}`$). The target objects, with a limiting magnitude $`b_J19.4`$, were selected from the Edinburgh–Durham Southern Galaxy Catalogue (EDSGC, Heydon–Dumbleton et al. 1989). A total of 4044 objects were observed, corresponding to $`90\%`$ of the parent photometric sample and selected to be a random subset of the total catalogue with respect to both magnitude and surface brightness. The total number of confirmed galaxies with reliable redshift measurement is 3342, while 493 objects turned out to be stars and 1 object is a quasar at redshift $`z1.174`$. No redshift measurement could be obtained for the remaining 208 spectra. As discussed in Vettolani et al. (1998), the magnitude distribution of the missed galaxies is consistent with a random extraction of the parent population. About half of the ESP galaxies present spectra with emission lines. Particular attention was paid to the redshift quality and several checks were applied to the data, using 1) multiple observations of $`200`$ galaxies, 2) $`750`$ galaxies for which the redshift from both absorption and emission line is available (Vettolani et al. 1998, Cappi et al. 1998). More details about the data reduction and sample completeness are reported in Vettolani et al. (1997, 1998).
Given the magnitude–limited nature of the survey, the computation of a clustering statistics like the power spectrum requires the knowledge of the selection function. This is defined as the expected probability to detect a galaxy at a redshift $`z`$ and can be expressed as
$$s(z)=\frac{_{\mathrm{max}[L_1,L_{\mathrm{min}}(z)]}^+\mathrm{}\varphi (L)𝑑L}{_{L_1}^+\mathrm{}\varphi (L)𝑑L},$$
(1)
where $`\varphi (L)`$ is the luminosity function, $`L_1`$ is the minimum luminosity of the sample and $`L_{\mathrm{min}}(z)`$ is the minimum luminosity detectable at redshift $`z`$, given the sample limiting magnitude.
In the ESP survey the minimum luminosity corresponds to an absolute magnitude $`M_{b_J,1}=12.4+5\mathrm{log}h`$ ($`h`$ is the Hubble constant in units of $`100\mathrm{k}ms^1\mathrm{M}pc^1`$). $`L_{\mathrm{min}}(z)`$ is the luminosity of a galaxy at redshift $`z`$ with an apparent magnitude equal to the apparent magnitude limit $`b_J=19.4`$. The corresponding absolute magnitude is given by
$$b_JM_{b_J}=25+5\mathrm{log}D_L(z)+K(z),$$
(2)
where $`D_L`$ is the luminosity distance in Mpc and $`K(z)`$ is the K-correction. $`D_L`$ is given by the Mattig formula (1958), which depends on the assumed cosmological model. For all ESP computations we assume a flat universe with $`\mathrm{\Omega }_\mathrm{o}=1`$ and $`\mathrm{\Lambda }=0`$. Before proceeding to the computation of luminosity distances, we have converted the observed heliocentric redshifts in the catalogue to the Cosmic Microwave Background (CMB) rest frame using a standard procedure, as described in Carretti (1999). The luminosity distance is then given by
$$D_L(z)=\frac{2c}{H_\mathrm{o}}(1+z)\left(1\frac{1}{\sqrt{1+z}}\right)$$
(3)
The K-correction is a function of redshift and morphological type, but the latter is not directly available for ESP galaxies. Following Zucca et al. (1997), we use an average K-correction, weighted over the expected morphological mixture at each $`z`$. See Zucca et al. (1997, cfr. their figure 1) for the details of this computation. A recent principal component analysis of the spectra (Scaramella, priv. comm.) confirms the adequacy of this mean correction.
The luminosity function is such that $`\varphi (L)dL`$ gives the density of galaxies with luminosity $`L[L,L+dL[`$. The ESP luminosity function is well approximated by a Schechter (1976) function (Zucca et al. 1997)
$$\varphi (L)dL=\varphi ^{}\left(\frac{L}{L^{}}\right)^\alpha e^{L/L^{}}d\left(\frac{L}{L^{}}\right)$$
(4)
with best fit parameters $`\alpha =1.22`$, $`M_{b_J}^{}=19.61+5\mathrm{log}h`$ and $`\varphi ^{}=0.020h^3`$ Mpc<sup>-3</sup>. In reality, as shown by Zucca et al. (1997), for $`M_{b_J}>16+5\mathrm{log}h`$ the faint end steepens with respect to the Schechter form and the overall shape is better described by adding an extra power law. Nevertheless, this is relevant only for the very local part of the sample and a description of the selection function using a simple Schechter fit is perfectly adequate for our purposes.
Another quantity to be taken into account for clustering analyses is the redshift completeness of the 107 fields, as not all target galaxies at the photometric limit were succesfully measured.
This can be expressed as (Vettolani et al. 1998)
$$C=\frac{N_Z}{N_TN_S0.122N_{NO}},$$
(5)
where, for each field, $`N_T`$ is the total number of objects in the photometric catalogue, $`N_Z`$ is the number of reliable galaxy redshifts, $`N_{NO}`$ is the number of not observed objects, $`N_S`$ is the number of stars and 0.122 is the fraction of stars in the spectroscopic sample. In Figure 3 we plot the completeness values for each field. Field numbers $`<100`$ denote fields in the northern row, while the others refer to the southern one.
The power spectrum analysis has been performed on both volume–limited and magnitude–limited subsamples of the survey. Volume–limited samples include all galaxies intrinsically more luminous than a given absolute magnitude $`M_{\mathrm{lim}}`$ and within the maximum redshift $`z_{\mathrm{max}}`$ at which such magnitude can still be detected within the survey apparent magnitude limit. In such a case, the expected mean density of galaxies does not vary with distance. Magnitude–limited catalogues, by definition, are simply subsets of all galaxies in the survey to a given apparent magnitude, possibly with the addition of an upper distance cut $`z_{\mathrm{max}}`$ above which the value of the selection function becomes too small. Magnitude–limited samples contain more objects, but the mean ensemble properties (as e.g. the galaxy luminosity distribution) vary with distance. We extract from the ESP survey two magnitude–limited samples with different $`z_{\mathrm{max}}`$ limit, plus one volume–limited sample with $`M_{\mathrm{lim}}20.1+5\mathrm{log}h`$. (For simplicity, we shall omit hereafter the $`5\mathrm{log}h`$ term).
For the estimate of the power spectrum, comoving distances are computed for each galaxy as $`D_c(z)=D_L(z)/(1+z)`$. The uncertainty introduced in $`D_c`$ because of our ignorance of the correct cosmological model amounts to less than 5% for a typical redshift $`z=0.20`$, when the value of $`\mathrm{\Omega }_\mathrm{o}`$ is changed from 1 to 0.2.
## 3 Power Spectrum Estimator
The galaxy power spectrum can be defined as
$$P(k)=\xi (x)e^{i𝐤𝐱}𝑑𝐱,$$
(6)
where $`\xi (x)`$ is the two–point correlation function, $`𝐱`$ and $`𝐤`$ are the comoving position and wavenumber vectors respectively, while $`x=|𝐱|`$ and $`k=|𝐤|`$. Under the hypothesis of homogeneity and isotropy, $`P(k)`$ and $`\xi (x)`$ are functions only of $`k`$ and $`x`$ respectively. By definition the two–point correlation function can be also written
$$P(k)|\widehat{\delta }(k)|^2,$$
(7)
where $`\widehat{\delta }(𝐤)`$ is the Fourier transform of the density contrast of the galaxies.
In this paper we follow the Fourier notation
$$\widehat{f}(𝐤)=f(𝐱)e^{i𝐤𝐱}𝑑𝐱,$$
(8)
$$f(𝐱)=\frac{1}{(2\pi )^3}\widehat{f}(𝐤)e^{i𝐤𝐱}𝑑𝐤.$$
(9)
To compute the power spectrum of galaxy density fluctuations from the observed galaxy distribution, we use a traditional Fourier method (cfr. Carretti 1999 for details), as developed by several authors (e. g. cfr. Peebles 1980, Fisher et al. 1993, Feldman et al. 1994, Park et al. 1994, Lin et al. 1996). We also apply a correction (Tegmark et al. 1998), that accounts for our ignorance on the true value of the mean density of galaxies (Peacock & Nicholson 1991).
Given a sample of $`N`$ galaxies of positions $`𝐱_j`$ and weights $`w_j`$, an estimate of the Fourier transform of density contrast is given by
$$\widehat{\stackrel{~}{\delta }}(𝐤)=\frac{V}{_{j=1}^Nw_j}\underset{j=1}{\overset{N}{}}w_je^{i𝐤𝐱_j}\widehat{W}(𝐤),$$
(10)
where $`V`$ is the volume of the sample and $`\widehat{W}(𝐤)`$ is the Fourier transform of the survey window function (hereafter a $``$ will denote the quantities estimated from the data). The window function $`W(𝐱)`$ is 1 within the volume covered by the sample and 0 elsewhere, so it can be described as an ensemble of 107 cones. This geometry allows us to obtain analytically the Fourier transform $`\widehat{W}(𝐤)`$ as the sum of the Fourier transform of each cone. In equation 10 each galaxy contributes with some weight $`w_j`$. In a volume–limited catalogue all galaxies have equal weight, i.e. $`w_j1`$. In a magnitude–limited catalogue the expected galaxy density decreases with the distance according to the selection function. Thus, the simplest form for the weight for a galaxy is given by the inverse of the selection function at its redshift $`z_j`$
$$w_j=\frac{1}{s(z_j)}.$$
(11)
If the catalogue completeness is $`C<1`$, the previous weight should be modified as
$$w_j=\frac{1}{C(𝐱_j)}$$
(12)
for volume–limited catalogues, and as
$$w_j=\frac{1}{s(z_j)C(𝐱_j)},$$
(13)
for magnitude–limited catalogues. $`C(𝐱_j)`$ is the completeness of the sample at the position of the $`j^{th}`$ galaxy.
Our adopted power spectrum estimator is defined with respect to $`\widehat{\stackrel{~}{\delta }}(𝐤)`$ by the following equation (Tegmark et al. 1998)
$$\stackrel{~}{P_c}(𝐤)=\frac{|\widehat{\stackrel{~}{\delta }}(𝐤)|^2\stackrel{~}{b}(𝐤)}{A(𝐤)},$$
(14)
where
$$A(𝐤)=\left(1\left|\frac{\widehat{W}(𝐤)}{\widehat{W}(\mathrm{𝟎})}\right|^2\right)V$$
(15)
accounts for our ignorance of the mean galaxy density, while
$$\stackrel{~}{b}(𝐤)=\frac{V^2}{\left(_{j=1}^Nw_j\right)^2}\underset{j=1}{\overset{N}{}}w_j^2\left|e^{i𝐤𝐱}\frac{\widehat{W}(𝐤)}{\widehat{W}(\mathrm{𝟎})}\right|^2$$
(16)
is the shot noise correction due to the finite size of the sample.
The observed power spectrum estimated by equation 14 is related to the true power spectrum $`P(k)`$ by
$$\stackrel{~}{P_c}(𝐤)=\frac{1}{(2\pi )^3A(𝐤)}P(k^{})\varphi (𝐤,𝐤^{})𝑑𝐤^{},$$
(17)
where
$$\varphi (𝐤,𝐤^{})=\left|\widehat{W}(𝐤𝐤^{})\frac{\widehat{W}(𝐤)}{\widehat{W}(\mathrm{𝟎})}\widehat{W}(𝐤^{})\right|^2.$$
(18)
For wavenumbers $`𝐤`$ such that $`|\widehat{W}(𝐤)||\widehat{W}(\mathrm{𝟎})|`$ this equation reduces to the convolution between $`P(k)`$ and $`|\widehat{W}(𝐤)|^2`$.
To describe the convolved power spectrum we choose to average $`\stackrel{~}{P_c}(𝐤)`$ over all directions
$`\stackrel{~}{P_c}(k)`$ $`=`$ $`\stackrel{~}{P_c}(𝐤)`$ (19)
$`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle _{𝛀_k}}\stackrel{~}{P_c}(𝐤)𝑑𝛀_k`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}k_{}^{}{}_{}{}^{2}P(k^{})\chi (k,k^{})𝑑k^{},`$
where $`𝛀_k`$ is the sphere defined by wavenumbers of amplitude $`k`$. The kernel of this integral equation is given by
$$\chi (k,k^{})=\frac{1}{2(2\pi )^4V}_{𝛀_k}_{𝛀_k^{}}\psi (𝐤,𝐤^{})𝑑𝛀_k𝑑𝛀_k^{},$$
(20)
where
$$\psi (𝐤,𝐤^{})=\frac{\left|\widehat{W}(𝐤𝐤^{})\widehat{W}(𝐤^{})\widehat{W}(𝐤)/\widehat{W}(\mathrm{𝟎})\right|^2}{1\left|\widehat{W}(𝐤)/\widehat{W}(\mathrm{𝟎})\right|^2}$$
(21)
and $`𝛀_k^{}`$ is the sphere defined by wavenumbers of amplitude $`k^{}`$.
The Fourier transform of the window function has been analytically computed as the sum of the Fourier transforms of all the 107 cones. In reality, the cones are slightly overlapped but the small common volume (2.85%) allows us to make the assumption of disjoined cones.
The small width of one ESP cone allows us to analytically compute its Fourier transform. Let $`r_\mathrm{o}`$ be the cone height and $`\mathrm{\Delta }\theta 1`$ rad its width ($`\mathrm{\Delta }\theta =16^{}=0.00465`$ rad). In the simple case of a cone centered on the $`z`$ axis, the Fourier transform is
$$\widehat{W}_c(𝐤)=_0^{r_\mathrm{o}}𝑑rr^2_0^{2\pi }𝑑\varphi _0^{\mathrm{\Delta }\theta }𝑑\theta \mathrm{sin}\theta e^{i𝐤𝐫}.$$
(22)
Taking into account the small value of $`\mathrm{\Delta }\theta `$, the integrand can be approximated to first order in $`\theta `$, resulting in
$`\widehat{W}_c(𝐤)`$ $`=`$ $`{\displaystyle _0^{r_\mathrm{o}}}𝑑rr^2{\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _0^{\mathrm{\Delta }\theta }}𝑑\theta \mathrm{sin}\theta e^{ikr\mathrm{cos}\alpha }`$ (23)
$`=`$ $`2\pi [1\mathrm{cos}(\mathrm{\Delta }\theta )]{\displaystyle _0^{r_\mathrm{o}}}𝑑rr^2e^{ikr\mathrm{cos}\alpha }.`$
This expression depends on $`r_\mathrm{o}`$, $`\mathrm{\Delta }\theta `$ and $`\mathrm{cos}\alpha `$, where $`\alpha `$ is the angle between $`𝐤`$ and the $`z`$ axis. For a generic cone along the direction $`(\theta _\mathrm{o},\varphi _\mathrm{o})`$ a rotation can be applied in order to bring the cone on the $`z`$ axis. So, the Fourier transform is
$`\widehat{W}_c(𝐤)`$ $`=`$ $`{\displaystyle _0^{r_\mathrm{o}}}𝑑rr^2{\displaystyle _{𝛀_{(\theta _\mathrm{o},\varphi _\mathrm{o})}}}𝑑𝛀e^{i𝐤𝐫}`$ (24)
$``$ $`2\pi [1\mathrm{cos}(\mathrm{\Delta }\theta )]{\displaystyle _0^{r_\mathrm{o}}}𝑑rr^2e^{ikr\mathrm{cos}\gamma },`$
where $`𝛀_{(\theta _\mathrm{o},\varphi _\mathrm{o})}`$ is the solid angle of the cone centered on $`(\theta _\mathrm{o},\varphi _\mathrm{o})`$, and $`\gamma `$ is the angle between the wave number $`𝐤`$ and the direction $`(\theta _\mathrm{o},\varphi _\mathrm{o})`$. By solving the integral one gets
$`\widehat{W}_c(𝐤)`$ $`=`$ $`2\pi i{\displaystyle \frac{[1\mathrm{cos}(\mathrm{\Delta }\theta )]}{k\mathrm{cos}\gamma }}\times `$ (25)
$`\left[r_\mathrm{o}^2e^{ikr_\mathrm{o}\mathrm{cos}\gamma }i{\displaystyle \frac{2r_\mathrm{o}e^{ikr_\mathrm{o}\mathrm{cos}\gamma }}{k\mathrm{cos}\gamma }}2{\displaystyle \frac{e^{ikr_\mathrm{o}\mathrm{cos}\gamma }1}{(k\mathrm{cos}\gamma )^2}}\right].`$
Finally, the Fourier transform of the whole ESP window function is given by the sum of the Fourier transform of the cones normalized to the true volume of the survey to account for the overlapping zone
$$\widehat{W}(𝐤)=\frac{V}{107V_c}\underset{i=1}{\overset{107}{}}\widehat{W}_{c,i}(𝐤),$$
(26)
where $`i`$ is the cone index, $`V_c`$ is the volume of one cone and $`V`$ is the true survey volume
$$V=107V_c(1\beta ),$$
(27)
which accounts for the total volume fraction $`\beta =0.0285`$ lost in the cone overlaps.
To check the reliability of our analytic approximation, we perform a numerical Fourier computation. We sample the survey volume on a regular grid by assigning $`1`$ to the grid points inside the window function and $`0`$ outside. We then perform an FFT. This numerical computation is limited by the finite size of the grid cells, but avoids the overlapping zones and considers the true window function.
Figure 4 (filled points and solid line) compares both the analytical and numerical estimates of the window function power spectrum averaged over spherical shells. The difference is less than 5%.
The strongly anisotropic geometry of the ESP survey (see Figure 4) introduces important convolution effects between the survey window function and the galaxy distribution. To clean the observed power spectrum for these effects, we have adopted Lucy’s deconvolution method (Lucy 1974; see also Baugh & Efstathiou 1993 and Lin et al. 1996, for a discussion about its application to power spectrum estimates).
The Lucy technique is a general method to estimate the frequency distribution $`\psi (\eta )`$ of a quantity $`\eta `$, when we know the frequency distribution $`\varphi (y)`$ of a second quantity $`y`$, related to $`\eta `$ by
$$\varphi (y)=\psi (\eta )\mathrm{\Pi }(y|\eta )𝑑\eta ,$$
(28)
where $`\mathrm{\Pi }(y|\eta )dy`$ is the probability that $`y^{}[y,y+dy[`$ when $`\eta ^{}=\eta `$. The probability $`\mathrm{\Pi }(y|\eta )`$ must be known and the frequency distribution $`\varphi (y)`$ is the observed one.
The solution of equation 28 can be obtained by an iterative procedure. Let $`Q(\eta |y)d\eta `$ be the probability that $`\eta ^{}[\eta ,\eta +d\eta [`$ when $`y^{}=y`$. The probability that $`y^{}[y,y+dy[`$ and $`\eta ^{}[\eta ,\eta +d\eta [`$ can be written as $`\varphi (y)dyQ(\eta |y)d\eta `$ and $`\psi (\eta )d\eta \mathrm{\Pi }(y|\eta )dy`$. From these two expressions and equation 28 we obtain
$$Q(\eta |y)=\frac{\psi (\eta )\mathrm{\Pi }(y|\eta )}{\psi (\eta )\mathrm{\Pi }(y|\eta )𝑑\eta },$$
(29)
which provides the identity
$$\psi (\eta )\varphi (y)Q(\eta |y)𝑑y.$$
(30)
The latter equation cannot be solved directly, since $`Q(\eta |y)`$ depends on the unknown $`\psi (\eta )`$ as well. Given a fiducial model for $`\psi (\eta )`$ and the known probability $`\mathrm{\Pi }(y|\eta )`$, equation 29 provides an estimate for $`Q(\eta |y)`$. This and the identity 30 allows us to compute an improved estimate for $`\psi (\eta )`$. The process can then be repeated until convergency. In our specific case the equation to be solved is eq. 19, where $`k_{}^{}{}_{}{}^{2}\chi (k,k^{})`$ plays the role of the probability $`\mathrm{\Pi }(y|\eta )`$. If we sample $`k`$ on logarithmic intervals the convolution integral becomes
$$\stackrel{~}{P_c}(k)=k_{}^{}{}_{}{}^{3}P(k^{})\chi (k,k^{})\mathrm{ln}(10)d(\mathrm{log}_{10}k^{})$$
(31)
and an iterative scheme for the deconvolved spectrum can be written as
$$P^{m+1}(k_i)=P^m(k_i)\frac{\underset{j}{}\left[\stackrel{~}{P_c}(k_j)/\stackrel{~}{P_c}^m(k_j)\right]\chi (k_j,k_i)}{_j\chi (k_j,k_i)},$$
(32)
where
$$\stackrel{~}{P_c}^m(k_j)=\underset{r}{}k_{r}^{}{}_{}{}^{3}P^m(k_r)\chi (k_j,k_r)\mathrm{ln}(10)\mathrm{\Delta }$$
(33)
and $`\mathrm{\Delta }=(\mathrm{log}_{10}k_{i+1}\mathrm{log}_{10}k_i)`$ is the logarithmic interval, while $`P^m`$ denotes the $`m^{th}`$ estimate of the spectrum.
One problem with the Lucy method is that of producing a noiser and noiser solution as the iteration converges. To avoid this, Lucy suggests to stop the iteration after the first few steps. This is quite arbitrary and we prefer to follow Baugh & Efstathiou (1993; see also Lin et al. 1996) in applying a smoothing procedure at each step
$$P^m(k_i)=0.25P^m(k_{i1})+0.5P^m(k_i)+0.25P^m(k_{i+1}).$$
(34)
We use $`P^0(k_i)=constant`$ as initial guess for the power spectrum, but we checked that the solution is independent of the shape of $`P^0(k_i)`$. One consequence of this smoothing is that some degree of correlation is introduced among the bins of $`P(k)`$.
The importance of the convolution effects on different scales can be estimated by plotting the integrand of eq. (31) as a function of $`k^{}`$, for different values of $`k`$ (Figure 5).
If the window was a large and regular sample of the Universe, the plots would be sharply peaked at $`k=k^{}`$, as it actually happens for large values of $`k`$ (small scales).
On the other hand, for small values of $`k`$, i.e. for spatial wavelengths comparable to the typical scales of the window (which are quite small, due to the strongly anisotropic shape), the true power is spread over a wide range of wavenumbers.
## 4 Numerical Tests
We test the whole procedure for estimating the power spectrum through $`N`$-body simulations that we have run assuming some cosmological models (Carretti 1999). The results of these simulations can be considered as a Universe, from which we can extract mock catalogues with the same features of the ESP survey (geometry, galaxy density, field completeness, selection function). We then apply to such mock catalogues the whole power spectrum estimate procedure (convolved power spectrum estimator and deconvolution tecnique) and we compare the result with the true power spectrum obtained from the whole set of particles.
The simulations were performed on a Cray T3E at CINECA supercompunting center (Bologna) using a Particle-Mesh (PM) code (Carretti & Messina 1999) and adopting two cosmological models: an unbiased Standard Cold Dark Matter (SCDM: $`\mathrm{\Omega }_\mathrm{o}=1`$, $`h=0.5`$, $`\sigma _8=1`$) and an unbiased Open Cold Dark Matter with shape parameter $`\mathrm{\Gamma }=0.2`$ (OCDM: $`\mathrm{\Omega }_\mathrm{o}=0.4`$, $`h=0.5`$, $`\sigma _8=1`$). They were run with a box size of $`700h^1`$ Mpc, $`512^3`$ grid points and $`512^3`$ particles, in order to reproduce a volume which can contain all catalogues selected for the analysis (max depth $`633h^1`$ Mpc) and to select a realistic number of mock galaxies for the magnitude–limited catalogues. From each simulation box, we randomly choose a particle as origin and extract sets of particles with the same features of the three ESP catalogues. The magnitude–limited selection is then reproduced by simply assigning a weight corresponding to the observed selection function.
From each simulation and for each ESP subsample we construct 50 independent mock catalogues. The average number of particles over the 50 realisations is set to the number of galaxies observed in the corresponding true ESP sample. The power spectrum estimator is then applied to each of the 50 mock galaxy catalogues, producing an independent estimate of $`P(k)`$ for that specific model and sample geometry. From each set of realisations, a mean $`P(k)`$ and its standard deviation can finally be computed and compared to the true power spectrum obtained from all the particles.
A general result from this exercise is that for $`k>0.065h`$ Mpc<sup>-1</sup> the systematic power suppression by the window function convolution is properly corrected for by our procedure, i.e. we are able to fully recover the input $`P(k)`$. In figure 6 we show the reconstructed power spectra from the three sets of mock catalogues, compared to the “true” ones, both for the SCDM and OCDM simulations. In panel a), in particular, we also show (open circles) the raw power spectrum before applying the Lucy deconvolution, to emphasize the dramatic effect of the ESP window function on all scales. It is evident that for $`k>0.065h`$ Mpc<sup>-1</sup> the mean deconvolved power spectra are a very good reconstruction of the original ones. In particular, in the SCDM case where the spectrum turnover scale is well sampled by the simulation box, the technique is able to nicely follow the change of shape at small $`k`$’s. This is important, because guarantees that the deconvolution method has enough resolution as to follow possible features in the data power spectrum, while at the same time recovering the correct amplitude. At smaller $`k`$’s the error bars explode, and the results become meaningless. In the case of the deepest sample, ESPm633, the reconstruction shows a small systematic overestimate of the amplitude for the SCDM spectrum on the largest scales, i.e. the reconstruction algorithm seems to have difficulty in following the curvature of the spectrum accurately. This is probably due to the rather small value of the selection function in the most distant part of the sample, which puts too large a weight on the distant objects. Rather than indicating a difficulty in the technique, this is probably telling us that it is safer to truncate the data at smaller distances, as for sample ESPm523.
Comparing the results from the SCDM and OCDM mock catalogues, we have checked that the fractional errors for the two cases are quite similar. Not knowing a priori the correct cosmological model, rather than choosing one of the two models as representative, we prefer to average the fractional errors measured from the two models.
Using the mock catalogues, we can also evaluate directly the possible effects of the field incompleteness on the power spectrum estimate. Figure 7 compares the mean power spectra obtained from one set of 50 SCDM mock samples both in the ideal case (all fields complete) and when the ESP field–to–field incompleteness is introduced and corrected for. It is clear that the incompleteness is correctly taken into account by the weighting scheme. Error bars (not reported for clarity) are also very similar.
## 5 The Power Spectrum of ESP galaxies
The numerical tests performed have given us an estimate of the reliability of our method to reconstruct the true power spectrum, so that we can now apply it to the three galaxy subsamples ESPm523, ESP523 and ESPm633.
The final results of the computation are shown in figure 8 and Table 2.
The error bars are partially reported only for ESPm523 to avoid confusion (the errors are similar for the three samples). The three estimates of the power spectrum are well consistent with each other. Given the large amplitude of the errors ($`30`$% of $`P(k)`$ for $`k>0.15h`$ Mpc<sup>-1</sup>, 50% for $`k0.1`$ and 75% for $`k0.065`$) the small differences in the slopes are not significant. In general, we can safely say that the power spectrum of ESP galaxies follows a power law $`P(k)k^n`$ with $`n2.2`$ for $`k>0.2h`$ Mpc<sup>-1</sup>, and $`n1.6`$ for $`k<0.2h`$ Mpc<sup>-1</sup>. In the range $`0.065<k<0.6h`$ Mpc<sup>-1</sup> there is no meaningful difference between the three estimates, which are therefore independent of the catalogue type (magnitude– or volume–limited) and of the catalogue depth.
In figure 8 we have also plot, for comparison, redshift– and real–space power spectra computed from the simulations described in section 4 (SCDM and $`\mathrm{\Gamma }=0.2`$ OCDM). Note how the former compensate for non linear evolution at large $`k`$’s in this way steepening the slope of $`P(k)`$ over the whole observed range, which makes the global slope closer to the observed one. Despite this comparison to models is deliberately limited, one can safely say that the data points (especially below $`k=0.1`$$`0.2h`$ Mpc<sup>-1</sup>) are in better agreement with the power spectrum of the $`\mathrm{\Gamma }=0.2`$ OCDM model. This model would reproduce this observation without biasing (the normalisation adopted by the simulation is $`\sigma _8=1`$).
## 6 Consistency between Real and Fourier Space
It is interesting to compare the Fourier transform of the ESP power spectrum estimated in this work with the two–point correlation function measured independently from the same sample (Guzzo et al. 2000). This exercise is a further check of the robustness and self–consistency of the estimate of $`P(k)`$. In addition, it is of specific interest to verify the effect of the survey geometry/window function in real and Fourier space. To simplify the procedure, we have first fitted the observed $`P(k)`$ with a simple analytical form with two power laws connected by a smooth turnover (e.g. Peacock 1997)
$$P(k)=\frac{(k/k_o)^\alpha }{1+\left(k/k_c\right)^{\alpha n}},$$
(35)
where $`k_o`$ is a normalisation factor, $`k_c`$ gives essentially the turnover scale, $`n`$ is the large–scale primordial index (here fixed to $`n=1`$), and $`\alpha `$ gives the slope for $`kk_c`$. We have used this function to reproduce the global shape of both the convolved and de–convolved estimates of $`P(k)`$ from the ESP523 sample. In terms of selection function, this sample is the closest to one of the volume–limited samples used by Guzzo et al. (2000) to estimate $`\xi (s)`$ from the same data. Figure 9 shows how this form provides a good description of the ESP power spectrum, with the deconvolved one characterised by $`k_o=0.080h`$ Mpc<sup>-1</sup>, $`k_c=0.062h`$ Mpc<sup>-1</sup>, $`\alpha =2.2`$ (note that while the slope $`\alpha `$ is a stable value, the turnover scale $`k_c`$ is very poorly constrained, given the limited range covered by the data).
Figure 10 shows the Fourier transform of the two fits, compared to the direct estimate of $`\xi (s)`$ by Guzzo et al. (2000). Two main comments should be made here. First, our ”best” estimate of $`P(k)`$, deconvolved for the ESP window function according to our recipe, reproduces rather well the observed two–point correlation function (solid line). Note how the Fourier transform of the simple direct estimate suffers from a systematic lack of power as a function of scale (dashed line), as we expected from our results on the mock samples. The second, more general comment concerns the stability of the two–point correlation function. One might naively think that the narrowness of the explored volume, which gives rise to the window function in Fourier space, should affect in a similar way also the estimate of clustering by the two–point correlation function. Figure 10 shows that this is not the case. In fact, the points showed here have not been subject to any kind of correction (Guzzo et al. 2000), a part from those which are standard in the estimation technique to take into account the survey boundaries. Still, they seem to sample clustering to the largest available scales in a reasonably unbiased way, without basically being affected by the survey geometry.
## 7 Comparison to Other Redshift Surveys
In the six panels of figure 11, we compare the power spectrum for the ESPm523 and ESPm633 samples to a variety of results from previous surveys, both selected in the optical and infrared (IRAS) bands. In general, there is a good level of unanimity among the different surveys concerning the slope of $`P(k)`$ over the range sampled by the ESP estimate. Optically–selected surveys show a good agreement also in amplitude, with a possible minor differential biasing effect in the case of CfA2–SSRS2–130 (panel a, da Costa et al. 1994), which is a volume–limited sample containing galaxies brighter than $`M^{}1.5`$. The effect of different biasing values is more evident in the comparison to IRAS–based surveys (IRAS 1.2 Jy, Fisher et al. 1993; QDOT, Feldman et al. 1994; PSCz, Sutherland et al. 1999) in panels e and f.
Particularly relevant is the comparison to the results of the Durham/UKST galaxy redshift survey (DUKST, Hoyle et al., 1999) (panel d). This survey is selected from the same parent photometric catalogue as the ESP (the EDSGC) and contains a comparable number of redshifts. However, it is less deep (bj 17), while covering a much larger solid angle by measuring redshift in a sparse–sampling fashion, picking one galaxy in three. This produces a window function which is essentially complementary to that of the ESP survey, with a good sampling of long wavelengths and a poor description of small–scale clustering, which on the contrary is well sampled by the ESP 1-in-1 redshift measurements. The agreement between these two data sets is impressive. This is a further confirmation of the quality of the deconvolution procedure we have applied to the ESP data, given the rather 3D shape of the DUKST volume which makes the window function practically negligible for this survey. Significantly more noisy is the estimate from the similarly $`b_J`$–selected Stromlo-APM redshift survey (Tadros & Efstathiou 1996), most probably because of the very sparse sampling of this survey and the smaller number of galaxies.
Finally, panel b) shows a comparison with the data from the $`r`$–selected LCRS (Lin et al. 1996). The power spectrum from this survey has a flatter slope with respect to our estimate from the ESP. More in general, it is flatter than practically all other power spectra shown in the figure. This is somewhat suspicious, as the two–point correlation functions agree rather well for ESP, LCRS, Stromlo-APM and DUKST (Guzzo 1999), and might be an indication that the effect of the window function has not been fully removed from the estimated spectrum.
## 8 Summary and Conclusions
The main results obtained in this work can be summarised as follows.
* We have developed a technique to properly describe the ESP window function analytically, and then deconvolve it from the measured power spectrum, to obtain an estimate of the galaxy power spectrum. The tests performed on a number of mock catalogues drawn from large $`N`$–body simulations show that the technique is able to recover the correct shape of $`P(k)`$ down to wavenumbers $`k0.065h`$ Mpc<sup>-1</sup>. In general, this technique for describing the window function analytically can be applied to any redshift survey composed by circular patches on the sky (e.g. the ongoing 2dF survey). In addition to its mathematical elegance, it has some computational advantages over the traditional method for recovering the survey window function, normally based on the generation of large Montecarlo poissonian realisations.
* The final estimates of the ESP power spectrum, extracted from three subsamples of the survey, are in good agreement within the error bars. The bright volume–limited sample does not show a clear difference in amplitude with respect to the apparent–magnitude limited ones. This agrees with the similar behaviour found for the two–point correlation function, i.e. a negligible evidence for luminosity segregation even for limiting absolute magnitudes $`M_{b_J}20`$ (Guzzo et al. 2000). This is only apparently in contrast with the results of Park et al. (1994), who found evidence for luminosity segregation studying the amplitude of the power spectrum in the CfA2 survey. In fact, that analysis concentrates on a range of luminosities about 1.5 magnitude brighter than $`M^{}`$, which for the CfA2 survey has a value of -18.8 (Marzke et al. 1994), i.e. nearly one magnitude fainter than for the ESP. This also agrees with the results of Benoist et al. (1996), who studied the correlation function for the SSRS2 sample, finding negligible signs of luminosity segregation for $`M>M^{}`$.
* All three estimates of $`P(k)`$ show a similar shape, with a well defined power–law $`k^n`$ with $`n2.2`$ for $`k0.2h`$ Mpc<sup>-1</sup>, and a smooth bend to a flatter shape ($`n1.6`$) for smaller $`k`$’s. The smallest wavenumber where a meaningful reconstruction can be performed ($`k0.065h`$ Mpc<sup>-1</sup>), does not allow us to explore the range of scales where other power spectra seem to show a flattening and hints for a turnover. In the framework of CDM models, however, the well–sampled steep slope between 0.08 and 0.3 $`h`$ Mpc<sup>-1</sup> favours a low–$`\mathrm{\Gamma }`$ model ($`\mathrm{\Gamma }=0.2`$), consistently with the most recent CMB observation of BOOMERANG/MAXIMA experiments (Jaffe et al. 2000).
* We have verified that the two–point correlation function $`\xi (s)`$ is much less sensitive to the effect of a difficult window function as that of the ESP, than the power spectrum. In fact, the measured correlation function (without any correction), agrees with the Fourier transform of the power spectrum, only after this has been cleaned of the combination by the window function. This is an instructive example of how these two quantities, despite being mathematically equivalent, can be significantly different in their practical estimates and be very differently affected by the peculiarities of data samples.
* When compared to previous estimates from other surveys, the ESP power spectrum is virtually indistinguishable from that of the Durham-UKST survey over the common range of wavenumbers. In particular, between 0.1 and 1 $`h`$ Mpc<sup>-1</sup> our power spectrum has significantly smaller error bars with respect to the DUKST, by virtue of its superior small–scale sampling. The absence of any systematic amplitude difference between these two surveys – both selected from the EDSGC catalogue, but with complementary volume and sampling choices – is an important indirect indication of the quality of the deconvolution procedure applied here, and also of the accuracy of the two independent estimates. In this respect, a combination of the Durham-UKST and ESP surveys possibly provides the current best measure of $`P(k)`$ for blue–selected galaxies in the full range $`0.031h`$ Mpc<sup>-1</sup>. It will be very interesting to compare these combined results to the power spectrum of the forthcoming 2dF redshift survey, which is also selected in the same $`b_J`$ band to virtually the same limiting magnitude than the ESP.
## Acknowledgments
We thank an anonymous referee for suggestions that helped us to improve the paper. We thank Hume Feldman, Huan Lin, and Michael Vogeley for providing us with their power spectrum results in electronic form, and Stefano Borgani for his COBE normalisation routine. LG and EZ thank all their collaborators in the ESP survey team, for their contribution to the success of the survey. This work has been partially supported by a CNAA grant.
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# References
Magnetised tori were considered long ago by Witten in the first attempt to recover four-dimensional chiral spectra from the low-energy field theory of superstrings. More recently, Bachas analysed the effect of Fradkin-Tseytlin deformations on open strings, and showed how their universal magnetic couplings can lead to chiral spectra with broken supersymmetry. However, these models have in general Nielsen-Olesen instabilities , that reflect themselves in the emergence of tachyonic modes. This complicates the analysis, and brings about some surprises. For instance, in some cases with extended ($`𝒩=2,4`$) supersymmetry, where one can analyse the potentials of the tachyonic modes, at the resulting minima supersymmetry is actually restored . The constructions in were both based on the assumption, natural at the time, of a vanishing instanton density for the internal magnetic field. However, we are now accustomed to more general settings, that have naturally emerged from type-I vacua , where a non-vanishing instanton density is compensated by the presence of additional branes . This letter is thus devoted to elucidate some peculiar effects of magnetic deformations with non-vanishing instanton number on toroidal and orbifold compactifications of type-I strings. As we shall see, these can result in new vacua with unbroken supersymmetry and Chan-Paton groups of reduced rank, where magnetised D9 branes effectively mimic BPS D5 (anti)branes.
It is by now well appreciated that, in non-trivial gravitational and gauge backgrounds, the Wess-Zumino coupling of endows D branes with R-R charges for forms of different degrees. It is perhaps less appreciated, however, that the Born-Infeld action can turn the non-vanishing vacuum energy of suitable internal magnetic fields into a positive tension capable of recovering the BPS bound for the additional charges. An indirect manifestation of this phenomenon was recently met in , where the open descendants of some asymmetric orbifolds with “brane supersymmetry breaking” were built using a magnetised internal space, and where a suitable choice of internal fields played an essential role in saturating all R-R tadpoles with only D9 branes.
Let us begin with some intuitive field theory arguments, well captured by the low-energy effective action for D9 branes in an internal abelian background <sup>5</sup><sup>5</sup>5The (dimensionless) magnetic fields used in this letter differ from the conventional ones by a $`2\pi \alpha ^{}`$ rescaling.,
$$𝒮_9=T_{(9)}_{_{10}}\mathrm{d}^{10}xe^\varphi \underset{a=1}{\overset{32}{}}\sqrt{det\left(g_{10}+q_aF\right)}\mu _{(9)}\underset{p,a}{}_{_{10}}e^{q_aF}C_{p+1}+\mathrm{},$$
(1)
where $`a`$ labels the types of Chan-Paton charges that couple to the magnetic fields with strength $`q_a`$,
$$T_{(p)}=\sqrt{\frac{\pi }{2\kappa ^2}}\left(2\pi \sqrt{\alpha ^{}}\right)^{3p}=|\mu _{(p)}|,$$
(2)
with $`T`$ and $`\mu `$ the tension and the R-R charge for a type-I D$`p`$ brane , and where $`\kappa `$ defines the ten-dimensional Newton constant $`G_N^{(10)}=\kappa ^2/8\pi `$. To illustrate the phenomenon, anticipating the string construction, it suffices to consider the geometry $`_{10}=_6\times T^2\times T^2`$ with constant abelian magnetic fields $`H_1`$ and $`H_2`$ lying in the two internal tori. These are effectively monopole fields, and thus satisfy the Dirac quantisation conditions
$$qH_iv_i=k_i(i=1,2),$$
(3)
where, aside from powers of $`2\pi `$, $`v_i=R_i^{(1)}R_i^{(2)}/\alpha ^{}`$ are the dimensionless volumes of the two tori of radii $`R_i^{(1)}`$ and $`R_i^{(2)}`$, $`k_i`$ are the degeneracies of the corresponding Landau levels and $`q`$ is the elementary electric charge for the system. As anticipated, we forego the restriction in and actually pick a pair of abelian fields aligned with the same U(1) subgroup, so that
$`𝒮_9`$ $`=`$ $`T_{(9)}{\displaystyle _{_{10}}}\mathrm{d}^{10}xe^\varphi \sqrt{g_6}{\displaystyle \underset{a=1}{\overset{32}{}}}\sqrt{(1+q_a^2H_1^2)(1+q_a^2H_2^2)}`$ (4)
$`\mathrm{\hspace{0.33em}32}\mu _{(9)}{\displaystyle _{_{10}}}C_{10}\left(2\pi \sqrt{\alpha ^{}}\right)^4\mu _{(9)}v_1v_2H_1H_2{\displaystyle \underset{a=1}{\overset{32}{}}}q_a^2{\displaystyle __6}C_6,`$
where $`g_6`$ denotes the six-dimensional space-time metric, and for simplicity we have chosen an identity metric in the internal space. In particular, if the two internal fields have identical magnitudes, for the resulting (anti)self-dual configuration the action becomes
$`𝒮_9`$ $`=`$ $`\mathrm{\hspace{0.33em}32}{\displaystyle _{_{10}}}\left(\mathrm{d}^{10}x\sqrt{g_6}T_{(9)}e^\varphi +\mu _{(9)}C_{10}\right)`$ (5)
$`{\displaystyle \underset{a=1}{\overset{32}{}}}\left({\displaystyle \frac{q_a}{q}}\right)^2{\displaystyle __6}\left(\mathrm{d}^6x\sqrt{g_6}|k_1k_2|T_{(5)}e^\varphi +k_1k_2\mu _{(5)}C_6\right).`$
Notice that the Dirac quantisation conditions (3) have compensated the integration over the internal tori, while in the second line the additional powers of $`\alpha ^{}`$ have nicely converted $`T_{(9)}`$ and $`\mu _{(9)}`$ into $`T_{(5)}`$ and $`\mu _{(5)}`$. Thus, a D9 brane on a magnetised $`T^2\times T^2`$ indeed mimics a D5 brane or a D5 antibrane according to whether the orientations of $`H_1`$ and $`H_2`$, reflected by the relative sign of $`k_1`$ and $`k_2`$, are identical or opposite.
We can now turn to the open-string description of this phenomenon. In order to obtain a supersymmetric configuration, we should start from an orbifold that normally requires the introduction of D5 branes. The simplest such instance is the six-dimensional compactification on $`(T^2\times T^2)/Z_2`$ with Klein-bottle projection
$$𝒦=\frac{1}{4}\left\{(Q_o+Q_v)(0;0)\left[P_1P_2+W_1W_2\right]+16\times 2(Q_s+Q_c)(0;0)\left(\frac{\eta }{\vartheta _4(0)}\right)^2\right\},$$
(6)
that corresponds to the introduction of $`\mathrm{O9}_+`$ and $`\mathrm{O5}_+`$ planes, and thus to a projected $`𝒩=(1,0)`$ supersymmetric closed spectrum with one tensor multiplet and 20 hypermultiplets.
In writing this expression, we have endowed the six-dimensional characters of with a pair of arguments, anticipating the effect of the magnetic deformations in the two internal tori. In general
$`Q_o(\eta ;\zeta )`$ $`=`$ $`V_4(0)\left[O_2(\eta )O_2(\zeta )+V_2(\eta )V_2(\zeta )\right]C_4(0)\left[S_2(\eta )C_2(\zeta )+C_2(\eta )S_2(\zeta )\right],`$
$`Q_v(\eta ;\zeta )`$ $`=`$ $`O_4(0)\left[V_2(\eta )O_2(\zeta )+O_2(\eta )V_2(\zeta )\right]S_4(0)\left[S_2(\eta )S_2(\zeta )+C_2(\eta )C_2(\zeta )\right],`$
$`Q_s(\eta ;\zeta )`$ $`=`$ $`O_4(0)\left[S_2(\eta )C_2(\zeta )+C_2(\eta )S_2(\zeta )\right]S_4(0)\left[O_2(\eta )O_2(\zeta )+V_2(\eta )V_2(\zeta )\right],`$
$`Q_c(\eta ;\zeta )`$ $`=`$ $`V_4(0)\left[S_2(\eta )S_2(\zeta )+C_2(\eta )C_2(\zeta )\right]C_4(0)\left[V_2(\eta )O_2(\zeta )+O_2(\eta )V_2(\zeta )\right],`$ (7)
where the four level-one O($`2n`$) characters are related to the four Jacobi theta functions according to
$`O_{2n}(\zeta )`$ $`=`$ $`{\displaystyle \frac{1}{2\eta ^n(\tau )}}\left(\vartheta _3^n(\zeta |\tau )+\vartheta _4^n(\zeta |\tau )\right),S_{2n}(\zeta )={\displaystyle \frac{1}{2\eta ^n(\tau )}}\left(\vartheta _2^n(\zeta |\tau )+i^n\vartheta _1^n(\zeta |\tau )\right),`$
$`V_{2n}(\zeta )`$ $`=`$ $`{\displaystyle \frac{1}{2\eta ^n(\tau )}}\left(\vartheta _3^n(\zeta |\tau )\vartheta _4^n(\zeta |\tau )\right),C_{2n}(\zeta )={\displaystyle \frac{1}{2\eta ^n(\tau )}}\left(\vartheta _2^n(\zeta |\tau )i^n\vartheta _1^n(\zeta |\tau )\right).`$ (8)
Whereas in the internal magnetic two-forms were chosen to satisfy
$$\mathrm{tr}H_iH_j=0,$$
(9)
here we allow for a non-vanishing instanton density, that in String Theory is naturally compensated by additional unpaired defects (an excess of D5 (anti)branes and/or O5 planes). In particular, as in our field theory considerations, we take the two internal fields aligned with the same U(1) subgroup of SO(32), a choice that in this $`Z_2`$ orbifold can preserve at most a $`\mathrm{U}(m)\times \mathrm{U}(n)`$ gauge group, with $`m+n=16`$. In the following, we actually restrict our attention to this maximal case, from which other examples can be obtained via Wilson lines or brane displacements.
In writing the direct-channel annulus amplitude, let us begin by recalling that a uniform magnetic field with components $`H_1`$ and $`H_2`$ in the two internal tori alters the boundary conditions for open strings, shifting their mode frequencies by
$$z_i^{\mathrm{L},\mathrm{R}}=\frac{1}{\pi }\left[\mathrm{tan}^1(q_\mathrm{L}H_i)+\mathrm{tan}^1(q_\mathrm{R}H_i)\right],$$
(10)
where $`q_\mathrm{L}`$ ($`q_\mathrm{R}`$) denote the charges of the left (right) end of the open string with respect to the U(1) fields $`H_i`$. A further novelty is displayed by “dipole” strings, with opposite end charges, whose oscillator modes are unaffected, but whose world-sheet coordinates undergo a complex “boost”, so that their Kaluza-Klein momenta $`m_i`$ are rescaled according to
$$m_i\frac{m_i}{\sqrt{1+q_a^2H_i^2}}.$$
(11)
This rescaling ensures the consistency of the transverse-channel amplitudes, whose lowest-level contributions, aside from a subtlety that we shall discuss later, are to group as usual into perfect squares.
The techniques of determine the direct-channel annulus amplitude
$`𝒜`$ $`=`$ $`\frac{1}{4}\{(Q_o+Q_v)(0;0)[(m+\overline{m})^2P_1P_2+(d+\overline{d})^2W_1W_2+2n\overline{n}\stackrel{~}{P}_1\stackrel{~}{P}_2]`$
$``$ $`2(m+\overline{m})(n+\overline{n})(Q_o+Q_v)(z_1\tau ;z_2\tau ){\displaystyle \frac{k_1\eta }{\vartheta _1(z_1\tau )}}{\displaystyle \frac{k_2\eta }{\vartheta _1(z_2\tau )}}`$
$``$ $`(n^2+\overline{n}^2)(Q_o+Q_v)(2z_1\tau ;2z_2\tau ){\displaystyle \frac{2k_1\eta }{\vartheta _1(2z_1\tau )}}{\displaystyle \frac{2k_2\eta }{\vartheta _1(2z_2\tau )}}`$
$``$ $`\left[(m\overline{m})^22n\overline{n}+(d\overline{d})^2\right](Q_oQ_v)(0;0)\left({\displaystyle \frac{2\eta }{\vartheta _2(0)}}\right)^2`$
$``$ $`2(m\overline{m})(n\overline{n})(Q_oQ_v)(z_1\tau ;z_2\tau ){\displaystyle \frac{2\eta }{\vartheta _2(z_1\tau )}}{\displaystyle \frac{2\eta }{\vartheta _2(z_2\tau )}}`$
$``$ $`(n^2+\overline{n}^2)(Q_oQ_v)(2z_1\tau ;2z_2\tau ){\displaystyle \frac{2\eta }{\vartheta _2(2z_1\tau )}}{\displaystyle \frac{2\eta }{\vartheta _2(2z_2\tau )}}`$
$`+`$ $`2(m+\overline{m})(d+\overline{d})(Q_s+Q_c)(0;0)\left({\displaystyle \frac{\eta }{\vartheta _4(0)}}\right)^2`$
$`+`$ $`2(d+\overline{d})(n+\overline{n})(Q_s+Q_c)(z_1\tau ;z_2\tau ){\displaystyle \frac{\eta }{\vartheta _4(z_1\tau )}}{\displaystyle \frac{\eta }{\vartheta _4(z_2\tau )}}`$
$``$ $`2(m\overline{m})(d\overline{d})(Q_sQ_c)(0;0)\left({\displaystyle \frac{\eta }{\vartheta _3(0)}}\right)^2`$
$``$ $`2(d\overline{d})(n\overline{n})(Q_sQ_c)(z_1\tau ;z_2\tau ){\displaystyle \frac{\eta }{\vartheta _3(z_1\tau )}}{\displaystyle \frac{\eta }{\vartheta _3(z_2\tau )}}\},`$
and the corresponding Möbius amplitude
$``$ $`=`$ $`\frac{1}{4}\{(\widehat{Q}_o+\widehat{Q}_v)(0;0)[(m+\overline{m})P_1P_2+(d+\overline{d})W_1W_2]`$
$``$ $`(n+\overline{n})(\widehat{Q}_o+\widehat{Q}_v)(2z_1\tau ;2z_2\tau ){\displaystyle \frac{2k_1\widehat{\eta }}{\widehat{\vartheta }_1(2z_1\tau )}}{\displaystyle \frac{2k_2\widehat{\eta }}{\widehat{\vartheta }_1(2z_2\tau )}}`$
$``$ $`\left(m+\overline{m}+d+\overline{d}\right)(\widehat{Q}_o\widehat{Q}_v)(0;0)\left({\displaystyle \frac{2\widehat{\eta }}{\widehat{\vartheta }_2(0)}}\right)^2`$
$``$ $`(n+\overline{n})(\widehat{Q}_o\widehat{Q}_v)(2z_1\tau ;2z_2\tau ){\displaystyle \frac{2\widehat{\eta }}{\widehat{\vartheta }_2(2z_1\tau )}}{\displaystyle \frac{2\widehat{\eta }}{\widehat{\vartheta }_2(2z_2\tau )}}\}.`$
Here we have actually resorted to a shorthand notation, where the arguments $`z_i`$ ($`2z_i`$) are associated to strings with one (two) charged ends. Moreover, both the imaginary modulus $`\frac{1}{2}it`$ of $`𝒜`$ and the complex modulus $`\frac{1}{2}+\frac{1}{2}it`$ of $``$ are denoted by the same symbol $`\tau `$, although the proper “hatted” contributions to the Möbius amplitude are explicitly indicated. $`P_i`$ and $`W_i`$ are conventional momentum and winding sums for the two-tori, while a “tilde” denotes a sum with momenta “boosted” as in (11). Finally, $`d`$ (together with its conjugate $`\overline{d}`$) is the Chan-Paton multiplicity for the D5 branes, while $`m`$ and $`n`$ (together with their conjugates $`\overline{m}`$ and $`\overline{n}`$) are Chan-Paton multiplicities for the D9 branes. For the sake of brevity, several terms with opposite U(1) charges, and thus with opposite $`z_i`$ arguments, have been grouped together, using the symmetries of the Jacobi theta-functions.
For generic magnetic fields, the open spectrum is indeed non-supersymmetric and develops Nielsen-Olesen instabilities . As emphasised in , the emergence of these tachyonic modes can be ascribed to the magnetic couplings of the internal components of gauge fields. For instance, small magnetic fields affect the mass formula for the untwisted string modes according to
$$\mathrm{\Delta }M^2=\frac{1}{2\pi \alpha ^{}}\underset{i=1,2}{}\left[(2n_i+1)|(q_\mathrm{L}+q_\mathrm{R})H_i|+2(q_\mathrm{L}+q_\mathrm{R})\mathrm{\Sigma }_iH_i\right],$$
(14)
where the first term originates from the Landau levels and the second from the magnetic moments of the spins $`\mathrm{\Sigma }_i`$. For the internal components of the vectors, the magnetic moment coupling generally overrides the zero-point contribution, leading to tachyonic modes, unless $`|H_1|=|H_2|`$, while for spin-$`\frac{1}{2}`$ modes it can at most compensate it. On the other hand, for twisted modes the zero-point contribution is absent, since ND strings have no Landau levels. In this case the low-lying space-time fermions, that originate from the fermionic part $`S_4O_4`$ of $`Q_s`$, are scalars in the internal space and have no magnetic moment couplings. However, their bosonic partners, that originate from $`O_4C_4`$, are affected by the magnetic deformations and have mass shifts $`\mathrm{\Delta }M^2\pm (H_1H_2)`$. Therefore, if $`H_1=H_2`$ all tachyonic instabilities are indeed absent. Actually, with this choice the supersymmetry charge, that belongs to $`C_4C_4`$, is also unaffected<sup>6</sup><sup>6</sup>6Type-II branes at angles preserving some supersymmetry were previously considered in . After T-dualities, these can be related to special choices for the internal magnetic fields. Type I toroidal models, however, can not lead to supersymmetric configurations, since the resulting R-R tadpoles require the introduction of antibranes.. Therefore, a residual supersymmetry is present for the entire string spectrum, and indeed, using the Jacobi identities for non-vanishing arguments , one can see that for $`z_1=z_2`$ both $`𝒜`$ and $``$ vanish identically. Still, the resulting supersymmetric models are rather peculiar, as can be seen from the deformed tadpole conditions, to which we now turn.
Let us begin by examining the untwisted R-R tadpole conditions. For $`C_4S_2C_2`$ one finds
$$\left[m+\overline{m}+n+\overline{n}32+q^2H_1H_2(n+\overline{n})\right]\sqrt{v_1v_2}+\frac{1}{\sqrt{v_1v_2}}\left[d+\overline{d}32\right]=0,$$
(15)
aside from terms that vanish after identifying the multiplicities of conjugate representations $`(m,\overline{m})`$, $`(n,\overline{n})`$ and $`(d,\overline{d})`$. The additional (untwisted) R-R tadpole conditions from $`Q_o`$ and $`Q_v`$ are compatible with (15) and do not add further constraints. This expression reflects the familiar Wess-Zumino coupling of eq. (1), and therefore the various powers of $`H`$ correspond to R-R forms of different degrees. In particular, as we anticipated in our field theory discussion, the term bilinear in the magnetic fields has a very neat effect: it charges the D9 brane with respect to the six-form potential. This can be seen very clearly making use of the quantisation condition (3), that turns the tadpole conditions (15) into
$`m+\overline{m}+n+\overline{n}`$ $`=`$ $`32,`$
$`k_1k_2(n+\overline{n})+d+\overline{d}`$ $`=`$ $`32.`$ (16)
Thus, if $`k_1k_2>0`$ the D9 branes indeed acquire the R-R charge of $`|k_1k_2|`$ D5 branes, while if $`k_1k_2<0`$ they acquire the R-R charge of as many D5 antibranes, in agreement with eq. (5).
The untwisted NS-NS tadpoles exhibit very nicely their relation to the Born-Infeld term in (1). For instance, the dilaton tadpole
$$\left[m+\overline{m}+(n+\overline{n})\sqrt{\left(1+q^2H_1^2\right)\left(1+q^2H_2^2\right)}32\right]\sqrt{v_1v_2}+\frac{1}{\sqrt{v_1v_2}}\left[d+\overline{d}32\right]$$
(17)
originates from $`V_4O_2O_2`$, and can be clearly linked to the derivative of the integrand of $`𝒮_9`$, specialised to the form (4), with respect to $`\varphi `$. On the other hand, the volume of the first internal torus originates from $`O_4V_2O_2`$, and the corresponding tadpole,
$$\left[m+\overline{m}+(n+\overline{n})\frac{1q^2H_1^2}{\sqrt{1+q^2H_1^2}}\sqrt{1+q^2H_2^2}32\right]\sqrt{v_1v_2}\frac{1}{\sqrt{v_1v_2}}\left[d+\overline{d}32\right],$$
(18)
can be linked to the derivative of the Born-Infeld action in (1) with respect to the corresponding breathing mode. A similar result holds for the volume of the second torus, with the proper interchange of $`H_1`$ and $`H_2`$. For the sake of brevity, we have omitted in these NS-NS tadpoles all terms that vanish using the constraint $`n=\overline{n}`$. However, the full expression of (18) is rather interesting, since, in contrast with the usual structure of unoriented string amplitudes, it is not a perfect square. This unusual feature can be ascribed to the behaviour of the internal magnetic fields under time reversal. Indeed, as stressed long ago in , these transverse-channel amplitudes involve a time-reversal operation $`𝒯`$, and are thus of the form $`𝒯(B)|q^{L_0}|B`$. Differently from the usual quantum mechanical amplitudes, this type of expression is generally a bilinear, rather than a sesquilinear, form. This, however, is not true in the present examples, where additional signs are introduced by the magnetic fields, that are odd under time reversal. As a result, in deriving from factorisation the Möbius amplitudes for these models, it is crucial to add the two contributions $`𝒯(B)|q^{L_0}|C`$ and $`𝒯(C)|q^{L_0}|B`$, that are different and effectively eliminate the additional terms from the transverse-channel.
Both (18) and the dilaton tadpole (17) simplify drastically in the interesting case $`H_1=H_2`$ where, using the Dirac quantisation conditions (3), they become
$$\left[m+\overline{m}+n+\overline{n}32\right]\sqrt{v_1v_2}\frac{1}{\sqrt{v_1v_2}}\left[k_1k_2(n+\overline{n})+d+\overline{d}32\right].$$
(19)
Thus, they both vanish, as they should, in these supersymmetric configurations, once the corresponding R-R tadpole conditions (16) are enforced.
The twisted R-R tadpoles
$$15\left[\frac{1}{4}(m\overline{m}+n\overline{n})\right]^2+\left[\frac{1}{4}(m\overline{m}+n\overline{n})(d\overline{d})\right]^2$$
(20)
originate from the sector $`S_4O_2O_2`$, whose states are scalars in the internal space. They reflect very neatly the distribution of branes among the sixteen fixed points, only one of which accommodates D5 branes in our examples, are not affected by the magnetic fields, and vanish identically for the given unitary gauge groups. In general these breaking terms, that originate from twisted modes flowing in the transverse channel, can be linked to internal curvature contributions to the Wess-Zumino term, here localised at the fixed points: this is actually the reason for the presence of D9 and D5 terms in the same expression in orbifold models. The corresponding NS-NS tadpoles, originating from the $`O_4S_2C_2`$ and $`O_4C_2S_2`$ sectors, are somewhat more involved, and after the identification of conjugate multiplicities are proportional to
$$\frac{q(H_1H_2)}{\sqrt{(1+q^2H_1^2)(1+q^2H_2^2)}}.$$
(21)
They clearly display new couplings for twisted NS-NS fields that, to the best of our knowledge, were not previously exhibited. Notice that, as expected, for $`H_1=H_2`$ these twisted tadpoles also vanish.
We can now describe some supersymmetric models corresponding to the special choice $`H_1=H_2`$. It suffices to confine our attention to the case $`k_1=k_2=2`$, the minimal Landau-level degeneracies allowed on this $`Z_2`$ orbifold. Although the projected closed spectra of all the resulting models are identical, and comprise the $`𝒩=(1,0)`$ gravitational multiplet, together with one tensor multiplet and twenty hypermultiplets, the corresponding open spectra are quite different from the standard ones, with a maximal gauge group of rank 32, $`\mathrm{U}(16)|_9\times \mathrm{U}(16)|_5`$ . Still, they are all free of irreducible gauge and gravitational anomalies, consistently with the vanishing of all R-R tadpoles .
A possible solution to the R-R tadpole conditions is $`m=13`$, $`n=3`$, $`d=4`$, that corresponds to a gauge group of rank 20, $`\mathrm{U}(13)|_9\times \mathrm{U}(3)|_9\times \mathrm{U}(4)|_5`$, with charged hypermultiplets in the representations $`(78+\overline{78},1;1)`$, in five copies of the $`(1,3+\overline{3};1)`$, in one copy of the $`(1,1;6+\overline{6})`$, in four copies of the $`(\overline{13},3;1)`$, in one copy of the $`(13,1;\overline{4})`$ and in one copy of the $`(1,\overline{3};4)`$. Alternatively, one can take $`m=14`$, $`n=2`$, $`d=8`$, obtaining a gauge group of rank 24, $`\mathrm{U}(14)|_9\times \mathrm{U}(2)|_9\times \mathrm{U}(8)|_5`$. The corresponding matter comprises charged hypermultiplets in the $`(91+\overline{91},1;1)`$, in one copy of the $`(1,1;28+\overline{28})`$, in four copies of the $`(\overline{14},2;1)`$, in one copy of the $`(14,1;\overline{8})`$, in one copy of the $`(1,2;8)`$, and in five copies of the $`(1,1+\overline{1},1)`$. On the other hand, for $`m=12`$, $`n=4`$, and thus $`d=0`$. This is a rather unusual supersymmetric $`Z_2`$ model without D5 branes, with a gauge group of rank 16, $`\mathrm{U}(12)\times \mathrm{U}(4)`$, and charged hypermultiplets in the representations $`(66+\overline{66},1)`$, in five copies of the $`(1,6+\overline{6})`$, and in four copies of the $`(\overline{12},4)`$. A distinctive feature of these spectra is that some of the matter occurs in multiple families. This peculiar phenomenon is a consequence of the multiplicities of Landau levels, that in these $`Z_2`$ orbifolds are multiples of two for each magnetised torus. Moreover, it should be appreciated that, in general, the rank reduction for the gauge group is not by powers of two as in the presence of a quantised antisymmetric tensor . Actually, these are not the first concrete examples of brane transmutation in type I vacua but, to the best of our knowledge, they are the first supersymmetric ones. $`Z_2`$ orientifolds without D5 branes have recently appeared in , where magnetised fractional D9 branes have been used to build six-dimensional asymmetric orientifolds with “brane supersymmetry breaking”.
One can also consider similar deformations of the model of , that has an $`𝒩=(1,0)`$ supersymmetric bulk spectrum with 17 tensor multiplets and four hypermultiplets. This alternative projection, allowed by the constraints in , introduces $`\mathrm{O9}_+`$ and $`\mathrm{O5}_{}`$ planes and thus requires, for consistency, an open sector resulting from the simultaneous presence of D9 branes and D5 antibranes, with “brane supersymmetry breaking”. A magnetised torus can now mimic D5 antibranes provided $`H_1=H_2`$, and one can then build several non-tachyonic configurations as in the previous case<sup>7</sup><sup>7</sup>7 There is a subtlety here. The different GSO projections for strings stretched between a D9 brane and a D5 antibrane would associate the low-lying twisted ND bosons to the characters $`O_4S_2(z_1)S_2(z_2)`$ and $`O_4C_2(z_1)C_2(z_2)`$, and thus now the choice $`H_1=H_2`$ would eliminate all tachyons even in the presence of D5 antibranes.. A particularly interesting one corresponds to a vacuum configuration without D5 antibranes, where the $`\mathrm{O5}_{}`$ charge is fully saturated by magnetised D9 branes. The resulting annulus and Möbius amplitudes can be obtained deforming the corresponding ones in , and read
$`𝒜`$ $`=`$ $`\frac{1}{4}\{(Q_o+Q_v)(0;0)[(m_1+m_2)^2P_1P_2+2n\overline{n}\stackrel{~}{P}_1\stackrel{~}{P}_2]`$ (22)
$``$ $`2(m_1+m_2)(n+\overline{n})(Q_o+Q_v)(z_1\tau ;z_2\tau ){\displaystyle \frac{k_1\eta }{\vartheta _1(z_1\tau )}}{\displaystyle \frac{k_2\eta }{\vartheta _1(z_2\tau )}}`$
$``$ $`(n^2+\overline{n}^2)(Q_o+Q_v)(2z_1\tau ;2z_2\tau ){\displaystyle \frac{2k_1\eta }{\vartheta _1(2z_1\tau )}}{\displaystyle \frac{2k_2\eta }{\vartheta _1(2z_2\tau )}}`$
$`+`$ $`\left[(m_1m_2)^2+2n\overline{n}\right](Q_oQ_v)(0;0)\left({\displaystyle \frac{2\eta }{\vartheta _2(0)}}\right)^2`$
$`+`$ $`2(m_1m_2)(n+\overline{n})(Q_oQ_v)(z_1\tau ;z_2\tau ){\displaystyle \frac{2\eta }{\vartheta _2(z_1\tau )}}{\displaystyle \frac{2\eta }{\vartheta _2(z_2\tau )}}`$
$`+`$ $`(n^2+\overline{n}^2)(Q_oQ_v)(2z_1\tau ;2z_2\tau ){\displaystyle \frac{2\eta }{\vartheta _2(2z_1\tau )}}{\displaystyle \frac{2\eta }{\vartheta _2(2z_2\tau )}}\},`$
and
$``$ $`=`$ $`\frac{1}{4}\{(m_1+m_2)(\widehat{Q}_o+\widehat{Q}_v)(0;0)P_1P_2`$
$``$ $`(n+\overline{n})(\widehat{Q}_o+\widehat{Q}_v)(2z_1\tau ;2z_2\tau ){\displaystyle \frac{2k_1\widehat{\eta }}{\widehat{\vartheta }_1(2z_1\tau )}}{\displaystyle \frac{2k_2\widehat{\eta }}{\widehat{\vartheta }_1(2z_2\tau )}}`$
$`+`$ $`\left(m_1+m_2\right)(\widehat{Q}_o\widehat{Q}_v)(0;0)\left({\displaystyle \frac{2\widehat{\eta }}{\widehat{\vartheta }_2(0)}}\right)^2`$
$`+`$ $`(n+\overline{n})(\widehat{Q}_o\widehat{Q}_v)(2z_1\tau ;2z_2\tau ){\displaystyle \frac{2\widehat{\eta }}{\widehat{\vartheta }_2(2z_1\tau )}}{\displaystyle \frac{2\widehat{\eta }}{\widehat{\vartheta }_2(2z_2\tau )}}\}.`$
In extracting the massless spectra of this class of models, it is important to notice that, at the special point $`H_1=H_2`$, all bosons from $`Q_o`$ with non-vanishing arguments and all fermions from $`Q_v`$ with non-vanishing arguments become massive. As a result, all massless fermions arising from strings affected by the internal magnetic fields have a reversed chirality, precisely as demanded by the cancellation of all irreducible anomalies. For $`|k_1|=|k_2|=2`$, one can obtain a gauge group $`\mathrm{SO}(8)\times \mathrm{SO}(16)\times \mathrm{U}(4)`$ and, aside from the corresponding $`𝒩=(1,0)`$ vector multiplets, the massless spectrum contains a hypermultiplet in the representation $`(8,16,1)`$, eight scalars in the $`(1,16,4+\overline{4})`$, two left-handed spinors in the $`(8,1,4+\overline{4})`$, and twelve scalars and five left-handed spinors in the $`(1,1,6+\overline{6})`$. Clearly, supersymmetry is explicitly broken on the magnetised D9 branes. Still, the resulting dilaton potential is effectively localised on the $`\mathrm{O5}_{}`$ plane, since it scales with the internal volume as in the undeformed model of .
These configurations present another interesting novelty: they have generalised Green-Schwarz couplings involving gauge fields and untwisted R-R forms, of the type
$$𝒮_{\mathrm{GS}}\underset{i}{}ϵ^{\mu _1\mathrm{}\mu _6}ϵ^{I_1\mathrm{}I_4}C_{I_1I_2\mu _3\mu _4\mu _5\mu _6}\mathrm{tr}\left(F_{\mu _1\mu _2}H_{I_3I_4}^i\right),$$
(24)
while standard orientifolds do not . In six dimensions, these four-forms are actually dual to axions $`a_{IJ}`$, and therefore this coupling can be rewritten in the form
$$𝒮_{\mathrm{GS}}\underset{i}{}\mathrm{tr}(A_\mu Q^i)H_{IJ}^i^\mu a^{IJ},$$
(25)
where $`Q^i`$ denote the group generators associated to the internal magnetic fields. Thus, additional U(1) gauge fields can acquire mass by a generalisation of the mechanism in , that in type-I strings generally involves several R-R forms. The (non-universal) axions involved in these Higgs mechanisms are the linear combinations $`H_{IJ}^ia^{IJ}`$.
A convenient way to recover these couplings uses, as in , a space-time magnetic background $``$ that, when introduced in the string amplitudes (S0.Ex15) and (S0.Ex18), deforms the space-time theta-functions according to
$$\frac{1}{\eta ^2}\frac{\vartheta _\alpha (0|\tau )}{\eta (\tau )}(q_\mathrm{L}+q_\mathrm{R})\tau \frac{\vartheta _\alpha (ϵ\tau |\tau )}{\vartheta _1(ϵ\tau |\tau )},$$
(26)
with $`\pi ϵ=\mathrm{tan}^1(q_\mathrm{L})+\mathrm{tan}^1(q_\mathrm{R})`$. As a result, the untwisted R-R tadpoles are modified, and become
$$\left[m+\overline{m}+n+\overline{n}32+q^2(H_1H_2+H_1+H_2)(n+\overline{n})\right]\sqrt{v_1v_2}\pm \frac{1}{\sqrt{v_1v_2}}\left[d+\overline{d}32\right].$$
(27)
Using the tadpole conditions (16) and the Dirac quantisation conditions (3), the terms linear in the space-time magnetic field identify the new Green-Schwarz couplings of eq. (24), needed to dispose of the new anomalous U(1) factors.
In conclusion, we have seen how in type I vacua a non-vanishing (anti)instanton density can be used to mimic BPS D5 (anti)branes, and we have exhibited some models with new distinctive features. These include supersymmetric $`T^4/Z_2`$ compactifications without D5 branes, or with gauge groups of unusual rank, that display new Green-Schwarz couplings of untwisted R-R forms. Several examples of this type can be constructed, both in six and in four dimensions. For instance, in the $`Z_3`$ orientifold of magnetic deformations allow the introduction of a net number of D5 (anti)branes, a setting to be contrasted with the models of , that only involve D5 brane-antibrane pairs. Models with “brane supersymmetry breaking”, in particular with additional brane-antibrane pairs, develop NS-NS tadpoles. These tadpoles are not eliminated by the magnetic deformation, and typically result in potentials that, although of run-away type for the dilaton, can in some cases stabilise some geometric moduli . Their presence requires a background redefinition , that was recently constructed explicitly in for the model in . In general, these vacua correspond to supergravity models frozen in phases of broken supersymmetry, where the presence of (lower-dimensional) non-supersymmetric couplings renders the field equations naively inconsistent, in complete analogy with ordinary gauge theories frozen in a Higgs phase. Although similar features were previously met in the anomalous Green-Schwarz couplings of , the peculiar supergravity models resulting from “brane supersymmetry breaking” clearly deserve further investigation.
Acknowledgements. We are grateful to P. Bain, F.S. Hassan, E. Kiritsis and R. Minasian for interesting discussions. C.A. is supported by the “Marie-Curie” fellowship HPMF-CT-1999-00256. This work was supported in part by the EEC contract HPRN-CT-2000-00122, in part by the EEC contract ERBFMRX-CT96-0090, in part by the EEC contract ERBFMRX-CT96-0045 and in part by the INTAS project 991590.
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# Adaptive Control for a Class of Nonlinear Systems with a Time-Varying Structure
## 1 Introduction
The field of nonlinear adaptive control developed rapidly in the last decade. The paper and others gave birth to an important branch of adaptive control theory, the nonlinear on-line function approximation based control, which includes neural (e.g., in ) and fuzzy (e.g., in ) approaches (note that there are several other relevant works on neural and fuzzy control, many of them cited in the references within the above papers). The neural and fuzzy approaches are most of the time equivalent, differing between each other only for the structure of the approximator chosen . Most of the papers deal with indirect adaptive control, trying first to identify the dynamics of the systems and eventually generating a control input according to the certainty equivalence principle (with some modification to add robustness to the control law), whereas very few authors (e.g., in ) use the direct approach, in which the controller directly generates the control input to guarantee stability.
Plants whose dynamics can be expressed in the so called “strict feedback form” have been considered, and techniques like backstepping and adaptive backstepping have emerged for their control. The papers present an extension of the tuning functions approach in which the nonlinearities of the strict feedback system are not assumed to be parametric uncertainties, but rather completely unknown nonlinearities to be approximated on-line with nonlinearly parameterized function approximators. Both the adaptive methods in and in attempt to approximate the dynamics of the plant on-line, so they may be classified as indirect adaptive schemes.
In this paper, we have combined an extension of the class of strict feedback systems considered in with the concept of a dynamic structure that depends on time, so as to propose a class of nonlinear systems with a time-varying structure, for which we develop a *direct* adaptive control approach. This class of systems is a generalization of the class of strict feedback systems traditionally considered in the literature. Moreover, the direct adaptive control developed here is, to our knowledge, the first of its kind in this context, and it presents several advantages with respect to indirect adaptive methods, including the fact that it needs less plant information to be implemented.
## 2 Direct Adaptive Control
Consider the class of continuous time nonlinear systems given by
$`\dot{x}_i`$ $`={\displaystyle \underset{j=1}{\overset{R}{}}}\rho _j(v)\left(\varphi _i^j(X_i)+\psi _i^j(X_i)x_{i+1}\right)`$
$`\dot{x}_n`$ $`={\displaystyle \underset{j=1}{\overset{R}{}}}\rho _j(v)\left(\varphi _n^j(X_n)+\psi _n^j(X_n)u\right)`$ (1)
where $`i=1,2,\mathrm{},n1`$, $`X_i=[x_1,\mathrm{},x_i]^{}`$, and $`X_n^n`$ is the state vector, which we assume measurable, and $`u`$ is the control input. The variable $`v^q`$ may be an additional input or a possibly exogenous “scheduling variable.” We assume that $`v`$ and its derivatives up to and including the $`(n1)^{th}`$ one are bounded and available for measurement, which may imply that $`v`$ is given by an external dynamical system. The functions $`\rho _j`$, $`j=1,\mathrm{},R`$ may be considered to be “interpolating functions” that produce the time-varying structural nature of system (2), since they combine $`R`$ systems in strict feedback form (given by the $`\varphi _i^j`$ and $`\psi _i^j`$ functions, $`i=1,\mathrm{},n`$, $`j=1\mathrm{},R`$) and the combination depends on time through the variable $`v`$ (thereby, the dynamics of the plant may be different at each time point depending on the scheduling variable). Here, we assume that the functions $`\rho _j`$ are $`n`$ times continuously differentiable, and that they satisfy, for all $`v^q`$, $`_{j=1}^R\rho _j(v)<\mathrm{}`$ and $`\left|\frac{^i\rho _j(v)}{v^i}\right|<\mathrm{}`$. Denote for convenience $`\varphi _i^c(X_i,v)=_{j=1}^R\rho _j(v)\varphi _i^j(X_i)`$ and $`\psi _i^c(X_i,v)=_{j=1}^R\rho _j(v)\psi _i^j(X_i).`$ We will assume that $`\varphi _i^c`$ and $`\psi _i^c`$ are sufficiently smooth in their arguments, and that they satisfy, for all $`X_i^i`$ and $`v^q`$, $`i=1,\mathrm{},n`$, $`\varphi _i^c(0,v)=0`$ and $`\psi _i^c(X_i,v)0.`$
Here, we will develop a direct adaptive control method for the class of systems (2). We assume that the interpolation functions $`\rho _j`$ are known, but the functions $`\varphi _i^j`$ and $`\psi _i^j`$ (which constitute the underlying time-varying dynamics of the system) are unknown. In an indirect adaptive methodology one would attempt to identify the unknown functions and then construct a stabilizing control law based on the approximations to the plant dynamics. Here, however, we will postulate the existence of an ideal control law (based on the assumption that the plant belongs to the class of systems (2)) which possesses some desired stabilizing properties, and then we devise adaptation laws that attempt to approximate the ideal control equation. This approximation will be performed within a compact set $`𝒮_{x_n}^n`$ of arbitrary size which contains the origin. In this manner, the results obtained are semi-global, in the sense that they are valid as long as the state remains within $`𝒮_{x_n}`$, but this set can be made as large as desired by the designer. In particular, with enough plant information it can be made large enough that the state never exits it, since, as will be shown a bound can be placed on the state transient. Furthermore, as will be indicated below, the stability can be made global by using bounding control terms.
For each vector $`X_i`$ we will assume the existence of a compact set $`𝒮_{x_i}^i`$ specified by the designer. We will consider trajectories within the compact sets $`𝒮_{x_i}`$, $`i=1,\mathrm{},n`$, where the sets are constructed such that $`𝒮_{x_i}𝒮_{x_{i+1}}`$, for $`i=1,\mathrm{},n1`$. We assume the existence of bounds $`\underset{¯}{\psi }_i^c`$, $`\overline{\psi }_i^c`$, and $`\psi _{i_d}^c`$, $`i=1,\mathrm{},n`$ (*not necessarily known*), such that for all $`v^q`$ and $`X_i𝒮_{x_i}`$, $`i=1,\mathrm{},n`$,
$`0`$ $`<\underset{¯}{\psi }_i^c\psi _i^c(X_i,v)\overline{\psi }_i^c<\mathrm{}`$
$`\left|\dot{\psi }_i^c\right|`$ $`=\left|{\displaystyle \underset{j=1}{\overset{R}{}}}\left({\displaystyle \frac{\rho _j(v)}{v}}\dot{v}\psi _i^j(X_i)+\rho _j(v){\displaystyle \frac{\psi _i^j(X_i)}{X_i}}\dot{X}_i\right)\right|\psi _{i_d}^c.`$ (2)
This assumption implies that the affine terms in the plant dynamics have a bounded gain and a bounded rate of change. Since the functions $`\psi _i^c`$ are assumed continuous, they are therefore bounded within $`𝒮_{x_i}`$. Similarly, note that even though the term $`|\dot{X}_i|`$ may not necessarily be globally bounded, it will have a constant bound within $`𝒮_{x_i}`$ due to the continuity assumptions we make. Therefore, assumption (2) will always be satisfied within $`𝒮_{x_n}`$. Moreover, in the simplest of cases, the first part of assumption (2) is satisfied globally when the functions $`\psi _i^j`$ are constant or sector bounded for all $`X_i^i`$.
The class of plants (2) is, to our knowledge, the most general class of systems considered so far within the context of adaptive control based on backstepping. In particular, in both and , which are indirect adaptive approaches, the input functions $`\psi _i^j`$ are assumed to be constant for $`i=1,\mathrm{},n`$. This assumption allows the authors of those works to perform a simpler stability analysis, which becomes more complex in the general case . Also, the addition of the interpolation functions $`\rho _j`$, $`j=1,\mathrm{},R`$, extends the class of strict feedback systems to one including systems with a time-varying structure , as well as systems falling in the domain of gain scheduling (where the plant dynamics are identified at different operating points and then interpolated between using a scheduling variable). Note that if we let $`R=1`$ and $`\rho _1(v)=1`$ for all $`v`$, together with $`\psi _i^c=1`$, $`i=1,\mathrm{},n`$, we have the particular case considered in .
The direct approach presented here has several advantages with respect to indirect approaches such as in . In particular, bounds on the input functions $`\psi _i^j`$ are only assumed to exist, but need *neither to be known nor to be estimated*. This is because the ideal law is formulated so that there is not an explicit need to include information about the bounds in the actual control law. Moreover, although assumption (2) appears to be more restrictive than what is needed in the indirect adaptive case, it is in fact not so due to the fact that the stability results are semi-global (i.e., since we are operating within the compact sets $`𝒮_{x_n}`$, continuity of the affine terms automatically implies the satisfaction of the second part of assumption (2)).
### 2.1 Direct Adaptive Control Theorem
Next, we state our main result and then show its proof<sup>1</sup><sup>1</sup>1We will generally omit the arguments of functions for brevity.. For convenience, we use the notation $`\nu _i=[v,\dot{v},\mathrm{},v^{(i1)}]^{q\times i}`$, $`i=1,\mathrm{},n`$.
Theorem 1: Consider system (2) with the state vector $`X_n`$ measurable and the scheduling matrix $`\nu _{n1}`$ measurable and bounded, together with the above stated assumptions on $`\varphi _i^c`$, $`\psi _i^c`$ and $`\rho _j`$, and (2). Assume also that $`\nu _i(0)𝒮_{v_i}^{q\times i}`$, $`X_i(0)𝒮_{x_i}^i`$, $`i=1,\mathrm{},n`$, where $`𝒮_{v_i}`$ and $`𝒮_{x_i}`$ are compact sets specified by the designer, and large enough that $`\nu _i`$ and $`X_i`$ do not exit them. Consider the diffeomorphism $`z_1=x_1`$, $`z_i=x_i\widehat{\alpha }_{i1}\alpha _{i1}^s`$, $`i=2,\mathrm{},n`$, with $`\widehat{\alpha }_i(X_i,\nu _i)=_{j=1}^R\rho _j(v)\widehat{\theta }_{\alpha _i^j}^{}\zeta _{\alpha _i^j}(X_i,\nu _i)`$ and $`\alpha _i^s(z_i,z_{i1})=k_iz_iz_{i1}`$, with $`k_i>0`$ and $`z_0=0`$. Assume the functions $`\zeta _{\alpha _i^j}(X_i,\nu _i)`$ to be at least $`ni`$ times continuously differentiable, and to satisfy, for $`i=1,\mathrm{},n`$, $`j=1,\mathrm{},R`$,
$$\left|\frac{^{ni}\zeta _{\alpha _i^j}}{[X_i,\nu _i]^{ni}}\right|<\mathrm{}.$$
(3)
Consider the adaptation laws for the parameter vectors $`\widehat{\theta }_{\alpha _i^j}^{N_{\alpha _i^j}}`$, $`N_{\alpha _i^j}`$, $`\dot{\widehat{\theta }}_{\alpha _i^j}=\rho _j\gamma _{\alpha _i^j}\zeta _{\alpha _i^j}z_i\sigma _{\alpha _i^j}\widehat{\theta }_{\alpha _i^j}`$, where $`\gamma _{\alpha _i^j}>0`$, $`\sigma _{\alpha _i^j}>0`$, $`i=1,\mathrm{},n`$, $`j=1,\mathrm{},R`$ are design parameters. Then, the control law $`u=\widehat{\alpha }_n+\alpha _n^s`$ guarantees boundedness of all signals and convergence of the states to the residual set
$$𝒟_d=\{X_n\mathrm{}^n:\underset{i=1}{\overset{n}{}}z_i^2\frac{2\underset{¯}{\psi }_mW_d}{\beta _d}\}.$$
(4)
where $`\underset{¯}{\psi }_m=\mathrm{min}_{1in}\overline{\psi }_i^c`$, $`\beta _d`$ is a constant, and $`W_d`$ measures approximation errors and ideal parameter sizes, and its magnitude can be reduced through the choice of the design constants $`k_i`$, $`\gamma _{\alpha _i^j}`$ and $`\sigma _{\alpha _i^j}`$.
###### Proof.
The proof requires $`n`$ steps, and is performed inductively. First, let $`z_1=x_1`$, and $`z_2=x_2\widehat{\alpha }_1\alpha _1^s`$, where $`\widehat{\alpha }_1`$ is the approximation to an ideal signal $`\alpha _1^{}`$ (“ideal” in the sense that if we had $`\widehat{\alpha }_1=\alpha _1^{}`$ we would have a globally asymptotically stable closed loop without need for the stabilizing term $`\alpha _1^s`$), and $`\alpha _1^s`$ will be given below. Let $`c_1>0`$ be a constant such that $`c_1>\frac{\psi _{1_d}^c}{2\underset{¯}{\psi }_1^c}`$, and $`\alpha _1^{}(x_1,v)=\frac{1}{\psi _1^c}\left(\varphi _1^cc_1z_1\right).`$ Since the ideal control $`\alpha _1^{}`$ is smooth, it may be approximated with arbitrary accuracy for $`v`$ and $`x_1`$ within the compact sets $`𝒮_{v_1}^q`$ and $`𝒮_{x_1}`$, respectively, as long as the size of the approximator can be made arbitrarily large.
For approximators of finite size let $`\alpha _1^{}(x_1,v)=_{j=1}^R\rho _j(v)\theta _{\alpha _1^j}^{^{}}\zeta _{\alpha _1^j}(v,x_1)+\delta _{\alpha _1}(v,x_1),`$ where the parameter vectors $`\theta _{\alpha _1^j}^{}^{N_{\alpha _1^j}}`$, $`N_{\alpha _1^j}`$, are optimum in the sense that they minimize the representation error $`\delta _{\alpha _1}`$ over the set $`𝒮_{x_1}\times 𝒮_{v_1}`$ and suitable compact parameter spaces $`\mathrm{\Omega }_{\alpha _1^j}`$, and $`\zeta _{\alpha _1^j}(x_1,v)`$ are defined via the choice of the approximator structure (see for an example of a choice for $`\zeta _{\alpha _i^j}`$). The parameter sets $`\mathrm{\Omega }_{\alpha _1^j}`$ are simply mathematical artifacts. As a result of the stability proof the approximator parameters are bounded using the adaptation laws in Theorem 2.1, so $`\mathrm{\Omega }_{\alpha _1^j}`$ does not need to be defined explicitly, and no parameter projection (or any other “artificial” means of keeping the parameters bounded) is required. The representation error $`\delta _{\alpha _1}`$ arises because the sizes $`N_{\alpha _i^j}`$ are finite, but it may be made arbitrarily small within $`𝒮_{x_1}\times 𝒮_{v_1}`$ by increasing $`N_{\alpha _i^j}`$ (i.e., we assume the chosen approximator structures possess the “universal approximation property”). In this way, there exists a constant bound $`d_{\alpha _1}>0`$ such that $`|\delta _{\alpha _1}|d_{\alpha _1}<\mathrm{}`$. To make the proof logically consistent, however, we need to assume that some knowledge about this bound and a bound on $`\theta _{\alpha _1^j}^{}`$ are available (since in this case it becomes possible to guarantee a priori that $`𝒮_{x_1}\times 𝒮_{v_1}`$ is large enough). However, in practice some amount of redesign may be required, since these bounds are typically guessed by the designer
Let $`\mathrm{\Phi }_{\alpha _1^j}=\widehat{\theta }_{\alpha _1^j}\theta _{\alpha _1^j}^{}`$ denote the parameter error, and approximate $`\alpha _1^{}`$ with $`\widehat{\alpha }_1(x_1,v,\widehat{\theta }_{\alpha _1^j};j=1,\mathrm{},R)=_{j=1}^R\rho _j(v)\widehat{\theta }_{\alpha _1^j}^{}\zeta _{\alpha _1^j}(x_1,v).`$ Hence, we have a linear in the parameters approximator with parameter vectors $`\widehat{\theta }_{\alpha _1^j}`$. Note that the structural dependence on time of system (2) is reflected in the controller, because $`\widehat{\alpha }_1`$ can be viewed as using the functions $`\rho _j(v)`$ to interpolate between “local” controllers of the form $`\widehat{\theta }_{\alpha _1^j}^{}\zeta _{\alpha _1^j}(x_1,v)`$, respectively. Notice that since the functions $`\rho _j`$ are assumed continuous and $`v`$ bounded, the signal $`\widehat{\alpha }_1`$ is well defined for all $`v𝒮_{v_1}`$.
Consider the dynamics of the transformed state, $`\dot{z}_1=\varphi _1^c+\psi _1^c(z_2+\widehat{\alpha }_1+\alpha _1^s)+\psi _1^c(\alpha _1^{}\alpha _1^{})=c_1z_1+\psi _1^cz_2+\psi _1^c(\widehat{\alpha }_1\alpha _1^{})+\psi _1^c\alpha _1^s=c_1z_1+\psi _1^cz_2+\psi _1^c\left(_{j=1}^R\rho _j\mathrm{\Phi }_{\alpha _1^j}^{}\zeta _{\alpha _1^j}\delta _{\alpha _1^j}\right)+\psi _1^c\alpha _1^s.`$ Let $`V_1=\frac{1}{2\psi _1^c}z_1^2+\frac{1}{2}_{j=1}^R\frac{\mathrm{\Phi }_{\alpha _1^j}^{}\mathrm{\Phi }_{\alpha _1^j}}{\gamma _{\alpha _1^j}}`$, and examine its derivative, $`\dot{V}_1=\frac{2\psi _1^c(2z_1\dot{z}_1)2z_1^2\dot{\psi }_1^c}{4\psi _1^{c^2}}+_{j=1}^R\frac{\mathrm{\Phi }_{\alpha _1^j}^{}\dot{\mathrm{\Phi }}_{\alpha _1^j}}{\gamma _{\alpha _1^j}}.`$ Using the expression for $`\dot{z}_1`$, $`\dot{V}_1=\frac{c_1z_1^2}{\psi _1^c}+z_1z_2+z_1_{j=1}^R\rho _j\mathrm{\Phi }_{\alpha _1^j}^{}\zeta _{\alpha _1^j}z_1\delta _{\alpha _1^j}+z_1\alpha _1^s\frac{1}{2}z_1^2\frac{\dot{\psi }_1^c}{\psi _1^{c^2}}+_{j=1}^R\frac{\mathrm{\Phi }_{\alpha _1^j}^{}\dot{\mathrm{\Phi }}_{\alpha _1^j}}{\gamma _{\alpha _1^j}}.`$ Choose the adaptation law $`\dot{\widehat{\theta }}_{\alpha _1^j}=\dot{\mathrm{\Phi }}_{\alpha _1^j}=\rho _j\gamma _{\alpha _1^j}\zeta _{\alpha _1^j}z_1\sigma _{\alpha _1^j}\widehat{\theta }_{\alpha _1^j},`$ with design constants $`\gamma _{\alpha _1^j}>0`$, $`\sigma _{\alpha _1^j}>0`$, $`j=1,\mathrm{},R`$ (we think of $`\sigma _{\alpha _1^j}\widehat{\theta }_{\alpha _1^j}`$ as a “leakage term”). Also, note that for any constant $`k_1>0`$, $`z_1\delta _{\alpha _1^j}|z_1|d_{\alpha _1}k_1z_1^2+\frac{d_{\alpha _1}^2}{4k_1}.`$ We pick $`\alpha _1^s=k_1z_1.`$
Notice also that, completing squares, $`\mathrm{\Phi }_{\alpha _1^j}^{}\widehat{\theta }_{\alpha _1^j}=\mathrm{\Phi }_{\alpha _1^j}^{}(\mathrm{\Phi }_{\alpha _1^j}+\theta _{\alpha _1^j}^{})\frac{|\mathrm{\Phi }_{\alpha _1^j}|^2}{2}+\frac{|\theta _{\alpha _1^j}^{}|^2}{2}.`$ Finally, observe that $`\frac{z_1^2}{\psi _1^c}\left(c_1+\frac{\dot{\psi }_1^c}{2\psi _1^c}\right)\frac{z_1^2}{\psi _1^c}\left(c_1\frac{\psi _{1_d}^c}{2\underset{¯}{\psi }_1^c}\right)\frac{\overline{c}_1z_1^2}{\overline{\psi }_1^c},`$ with $`\overline{c}_1=c_1\frac{\psi _{1_d}^c}{2\underset{¯}{\psi }_1^c}>0`$. Then, we obtain $`\dot{V}_1\frac{\overline{c}_1z_1^2}{\overline{\psi }_1^c}\frac{1}{2}_{j=1}^R\sigma _{\alpha _1^j}\frac{|\mathrm{\Phi }_{\alpha _1^j}|^2}{\gamma _{\alpha _1^j}}+z_1z_2+\frac{d_{\alpha _1}^2}{4k_1}+\frac{1}{2}_{j=1}^R\sigma _{\alpha _1^j}\frac{\theta _{\alpha _1^j}^{}}{\gamma _{\alpha _1^j}}.`$ This completes the first step of the proof.
We may continue in this manner up to the $`n^{th}`$ step<sup>2</sup><sup>2</sup>2We omit intermediate steps for brevity., where we have $`z_n=x_n\widehat{\alpha }_{n1}\alpha _{n1}^s`$, with $`\widehat{\alpha }_{n1}`$ and $`\alpha _{n1}^s`$ defined as in Theorem 2.1. Consider the ideal signal $`\alpha _n^{}(X_n,\nu _n)=\frac{1}{\psi _n^c}\left(\varphi _n^cc_nz_n+\dot{\widehat{\alpha }}_{n1}+\dot{\alpha }_{n1}^s\right)`$ with $`c_n>\frac{\psi _{n_d}^c}{2\underset{¯}{\psi }_n^c}`$. Notice that, even though the terms $`\dot{\widehat{\theta }}_{\alpha _{n1}^j}`$ appear in $`\alpha _n^{}`$ through the partial derivatives in $`\dot{\widehat{\alpha }}_{n1}`$, $`\widehat{\theta }_{\alpha _{n1}^j}`$ does not need to be an input to $`\alpha _n^{}`$, since the resulting product of the partial derivatives and $`\dot{\widehat{\theta }}_{\alpha _{n1}^j}`$ can be expressed in terms of $`z_1,\mathrm{},z_{n1}`$, $`v`$ and $`\sigma _{\alpha _{n1}^j}\widehat{\alpha }_{n1}`$. To simplify the notation, however, we will omit the dependencies on inputs other than $`X_i`$ and $`\nu _i`$, but bearing in mind that, when implementing this method, more inputs may be required to satisfy the proof. Also, note that by assumption (3), $`|\alpha _n^{}|<\mathrm{}`$ for bounded arguments. Therefore, we may represent $`\alpha _n^{}`$ with $`\alpha _n^{}(X_n,\nu _n)=_{j=1}^R\rho _j(v)\theta _{\alpha _n^j}^{^{}}\zeta _{\alpha _n^j}(X_n,\nu _n)+\delta _{\alpha _n}(X_n,\nu _n)`$ for $`X_n𝒮_{x_n}^n`$ and $`\nu _n𝒮_{v_n}^{q\times n}`$. The parameter vector $`\theta _{\alpha _n^j}^{}^{N_{\alpha _n^j}}`$, $`N_{\alpha _n^j}`$ is an optimum within a compact parameter set $`\mathrm{\Omega }_{\alpha _n}`$, in a sense similar to $`\theta _{\alpha _1^j}^{}`$, so that for $`(X_n,\nu _n)𝒮_{x_n}\times 𝒮_{v_n}`$, $`|\delta _{\alpha _n}|d_{\alpha _n}<\mathrm{}`$ for some bound $`d_{\alpha _n}>0`$. Let $`\mathrm{\Phi }_{\alpha _n^j}=\widehat{\theta }_{\alpha _n^j}\theta _{\alpha _n^j}^{}`$, and consider the approximation $`\widehat{\alpha }_n`$ as given in Theorem 2.1. The control law $`u=\widehat{\alpha }_n+\alpha _n^s`$ yields $`\dot{z}_n=\varphi _n^c+\psi _n^c(\widehat{\alpha }_n+\alpha _n^s)\dot{\widehat{\alpha }}_{n1}\dot{\alpha }_{n1}^s+\psi _n^c(\alpha _n^{}\alpha _n^{})=c_nz_n+\psi _n^c\left(_{j=1}^R\rho _j(v)\mathrm{\Phi }_{\alpha _n^j}^{}\zeta _{\alpha _n^j}\delta _{\alpha _n}\right)+\psi _n^c\alpha _n^s.`$ Choose the Lyapunov function candidate $`V=V_{n1}+\frac{1}{2\psi _n^c}z_n^2+\frac{1}{2}_{j=1}^R\frac{\mathrm{\Phi }_{\alpha _n^j}^{}\mathrm{\Phi }_{\alpha _n^j}}{\gamma _{\alpha _n^j}}`$ and examine its derivative, $`\dot{V}=\dot{V}_{n1}\frac{c_nz_n^2}{\psi _n^c}+z_n_{j=1}^R\rho _j(v)\mathrm{\Phi }_{\alpha _n^j}^{}\zeta _{\alpha _n^j}z_n\delta _{\alpha _n}+z_n\alpha _n^s\frac{1}{2}z_n^2\frac{\dot{\psi }_n^c}{\psi _n^{c^2}}+_{j=1}^R\frac{\mathrm{\Phi }_{\alpha _n^j}^{}\dot{\mathrm{\Phi }}_{\alpha _n^j}}{\gamma _{\alpha _n^j}}`$. One can show inductively that $`\dot{V}_{n1}_{i=1}^{n1}\frac{\overline{c}_iz_i^2}{\overline{\psi }_i^c}\frac{1}{2}_{i=1}^{n1}_{j=1}^R\sigma _{\alpha _i^j}\frac{|\mathrm{\Phi }_{\alpha _i^j}|^2}{\gamma _{\alpha _i^j}}+z_{n1}z_n+_{i=1}^{n1}\frac{d_{\alpha _i}^2}{4k_i}+\frac{1}{2}_{i=1}^{n1}_{j=1}^R\sigma _{\alpha _i^j}\frac{|\theta _{\alpha _i^j}^{}|^2}{\gamma _{\alpha _i^j}}`$ with constants $`\overline{c}_i=c_i\frac{\psi _{i_d}}{2\underset{¯}{\psi }_i^c}>0`$, $`i=1,\mathrm{},n`$. The choice of adaptation laws for $`\theta _{\alpha _n^j}`$ and of $`\alpha _n^s`$ in Theorem 2.1, together with the observations that $`\frac{\sigma _{\alpha _n^j}}{\gamma _{\alpha _n^j}}\mathrm{\Phi }_{\alpha _n^j}^{}\widehat{\theta }_{\alpha _n^j}\frac{\sigma _{\alpha _n^j}}{\gamma _{\alpha _n^j}}\frac{|\mathrm{\Phi }_{\alpha _n^j}|^2}{2}+\frac{\sigma _{\alpha _n^j}}{\gamma _{\alpha _n^j}}\frac{|\theta _{\alpha _n^j}^{}|^2}{2}`$, $`z_n\delta _{\alpha _n^j}k_nz_n^2+\frac{d_{\alpha _n}}{4k_n}`$, with $`k_n>0`$ and $`\frac{z_n^2}{\psi _n^c}\left(c_n+\frac{\dot{\psi }_n^c}{2\psi _n^c}\right)\frac{\overline{c}_nz_n^2}{\overline{\psi }_n^c}`$ imply
$$\dot{V}\underset{i=1}{\overset{n}{}}\frac{\overline{c}_iz_i^2}{\overline{\psi }_i^c}\frac{1}{2}\underset{i=1}{\overset{n}{}}\underset{j=1}{\overset{R}{}}\sigma _{\alpha _i^j}\frac{|\mathrm{\Phi }_{\alpha _i^j}|^2}{\gamma _{\alpha _i^j}}+W_d,$$
(5)
where $`W_d`$ contains the combined effects of representation errors and ideal parameter sizes, and is given by $`W_d=_{i=1}^n\frac{d_{\alpha _i}^2}{4k_i}+\frac{1}{2}_{i=1}^n_{j=1}^R\sigma _{\alpha _i^j}\frac{|\theta _{\alpha _i^j}^{}|^2}{\gamma _{\alpha _i^j}}.`$ Note that if $`_{i=1}^n\frac{\overline{c}_iz_i^2}{\overline{\psi }_i^c}W_d`$ or $`\frac{1}{2}_{i=1}^n_{j=1}^R\sigma _{\alpha _i^j}\frac{|\mathrm{\Phi }_{\alpha _i^j}|^2}{\gamma _{\alpha _i^j}}W_d`$, then we have $`\dot{V}0`$. Furthermore, letting $`\underset{¯}{\psi }_m=\mathrm{min}_{1in}(\underset{¯}{\psi }_i^c)`$, $`\overline{\psi }_m=\mathrm{max}_{1in}(\overline{\psi }_i^c)`$, and defining $`\overline{c}_0=\mathrm{min}_{1in}(\overline{c}_i)`$, $`\psi _m=\frac{\underset{¯}{\psi }_m}{\overline{\psi }_m}`$ and $`\sigma _0=\mathrm{min}_{1in,1jR}\left(\sigma _{\alpha _i^j}\right)`$ we have $`_{i=1}^n\frac{\overline{c}_iz_i^2}{\overline{\psi }_i^c}\overline{c}_0_{i=1}^n\frac{z_i^2}{\overline{\psi }_i^c}=\overline{c}_0_{i=1}^n\frac{z_i^2}{\psi _i^c}\frac{\psi _i^c}{\overline{\psi }_i^c}\overline{c}_0_{i=1}^n\frac{z_i^2}{\psi _i^c}\frac{\underset{¯}{\psi }_i^c}{\overline{\psi }_i^c}\overline{c}_0\psi _m_{i=1}^n\frac{z_i^2}{\psi _i^c}`$ and $`\frac{1}{2}_{i=1}^n_{j=1}^R\sigma _{\alpha _i^j}\frac{|\mathrm{\Phi }_{\alpha _i^j}|^2}{\gamma _{\alpha _i^j}}\sigma _0\frac{1}{2}_{i=1}^n_{j=1}^R\frac{|\mathrm{\Phi }_{\alpha _i^j}|^2}{\gamma _{\alpha _i^j}}.`$ Then, letting $`\beta _d=\mathrm{min}(2\overline{c}_0\psi _m,\sigma _0)`$, we have that if
$$V=\frac{1}{2}\underset{i=1}{\overset{n}{}}\frac{z_i^2}{\psi _i^c}+\frac{1}{2}\underset{i=1}{\overset{n}{}}\underset{j=1}{\overset{R}{}}\frac{|\mathrm{\Phi }_{\alpha _i^j}|^2}{\gamma _{\alpha _i^j}}V_0$$
(6)
with $`V_0=\frac{W_d}{\beta _d}`$, then $`\dot{V}0`$ and all signals in the closed loop are bounded. Furthermore, we have $`\dot{V}\beta _dV+W_d`$, which implies that $`0V(t)\frac{W_d}{\beta _d}+\left(V(0)\frac{W_d}{\beta _d}\right)e^{\beta _dt}`$ so that both the transformed states and the parameter error vectors converge to a bounded set. Finally, we conclude from the upper bound on $`V(t)`$ that the state vector $`X_n`$ converges to the residual set (4). $`\mathrm{}`$
Remark 1: The representation error bounds and the size of the ideal parameter vectors are assumed known, since they affect the size of the residual set to which the states converge. It is possible to augment the direct adaptive algorithm with “auto-tuning” capabilities (similar to ), which would relax the need for these bounds.
Furthermore, note that the stability result of Theorem 2.1 is semi-global, in the sense that it is valid within the compact sets $`𝒮_{v_i}`$ and $`𝒮_{x_i}`$, $`i=1,\mathrm{},n`$, which can be made arbitrarily large. The stability result may be made global by adding a high gain bounding control term to the control law. Such a term may be particularly useful when, due to a complete lack of a priori knowledge, the control designer is unable to guarantee that the compact sets $`𝒮_{x_i}`$, $`i=1,\mathrm{},n`$, are large enough so that the state will not exit them before the controller has time to bring the state inside $`𝒟_d`$; moreover, it may also happen that due to a poor design and poor system knowledge, $`𝒟_d`$ is not contained in $`𝒮_{x_n}`$. In this case, too, bounding control terms may be helpful until the design is refined and improved. However, using bounding control requires explicit knowledge of functional upper bounds of $`|\psi _i^c(v,X_i)|`$, and also of the lower bounds $`\underset{¯}{\psi }_i^c`$, $`i=1,\mathrm{},n`$, whose knowledge we do not mandate in Theorem 2.1. Bounding terms may be added to the diffeomorphism in Theorem 2.1, but we do not present the analysis since it is similar to the one we present here and it is algebraically tedious; we simply note, though, that the bounding terms have to be smooth (because they need to be differentiable), so they need to be defined in terms of smooth approximations to the sign, saturation and absolute value functions that are typically used in this approach.
Remark 2: If the bounds $`\underset{¯}{\psi }_i^c`$, $`\overline{\psi }_i^c`$ and $`\psi _{i_d}^c`$ are known, it becomes possible for the designer to directly set the constants $`c_i`$ in the control law. Notice that with knowledge of these bounds, the term $`\underset{¯}{\psi }_m`$ is also known, and we can pick constants $`c_i`$ such that $`c_i>\frac{\psi _{i_d}^c}{2\underset{¯}{\psi }_i^c}`$. Define the auxiliary functions $`\eta _i=c_iz_i`$. We may explicitly set the constant $`c_i`$ in $`\alpha _i^{}`$ if we let $`\eta _i`$ be an input to the $`i^{th}`$ approximator structure, i.e., if we let $`\alpha _i^{}(X_i,\nu _i,\dot{X}_{r_i},\eta _i)=_{j=1}^R\rho _j(v)\theta _{\alpha _i^j}^{^{}}\zeta _{\alpha _i^j}(X_i,\nu _i,\dot{X}_{r_i},\eta _i)+\delta _{\alpha _i}`$. Then, the approximators used in the control procedure are given by $`\widehat{\alpha }_i(X_i,\nu _i,\dot{X}_{r_i},\eta _i)=_{j=1}^R\rho _j(v)\widehat{\theta }_{\alpha _i^j}^{}\zeta _{\alpha _i^j}(X_i,\nu _i,\dot{X}_{r_i},\eta _i)`$ and the stability analysis can be carried out as expected.
### 2.2 Performance Analysis: $`_{\mathrm{}}`$ Bounds and Transient Design
The stability result of Theorem 2.1 is useful in that it indicates conditions to obtain a stable closed-loop behavior for a plant belonging to the class given by (2). However, it is not immediately clear how to choose the several design constants to improve the control performance. Here we concentrate on the tracking problem, and present design guidelines with respect to an $`_{\mathrm{}}`$ bound on the tracking error. We are interested in having $`x_1`$ track the reference model state $`x_{r_1}`$ of the reference model $`\dot{x}_{r_i}=x_{r_{i+1}},\text{ }i=1,2,\mathrm{},n1`$, $`\dot{x}_{r_n}=f_r(X_{r_n},r)`$ with bounded reference input $`r(t)`$. Now, we need to use the diffeomorphism $`z_1=x_1x_{r_1}`$, $`z_i=x_i\widehat{\alpha }_{i1}\alpha _{i1}^s`$, $`i=2,\mathrm{},n`$ with $`\alpha _1^{}(x_1,v,\dot{x}_{r_1})=\frac{1}{\psi _1^c}\left(\varphi _1^cc_1z_1+x_{r_2}\right)`$ and $`\alpha _i^{}(X_i,\nu _i,\dot{X}_{r_i})=\frac{1}{\psi _i^c}\left(\varphi _i^cc_iz_i+\dot{\widehat{\alpha }}_i+\dot{\alpha }_i^s\right)`$ for $`i=2,\mathrm{},n`$. The stability proof needs to be modified accordingly, and it can be shown that the tracking error $`|x_1x_{r_1}|`$ converges to a neighborhood of size $`\sqrt{\frac{2\underset{¯}{\psi }_mW_d}{\beta _d}}`$.
From the upper bound on $`V(t)`$ we can write $`V(t)\frac{W_d}{\beta _d}+V(0)e^{\beta _dt}`$. From here, it follows that $`\frac{1}{2}_{i=1}^n\frac{z_i^2(t)}{\psi _i^c(t)}\frac{W_d}{\beta _d}+\left(\frac{1}{2}_{i=1}^n\frac{z_i^2(0)}{\psi _i^c(0)}+\frac{1}{2}_{i=1}^n_{j=1}^R\frac{|\mathrm{\Phi }_{\alpha _i^j}(0)|^2}{\gamma _{\alpha _i^j}}\right)e^{\beta _dt}.`$ The terms $`z_i(0)`$ depend on the design constants in a complex manner. For this reason, rather than trying to take them into account in the design procedure, we follow the trajectory initialization approach taken in , which allows the designer to set $`z_i(0)=0`$, $`i=1,\mathrm{},n`$ by an appropriate choice of the reference model’s initial conditions. In our case, in addition to the assumption that it is possible to set the initial conditions of the reference model, we will have to assume certain invertibility conditions on the approximators. In particular, since $`z_1(0)=x_1(0)x_{r_1}(0)`$, for $`z_1(0)=0`$ we need to set $`x_{r_1}(0)=x_1(0)`$.
For the $`i^{th}`$ transformed state $`z_i`$, $`i=2,\mathrm{},n`$, $`z_i(0)=x_i(0)\widehat{\alpha }_{i1}(0)\alpha _{i1}^s(0)`$. Notice that $`\alpha _{i1}^s(0)=\alpha _{i1}^s(z_{i1}(0),z_{i2}(0))`$, so that if $`z_{i1}(0)=0`$ and $`z_{i2}(0)=0`$ we have $`\alpha _{i1}^s(0)=0`$. In particular, notice that this holds for $`i=2`$. In this case, to set $`z_2(0)=0`$ we need to have $`\widehat{\alpha }_1(x_1(0),v(0),x_{r_2}(0))=x_2(0)`$. This equation can be solved analytically (or numerically) for $`x_{r_2}(0)`$ provided $`\frac{\widehat{\alpha }_1}{x_{r_2}}|_{t=0}0`$. This is not an unreasonable condition, since it depends on the choice of approximator structure the designer makes. The structure can be chosen so that it satisfies this condition. Granted this is the case, it clearly holds that $`\alpha _2^s(0)=0`$, and the same procedure can be inductively carried out for $`i=3,\mathrm{},n`$, with the choices $`\widehat{\alpha }_{i1}(X_{i1}(0),\nu _{i1}(0),x_{r_i}(0))=x_i(0)`$.
This procedure yields the simpler bound $`_{i=1}^nz_i^2(t)\frac{2\underset{¯}{\psi }_mW_d}{\beta _d}+\underset{¯}{\psi }_m\left(_{i=1}^n_{j=1}^R\frac{|\mathrm{\Phi }_{\alpha _i^j}(0)|^2}{\gamma _{\alpha _i^j}}\right)e^{\beta _dt}`$. We would like to make this bound small, so that the transient excursion of the tracking error is small. Notice that we do not have direct control on the size of $`\beta _d`$, since this term depends on the unknown constants $`c_i`$, which appear in the ideal signals $`\alpha _i^{}`$. Even though it is not necessary to be able to set $`\beta _d`$ to reduce the size of the bound, it is possible to do so if the bounds $`\underset{¯}{\psi }_i^c`$, $`\overline{\psi }_i^c`$ and $`\psi _{i_d}^c`$ are known.
At this point, it becomes more clear how to choose the constants to achieve a smaller bound. Recalling the expression of $`W_d`$, note that, first, one may want to have $`\beta _d>1`$, so that $`W_d`$ is not made larger when divided by $`\beta _d`$, and so that the convergence is faster. This may be achieved by setting $`c_i`$ such that $`2\overline{c}_i\psi _m>1`$ (if enough knowledge is available to do so) and $`\sigma _{\alpha _i^j}>1`$. However, having large $`\sigma _{\alpha _i^j}`$ makes $`W_d`$ larger; this can be offset, however, by also choosing the ratio $`\sigma _{\alpha _i^j}/\gamma _{\alpha _i^j}<1`$ or smaller. Finally, it is clear that making $`k_i`$ larger reduces the effects of the representation errors, and therefore makes $`W_d`$ smaller. Observe that there is enough design freedom to make $`W_d`$ small and $`\beta _d`$ large independently of each other.
These simple guidelines may become very useful when performing a real control design. Moreover, notice that the bound on $`_{i=1}^nz_i^2(t)`$ makes it possible to specify the compact sets of the approximators so that, even throughout the transient, it can be guaranteed that the states will remain within the compact sets without the need for a global bounding control term. This has been a recurrent shortcoming of many on-line function approximation based methods, and the explicit bound on the transient makes it possible to overcome it.
## 3 Conclusions
In this paper we have developed a direct adaptive control method for a class of uncertain nonlinear systems with a time-varying structure using a Lyapunov approach to construct the stability proofs. The systems we consider are composed of a finite number of “pieces,” or dynamic subsystems, which are interpolated by functions that depend on a possibly exogenous scheduling variable. We assume that each piece is in strict feedback form, and show that the methods yield stability of all signals in the closed-loop, as well as convergence of the state vector to a residual set around the equilibrium, whose size can be set by the choice of several design parameters
We argue that the direct adaptive method presents several advantages over indirect methods in general, including the need for a smaller amount of information about the plant and a simpler design. Finally, we provide design guidelines based on $`_{\mathrm{}}`$ bounds on the transient and argue that this bound makes it possible to precisely determine how large the compact sets for the function approximators should be so that the states do not exit them.
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# 1 Introduction
## 1 Introduction
There are two major problems in the QCD analysis of polarized inclusive deep inelastic scattering (DIS): i) the absence of neutrino data makes it impossible, in principle, to determine the non-strange polarized sea-quark densities $`\mathrm{\Delta }\overline{u}(x,Q^2)`$ and $`\mathrm{\Delta }\overline{d}(x,Q^2)`$, ii) the separate determination of the polarized strange quark density $`\mathrm{\Delta }s(x,Q^2)`$ and the polarized gluon density $`\mathrm{\Delta }G(x,Q^2)`$ relies heavily on the QCD evolution in $`Q^2`$ and use of the flavour $`SU(3)_F`$-invariance relation
$`{\displaystyle _0^1}𝑑x\left[\mathrm{\Delta }u+\mathrm{\Delta }\overline{u}+\mathrm{\Delta }d+\mathrm{\Delta }\overline{d}2\mathrm{\Delta }s2\mathrm{\Delta }\overline{s}\right]=3FD.`$ (1)
The absence of a long lever arm in $`Q^2`$, in the polarized case, and doubt concerning the reliability of $`SU(3)_F`$-invariance for hyperon $`\beta `$-decay means that $`\mathrm{\Delta }s`$ and $`\mathrm{\Delta }G`$ are still rather poorly known , despite the dramatic improvement in the quality of the data in the past few years.
The direct resolution of these problems must await a series of new machine development projects, based on very high intensity neutrino beams, which are most unlikely to come into operation before the year 2015.
In the meantime there is currently a major experimental effort at CERN , HERA and Jefferson Lab. to study semi-inclusive polarized DIS reactions
$`\stackrel{}{e}+\stackrel{}{N}e+h+X,`$ (2)
involving the detection of the produced hadron $`h`$.
The theoretical structure for the analysis of such reactions, in both leading order QCD (LO) and next-to-leading order (NLO), exists . However we are critical of the type of LO analysis carried out thus far by the experimental groups.
The analysis of semi-inclusive DIS involves both parton densities and fragmentation functions. In the LO treatments referred to above, the fragmentation functions (FF’s) are treated as known quantities from inclusive $`e^+e^{}hadrons`$, and are used in constructing an auxiliary quantity called ”purity”. However, it is well known that in $`e^+e^{}hadrons`$, both in LO and NLO, first, only the combinations $`D_q^h+D_{\overline{q}}^h`$ can be determined while for the analyses of semi-inclusive DIS both $`D_q^h`$ and $`D_{\overline{q}}^h`$ are needed separately, and second, the existing analysis are rather ambiguous: a detailed study in makes a 31 parameter fit to the data, and no errors are quoted, and, in a more recent study , the FF’s differ significantly from those in , by 40% or more in some regions of $`z`$. Under these circumstances it is unreasonable to pretend to have an absolute knowledge of the fragmentation functions.
Two NLO analyses based upon a global analysis of the inclusive and semi-inclusive data have been attempted , . In the more recent analysis the authors relax the equality of $`\mathrm{\Delta }\overline{u}(x)`$ and $`\mathrm{\Delta }\overline{d}(x)`$ and find a preference for a positive $`\mathrm{\Delta }\overline{u}(x)`$, but effectively no constraint on the sign of $`\mathrm{\Delta }\overline{d}(x)`$. However, in both these analyses it is again assumed that the FF’s are known exactly: those of being used in , and those of in .
In addition, in the above mentioned analysis, some simplifying assumptions are made about relations between various polarized parton densities. In the following, except where expressly indicated, we make no assumptions at all concerning the polarized or unpolarized parton densities. Indeed there are persuasive arguments from the large-$`N_c`$ limit of QCD that a significant difference should exist between $`\mathrm{\Delta }\overline{u}(x)`$ and $`\mathrm{\Delta }\overline{d}(x)`$ with $`|\mathrm{\Delta }\overline{u}\mathrm{\Delta }\overline{d}|>|\overline{u}\overline{d}|`$, and it has been argued that such a situation is compatible with all present day data . Further, bearing in mind recent arguments that $`s(x)\overline{s}(x)`$ and $`\mathrm{\Delta }s(x)\mathrm{\Delta }\overline{s}(x)`$, we even refrain from the very common assumption of the equality of these densities.
However, from experience gained in the analysis of inclusive polarized DIS, it appears that the parameter space is sufficiently complicated to be able to produce biases in the $`\chi ^2`$ analysis, which can lead to unphysical results. We thus believe it to be dangerous in either LO or NLO QCD, to put together all inclusive and semi-inclusive data in one global analysis. Rather, what is required, is a working strategy, making optimal use of selected parts of the data.
The aim of this paper is precisely to provide a strategy, in both LO and NLO QCD, for the analysis of the semi-inclusive data. We proceed as follows: Information about the FF’s is obtained from unpolarized semi-inclusive data. Then this information is used to determine the polarized parton densities from polarized semi-inclusive DIS. Appropriate use of the information from DIS and $`e^+e^{}hadrons`$ is made as well. We suggest, for example, that one should use as input not just a knowledge of the unpolarized parton densities and their errors, but rather the polarized isotopic combination
$`\mathrm{\Delta }q_3(x,Q^2)=(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u})(\mathrm{\Delta }d+\mathrm{\Delta }\overline{d})`$ (3)
which is by now very well determined from polarized inclusive data, and which is free from any influence of $`\mathrm{\Delta }s`$ and $`\mathrm{\Delta }G`$. Inclusive $`e^+e^{}hadrons`$ and DIS are used also to considerably simplify the NLO expressions for the cross sections of (2).
Further, in stead of dealing with each parton density and FF separately, as is often done we work with their singlet and non-singlet combinations. As both the parton densities and the FF’s enter the semi-inclusive cross sections, we single out such observables that are singlet (non-singlet) combinations on both the parton densities and the FF’s.
This leads to the fact that we often consider linear combinations of experimentally measured quantities, and it may be objected that thereby we are dealing with experimental observables with possibly large errors. It has to be understood that that is a reflection of the true situation and not an artifact of our approach. To take an absolutely trivial example, suppose $`E_1`$ and $`E_2`$ are two experimentally measured functions used in an attempt to determine the theoretical functions $`T_1`$ and $`T_2`$, where $`E_1=T_1+T_2`$ and $`E_2=T_1T_2`$. Now, if it happens that $`E_1E_2`$ and if we write $`T_{1,2}=(1/2)(E_1\pm E_2)`$ then $`T_2`$ will be very poorly determined. This is unavoidable. It does not help to do a best fit to $`E_1`$ and $`E_2`$ with some parametrisations of $`T_1`$ and $`T_2`$. Inherently the result for $`T_2`$ will have a large relative error.
Thus we believe that any relatively large errors occurring in our manipulation of the experimental quantities reflects a genuine and unavoidable imprecision in the determination of certain theoretical quantities.
In order to minimize systematic errors experimentalists prefer to consider asymmetries or ratios of cross sections where the detection efficiencies should roughly cancel out, e.g. ratios of polarized to unpolarized DIS or ratios of polarized to unpolarized semi-inclusive for a given detected hadron. We appreciate that this is a fact of experimental life. However we wish to stress that a large amount of information is lost in restricting oneself to only these ratios. It is vitally important to gain control of the systematic errors in detection efficiencies, and although we try as far as possible to deal with the favoured kind of ratios we will be forced also to consider other types of cross-section ratios.
Throughout the rest of this paper we assume that a kinematic separation is possible between hadrons produced in the current fragmentation region and those produced from the target remnants. We consider only the current fragmentation region so that our formulae apply only to this region and fracture functions play no role in our discussion.
It has to be stressed that there is a huge difference in complexity between the LO and NLO treatments. Thus it makes sense to utilize the LO approach, provided appropriate checks, (which we suggest) are carried out.
Our analysis proceeds in a step-by-step fashion. Firstly we describe a generic test for the reliability of a LO treatment. Assuming this to be successful we present the analysis in LO. In LO the information on the polarized valence quark densities $`\mathrm{\Delta }u_V`$ and $`\mathrm{\Delta }d_V`$ and on the breaking of SU(2) invariance of the polarized sea is obtained without any knowledge of the FF’s. However, it should be clear that any information on the sea quark densities in LO cannot be reliable, since they are expected to be small and thus comparable to the NLO corrections. We consider also the experimentally difficult case of $`\varphi `$ production, which seems to be the best way to get an accurate determination of $`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}`$ in LO. We discuss also how, in principle, one can test $`s(x)=\overline{s}(x)`$ and/or $`\mathrm{\Delta }s(x)=\mathrm{\Delta }\overline{s}(x)`$.
In the NLO treatment we show in a new way how information on $`e^+e^{}h+X`$ and inclusive DIS can be incorporated directly so as to considerably simplify the NLO expressions for semi-inclusive cross sections. This key result is presented in eqs. (69) and (70). Next we discuss the fragmentation combination $`D_u^{\pi ^+}D_u^\pi ^{}`$ and use it to evaluate $`\mathrm{\Delta }u_V`$ and $`\mathrm{\Delta }d_V`$ in NLO. Then we obtain expressions for $`D_u^{\pi ^+}+D_u^\pi ^{}`$, for $`D_G^{\pi ^+}`$ and for $`D_s^{\pi ^+}+D_s^\pi ^{}`$. With these we are able to obtain ($`\mathrm{\Delta }\overline{u}\mathrm{\Delta }\overline{d}`$) in NLO. Finally we have a set of 2 equations involving 3 unknown functions, ($`\mathrm{\Delta }u+\mathrm{\Delta }\overline{u}+\mathrm{\Delta }d+\mathrm{\Delta }\overline{d}`$), ($`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}`$) and $`\mathrm{\Delta }G`$. An accurate determination of all three functions would require data over a presently impossibly wide range of $`Q^2`$. We thus suggest using here for $`\mathrm{\Delta }G`$ its determination from charm production . It should then finally be possible to get an accurate assessment of $`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}`$ in NLO. Lastly we consider the evaluation of $`s(x)\overline{s}(x)`$ and $`\mathrm{\Delta }s(x)\mathrm{\Delta }\overline{s}(x)`$.
## 2 A strategy for semi-inclusive DIS
In NLO QCD, the expressions for semi-inclusive cross sections involve convolutions of parton densities and fragmentation functions with (known) Wilson coefficients . Our lack of knowledge of the errors on the FF’s will make it difficult to assess the accuracy of the parton densities which we are trying to determine. In the LO QCD approximation, on the other hand there are no convolutions, but simple products only (independent fragmentation), and it becomes possible to construct measurable combinations of cross sections in which the FF’s completely cancel out . However it is not clear how reliable the LO is.
We believe it is quite safe when determining the large $`\mathrm{\Delta }u`$ and $`\mathrm{\Delta }d`$ densities, but could be quite misleading for $`\mathrm{\Delta }\overline{u}`$, $`\mathrm{\Delta }\overline{d}`$ and $`\mathrm{\Delta }s`$. In any event it is absolutely essential to test independent fragmentation in order to have a feeling for the errors on parton densities obtained via the LO formalism.
In LO, the structure of the expressions is generally of the form
$`partondensity\mathrm{\Delta }q(x,Q^2)=experimentalobservableE(x,z,Q^2)`$ (4)
or
$`fragmentationfunctionD(z,Q^2)=experimentalobservableE(x,z,Q^2).`$ (5)
In both cases the characteristic feature of the LO treatment is that the RH sides, which can in principle depend on $`(x,z,Q^2)`$, should only depend on two of these, either $`(x,Q^2)`$ or $`(z,Q^2)`$ respectively, so that there is an independence of the third variable, which we shall call the passive variable.
Every expression of the form (4) or (5) should be tested for dependence on the passive variable. If a significant dependence is found it does not mean that the LO analysis must be abandoned, but it suggests that the variation with the passive variable be used as an estimate of the theoretical errors, $`\delta _{TH}[\mathrm{\Delta }q(x,Q^2)]`$, $`\delta _{TH}[D(z,Q^2)]`$ on the sought for quantities.
In Section 7 we discuss a strategy for the analysis in NLO.
## 3 Parton densities in LO QCD
It is useful to introduce the following notation for semi-inclusive processes:
$`\stackrel{~}{\sigma }^h{\displaystyle \frac{x(P+l)^2}{4\pi \alpha ^2}}\left({\displaystyle \frac{2y^2}{1+(1y)^2}}\right){\displaystyle \frac{d^3\sigma ^h}{dxdydz}}`$ (6)
and
$`\mathrm{\Delta }\stackrel{~}{\sigma }^h{\displaystyle \frac{x(P+l)^2}{4\pi \alpha ^2}}\left({\displaystyle \frac{y}{2y}}\right)\left[{\displaystyle \frac{d^3\sigma _{++}^h}{dxdydz}}{\displaystyle \frac{d^3\sigma _+^h}{dxdydz}}\right]`$ (7)
where $`P^\mu `$ and $`l^\mu `$ are the nucleon and lepton four momenta, and $`\sigma _{\lambda \mu }`$ refer to a lepton of helicity $`\lambda `$ and a nucleon of helicity $`\mu `$. The variables $`x`$, $`y`$, $`z`$ are the usual DIS kinematic variables.
Then in LO the cross sections for the semi-inclusive production of a hadron $`h`$ have the simple $`y`$-independent form
$`\mathrm{\Delta }\stackrel{~}{\sigma }^h(x,z,Q^2)`$ $`=`$ $`{\displaystyle \underset{q,\overline{q}}{}}e_q^2\mathrm{\Delta }q_i(x,Q^2)D_i^h(z,Q^2)`$ (8)
$`\stackrel{~}{\sigma }^h(x,z,Q^2)`$ $`=`$ $`{\displaystyle \underset{q,\overline{q}}{}}e_q^2q_i(x,Q^2)D_i^h(z,Q^2),`$ (9)
where the sum is over quarks and aniquarks, and where $`D_i^h`$ is the fragmentation function for quark or antiquark $`i`$ to produce $`h`$.
We consider sum and difference cross sections for producing $`h`$ and its charge conjugate $`\overline{h}`$ on both protons and neutrons, and define
$`\mathrm{\Delta }A_{p,n}^{h\pm \overline{h}}(x,z,Q^2)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }\stackrel{~}{\sigma }_{p,n}^h\pm \mathrm{\Delta }\stackrel{~}{\sigma }_{p,n}^{\overline{h}}}{\stackrel{~}{\sigma }_{p,n}^h\pm \stackrel{~}{\sigma }_{p,n}^{\overline{h}}}}{\displaystyle \frac{\mathrm{\Delta }\stackrel{~}{\sigma }_{p,n}^{h\pm \overline{h}}}{\stackrel{~}{\sigma }_{p,n}^{h\pm \overline{h}}}},`$ (10)
$`\mathrm{\Delta }A_{p\pm n}^{h\pm \overline{h}}(x,z,Q^2)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }\stackrel{~}{\sigma }_p^{h\pm \overline{h}}\pm \mathrm{\Delta }\stackrel{~}{\sigma }_n^{h\pm \overline{h}}}{\stackrel{~}{\sigma }_p^{h\pm \overline{h}}\pm \stackrel{~}{\sigma }_n^{h\pm \overline{h}}}}.`$ (11)
For inclusive unpolarized and polarized DIS cross sections we use the notation:
$`\stackrel{~}{\sigma }^{DIS}{\displaystyle \frac{x(P+l)^2}{4\pi \alpha ^2}}\left({\displaystyle \frac{2y^2}{1+(1y)^2}}\right){\displaystyle \frac{d^2\sigma ^{DIS}}{dxdy}}`$ (12)
and
$`\mathrm{\Delta }\stackrel{~}{\sigma }^{DIS}{\displaystyle \frac{x(P+l)^2}{4\pi \alpha ^2}}\left({\displaystyle \frac{y}{2y}}\right)\left[{\displaystyle \frac{d^3\sigma _{++}^{DIS}}{dxdy}}{\displaystyle \frac{d^3\sigma _+^{DIS}}{dxdy}}\right].`$ (13)
Then in LO we have the $`y`$-independent expressions:
$`\stackrel{~}{\sigma }^{DIS}(x,Q^2)`$ $`=`$ $`2F_1^N(x,Q^2)|_{LO}={\displaystyle \underset{q,\overline{q}}{}}e_q^2q_i(x,Q^2)`$ (14)
$`\mathrm{\Delta }\stackrel{~}{\sigma }^{DIS}(x,Q^2)`$ $`=`$ $`2g_1^N(x,Q^2)|_{LO}={\displaystyle \underset{q,\overline{q}}{}}e_q^2\mathrm{\Delta }q_i(x,Q^2).`$ (15)
In addition to (10) and (11) we consider the ratios of sum and difference hadron yields for the unpolarized semi-inclusive and inclusive processes:
$`R_{p,n}^{h\pm \overline{h}}(x,z,Q^2)={\displaystyle \frac{\stackrel{~}{\sigma }_{p,n}^h\pm \stackrel{~}{\sigma }_{p,n}^{\overline{h}}}{\stackrel{~}{\sigma }_{p,n}^{DIS}}},R_{p\pm n}^{h\pm \overline{h}}(x,z,Q^2)={\displaystyle \frac{\stackrel{~}{\sigma }_p^{h\pm \overline{h}}\pm \stackrel{~}{\sigma }_n^{h\pm \overline{h}}}{\stackrel{~}{\sigma }_p^{DIS}\pm \stackrel{~}{\sigma }_n^{DIS}}}.`$ (16)
(It is equally good to use a sum over any set of hadrons $`h`$ and their charge conjugate $`\overline{h}`$.)
### 3.1 Testing LO QCD
Using only charge conjugation invariance it is easy to show that
$`\mathrm{\Delta }A_{pn}^{h+\overline{h}}(x,z,Q^2)={\displaystyle \frac{\mathrm{\Delta }\stackrel{~}{\sigma }_p^{DIS}\mathrm{\Delta }\stackrel{~}{\sigma }_n^{DIS}}{\stackrel{~}{\sigma }_p^{DIS}\stackrel{~}{\sigma }_n^{DIS}}}(x,Q^2)={\displaystyle \frac{(g_1^pg_1^n)|_{LO}}{(F_1^pF_1^n)|_{LO}}}(x,Q^2).`$ (17)
This is a key relation for testing the reliability of the LO. The RHS is completely known from inclusive DIS and is independent of $`z`$. The LHS, in principle, depends also upon $`z`$ and upon the hadron $`h`$. Only in LO (or in the simple parton model) should it be independent of $`z`$ and of $`h`$. It is thus crucial to test this feature.
To help with statistics it is also possible to formulate an integrated version of (17). This is given in . For the rest of this section we assume that the test (17) has been successful and proceed with the analysis in LO.
### 3.2 The valence quark densities in LO
The polarized valence quark densities can be obtained from $`\pi ^\pm `$ production, assuming only isospin invariance:
$`\mathrm{\Delta }u_V`$ $`=`$ $`{\displaystyle \frac{1}{15}}\left\{4(4u_Vd_V)\mathrm{\Delta }A_p^{\pi ^+\pi ^{}}+(4d_Vu_V)\mathrm{\Delta }A_n^{\pi ^+\pi ^{}}\right\}`$
$`\mathrm{\Delta }d_V`$ $`=`$ $`{\displaystyle \frac{1}{15}}\left\{4(4d_Vu_V)\mathrm{\Delta }A_n^{\pi ^+\pi ^{}}+(4u_Vd_V)\mathrm{\Delta }A_p^{\pi ^+\pi ^{}}\right\}.`$ (18)
For the case of $`K^\pm `$ or $`\mathrm{\Lambda },\overline{\mathrm{\Lambda }}`$ production, if one makes also the conventional assumption that $`s=\overline{s}`$ and $`\mathrm{\Delta }s=\mathrm{\Delta }\overline{s}`$ one has in addition:
$`\mathrm{\Delta }u_V`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{(u_V+d_V)\mathrm{\Delta }A_{p+n}^{K^+K^{}}+(u_Vd_V)\mathrm{\Delta }A_{pn}^{K^+K^{}}\right\}`$
$`\mathrm{\Delta }d_V`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{(u_V+d_V)\mathrm{\Delta }A_{p+n}^{K^+K^{}}(u_Vd_V)\mathrm{\Delta }A_{pn}^{K^+K^{}}\right\}.`$ (19)
Note that we have not assumed $`D_d^{K^+}=D_d^K^{}`$, although that is suggested by the absence of a $`d`$ quark in the leading Fock state of $`K^\pm `$. Indeed the above equality can be tested (see Section 5.2).
For isoscalar hadrons like $`\mathrm{\Lambda },\overline{\mathrm{\Lambda }}`$ again assuming $`s=\overline{s}`$ and $`\mathrm{\Delta }s=\mathrm{\Delta }\overline{s}`$, one finds
$`\mathrm{\Delta }u_V`$ $`=`$ $`{\displaystyle \frac{1}{15}}\left\{4(4u_V+d_V)\mathrm{\Delta }A_p^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}(4d_V+u_V)\mathrm{\Delta }A_n^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}\right\}`$
$`\mathrm{\Delta }d_V`$ $`=`$ $`{\displaystyle \frac{1}{15}}\left\{4(4d_V+u_V)\mathrm{\Delta }A_n^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}(4u_V+d_V)\mathrm{\Delta }A_p^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}\right\}.`$ (20)
We shall comment in Section 5.5 on the situation if one does not assume $`s=\overline{s}`$ and $`\mathrm{\Delta }s=\mathrm{\Delta }\overline{s}`$, where we suggest a method for estimating if the failure of these equalities is serious or not. In any event, the safe way to obtain $`\mathrm{\Delta }u_V`$ and $`\mathrm{\Delta }d_V`$ is via $`\pi ^\pm `$ production, eq. (18).
Once again the reliability of these LO equations can be tested by checking that the RH sides of (18), (19) and (20) do not depend on $`z`$. If it is found that for a given $`x`$-bin the RH sides vary with $`z`$ by some amount $`\delta _{TH}[\mathrm{\Delta }u_V]`$, $`\delta _{TH}[\mathrm{\Delta }d_V]`$, then these could be regarded as an estimate of the theoretical error at this $`x`$ value.
## 4 Use of $`\mathrm{\Delta }q_3(x,Q^2)`$ from polarized inclusive DIS
At NLO the spin dependent structure functions for protons and neutrons are given by:
$`g_1^p(x,Q^2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{q=u,d,s}{}}e_q^2[(\mathrm{\Delta }q+\mathrm{\Delta }\overline{q})(1+{\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}\delta C_q)+`$ (21)
$`+{\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}\mathrm{\Delta }G\delta C_G]`$
involving a convolution of polarized parton densities with known Wilson coefficients. For the neutron, $`g_1^n`$ is obtained by the replacement
$`\mathrm{\Delta }u\mathrm{\Delta }d`$ (22)
in $`g_1^p`$.
Now it is clear that one can obtain information only on the combinations:
$`\mathrm{\Delta }u+\mathrm{\Delta }\overline{u},\mathrm{\Delta }d+\mathrm{\Delta }\overline{d},\mathrm{\Delta }s+\mathrm{\Delta }\overline{s},\mathrm{\Delta }G.`$ (23)
and that it is impossible to obtain separate information on the valence and non-strange sea quark polarizations from inclusive, neutral current, polarized DIS. <sup>3</sup><sup>3</sup>3However, sometimes it is convenient to parametrize separately the valence and sea quarks and to assume, for example, $`\mathrm{\Delta }\overline{u}=\mathrm{\Delta }\overline{d}=\lambda \mathrm{\Delta }s`$, where $`\lambda `$ is a free parameter. It should be obvious then, that any claim that the $`\chi ^2`$ analysis favors some particular value of $`\lambda `$ must be fictitious and a consequence of some hidden bias in the minimization procedure. Yet such claims have been made.
In our semi-inclusive analysis we shall use as a known qauntity only the non-singlet combination of the polarized parton densities $`\mathrm{\Delta }q_3`$, eq.(3). Now that there is such an improvement in the quality of the neutron data we believe this quantity is very well constrained directly by the inclusive data. For one has, from (21)
$`g_1^p(x,Q^2)g_1^n(x,Q^2)={\displaystyle \frac{1}{6}}\mathrm{\Delta }q_3\left(1+{\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}\delta C_q\right)`$ (24)
and $`\mathrm{\Delta }q_3(x,Q^2)`$ is determined without any influence from the less well known quantities $`\mathrm{\Delta }s`$ and $`\mathrm{\Delta }G`$, either in LO or in NLO. Of course if the semi-inclusive analysis is done in LO one must use $`\mathrm{\Delta }q_3(x,Q^2)|_{LO}`$. Note that we do not use information on $`(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u})`$ or $`(\mathrm{\Delta }d+\mathrm{\Delta }\overline{d})`$ from inclusive DIS, since these are subject to the less known strange quark and gluon effects.
### 4.1 SU(2) symmetry of the sea quark densities in LO
One has
$`[\mathrm{\Delta }\overline{u}(x,Q^2)\mathrm{\Delta }\overline{d}(x,Q^2)]_{LO}={\displaystyle \frac{1}{2}}\left[\mathrm{\Delta }q_3(x,Q^2)+\mathrm{\Delta }d_V(x,Q^2)\mathrm{\Delta }u_V(x,Q^2)\right]_{LO}.`$ (25)
Eq. (25) determines the SU(2) symmetry breaking of the polarized sea without requiring any knowledge of the unknown $`\mathrm{\Delta }\overline{q}`$ and $`\mathrm{\Delta }G`$. A possible test for SU(2) breaking for the polarized sea densities that does not require any knowledge of the polarized densities was given in .
In order to determine the polarized sea quark densities separately we need one more relation, namely the value of
$`\mathrm{\Delta }q_+(x,Q^2)\mathrm{\Delta }u+\mathrm{\Delta }\overline{u}+\mathrm{\Delta }d+\mathrm{\Delta }\overline{d}.`$ (26)
For then
$`(\mathrm{\Delta }\overline{u}+\mathrm{\Delta }\overline{d})_{LO}={\displaystyle \frac{1}{2}}(\mathrm{\Delta }q_+\mathrm{\Delta }u_V\mathrm{\Delta }d_V)_{LO}`$ (27)
which combined with (25) yields $`\mathrm{\Delta }\overline{u}`$ and $`\mathrm{\Delta }\overline{d}`$ separately.
Although each term on the RHS of (25) and (27) should be well determined in LO, the corresponding linear combinations are expected to be small and may thus be very sensitive to NLO corrections. An indication of the sensitivity may be inferred from the fact that in inclusive polarized DIS, $`\mathrm{\Delta }s(x,Q^2)`$ changes by roughly a factor of 2 in going from LO to NLO or when one changes factorisation schemes from $`\overline{MS}`$ to $`AB`$ or $`JET`$.
Note that determining, say, $`\mathrm{\Delta }\overline{u}`$ via $`\mathrm{\Delta }\overline{u}=\frac{1}{2}[(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u})\mathrm{\Delta }u_V]`$ is unreliable since determination of ($`\mathrm{\Delta }u+\mathrm{\Delta }\overline{u}`$) from inclusive DIS requires a knowledge of $`\mathrm{\Delta }s`$.
In order to determine $`\mathrm{\Delta }q_+(x,Q^2)`$ it will first be necessary to extract some information on the FF’s.
### 4.2 Fragmentation functions in LO
1. From measurements of the ratios $`R_{pn}^{h+\overline{h}}`$ of the semi-inclusive to inclusive DIS cross sections on protons and neutrons for any given hadron $`h`$, it is feasible in LO to learn a great deal about the FF’s $`D_q^h+D_q^{\overline{h}}D_q^{h+\overline{h}}`$. This combination is measured also in $`e^+e^{}hadrons`$.
Analogous to the polarized case, we define
$`q_3(x,Q^2)`$ $`=`$ $`u(x,Q^2)+\overline{u}(x,Q^2)d(x,Q^2)\overline{d}(x,Q^2)`$ (28)
$`q_+(x,Q^2)`$ $`=`$ $`u(x,Q^2)+\overline{u}(x,Q^2)+d(x,Q^2)+\overline{d}(x,Q^2).`$ (29)
which are well determined from inclusive DIS data and which thus can be taken as known quantities in the semi-inclusive analysis.
* Using data on unpolarized semi-inclusive DIS we have, in LO,
$`R_{pn}^{h+\overline{h}}={\displaystyle \frac{2}{3}}\left[4D_u^{h+\overline{h}}(z,Q^2)D_d^{h+\overline{h}}(z,Q^2)\right].`$ (30)
* For the case of pions, kaons and $`\mathrm{\Lambda }`$, when SU(2) invariance can be used this simplifies to
$`R_{pn}^{\pi ^++\pi ^{}}=2D_u^{\pi ^++\pi ^{}}(z,Q^2)=2D_d^{\pi ^++\pi ^{}}(z,Q^2).`$ (31)
* For $`K`$ mesons and $`\mathrm{\Lambda }`$ hyperons,
$`R_{pn}^{K,\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}}=2D_u^{K,\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}}(z,Q^2)=2D_d^{K,\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}}(z,Q^2),`$ (32)
where the superscript $`K`$ stands for the sum over all produced kaons:
$`KK^++K^{}+K^0+\overline{K}^0.`$ (33)
It will also be of great interest to compare these FF’s with those obtained from $`e^+e^{}hadrons`$ and those used in Monte Carlo models.
2. Given that $`u_V(x,Q^2)`$ and $`d_V(x,Q^2)`$ are well determined from inclusive DIS one can proceed further to obtain the other combinations of FF’s $`D_q^hD_q^{\overline{h}}D_q^{h\overline{h}}`$.
* One finds for $`\pi ^\pm `$
$`D_u^{\pi ^+\pi ^{}}`$ $`=`$ $`{\displaystyle \frac{9\stackrel{~}{\sigma }_p^{\pi ^+\pi ^{}}}{4u_Vd_V}}={\displaystyle \frac{18\left(F_1^p\right)_{LO}R_p^{\pi ^+\pi ^{}}}{4u_Vd_V}}`$ (34)
$`=`$ $`{\displaystyle \frac{9\stackrel{~}{\sigma }_n^{\pi ^+\pi ^{}}}{4d_Vu_V}}={\displaystyle \frac{18\left(F_1^n\right)_{LO}R_n^{\pi ^+\pi ^{}}}{4d_Vu_V}}`$
Combined with (31) we have expressions for $`D_u^{\pi ^+}`$ and $`D_u^\pi ^{}`$ separately.
* For $`K^\pm `$ one obtains
$`4D_u^{K^+K^{}}D_d^{K^+K^{}}`$ $`=`$ $`{\displaystyle \frac{9\left\{\stackrel{~}{\sigma }_p^{K^+K^{}}\stackrel{~}{\sigma }_n^{K^+K^{}}\right\}}{u_Vd_V}}`$ (35)
$`=`$ $`{\displaystyle \frac{18\left[\left(F_1^p\right)_{LO}R_p^{K^+K^{}}\left(F_1^n\right)_{LO}R_n^{K^+K^{}}\right]}{u_Vd_V}}`$
It is usually assumed, and this is presumably a very good approximation, that
$`D_d^{K^+}=D_d^K^{}`$ (36)
in which case (35) can be read as an expression for $`D_u^{K^+K^{}}`$. It is not possible to test relation (36) without taking $`s=\overline{s}`$ and/or $`\mathrm{\Delta }s=\mathrm{\Delta }\overline{s}`$, but such an approach is hard to justify given that any failure of (36) is presumably very small.
* For isoscalar hadrons like $`\mathrm{\Lambda },\overline{\mathrm{\Lambda }}`$
$`D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}=D_d^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}={\displaystyle \frac{6\left[\left(F_1^p\right)_{LO}R_p^{\mathrm{\Lambda }^+\mathrm{\Lambda }^{}}\left(F_1^n\right)_{LO}R_n^{\mathrm{\Lambda }^+\mathrm{\Lambda }^{}}\right]}{u_Vd_V}}.`$ (37)
Given that $`\mathrm{\Delta }u_V`$ and $`\mathrm{\Delta }d_V`$ are determined via (18) we can write analogous expression for the above $`D_q^h`$ functions using the polarized data.
Of course all the expressions (30), (31), (32), (34), (35) and (37), being LO results, must be tested by demonstrating that the RH sides are essentially independent of the passive variable $`x`$.
Now that we have determined $`D_u^{\pi ^++\pi ^{}}`$ we can determine $`D_s^{\pi ^++\pi ^{}}`$ in LO via
$`D_s^{\pi ^++\pi ^{}}={\displaystyle \frac{9\left(F_1^p+F_1^n\right)_{LO}R_{p+n}^{\pi ^++\pi ^{}}5q_+D_u^{\pi ^++\pi ^{}}}{2(s+\overline{s})}}.`$ (38)
Similarly, since $`D_u^{\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}}`$ is determined via (32), we can find $`D_s^{\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}}`$ from
$`D_s^{\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}}={\displaystyle \frac{9\left(F_1^p+F_1^n\right)_{LO}R_{p+n}^{\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}}5q_+D_u^{\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}}}{2(s+\overline{s})}}.`$ (39)
### 4.3 The non-strange sea quark densities revisited, in LO
Now that we have determined $`D_u^{\pi ^++\pi ^{}}`$ and $`D_s^{\pi ^++\pi ^{}}`$ in LO we can, in principle, determine $`\mathrm{\Delta }q_+(x,Q^2)`$ and $`\mathrm{\Delta }s(x,Q^2)+\mathrm{\Delta }\overline{s}(x,Q^2)`$ from the semi-inclusive and inclusive relations, in LO,
$`g_1^p+g_1^n`$ $`=`$ $`{\displaystyle \frac{5\mathrm{\Delta }q_++2(\mathrm{\Delta }s+\mathrm{\Delta }\overline{s})}{18}}`$ (40)
$`\mathrm{\Delta }A_{p+n}^{\pi ^++\pi ^{}}`$ $`=`$ $`{\displaystyle \frac{5\mathrm{\Delta }q_+D_u^{\pi ^++\pi ^{}}+2(\mathrm{\Delta }s+\mathrm{\Delta }\overline{s})D_s^{\pi ^++\pi ^{}}}{5q_+D_u^{\pi ^++\pi ^{}}+2(s+\overline{s})D_s^{\pi ^++\pi ^{}}}}.`$ (41)
Such a determination of $`\mathrm{\Delta }q_+`$ in LO is likely to be reliable, but $`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}`$ and the individual $`\mathrm{\Delta }\overline{u}`$ and $`\mathrm{\Delta }\overline{d}`$ obtained in LO via (27) may be subject to significant uncertainty.
We note that there is an alternative way to determine $`\mathrm{\Delta }q_+`$, but it requires the ability to detect $`K^0`$. In that case one has, in LO,
$`\mathrm{\Delta }q_+=q_+{\displaystyle \frac{\mathrm{\Delta }\stackrel{~}{\sigma }_p^{K^++K^{}(K^0+\overline{K}^0)}+\mathrm{\Delta }\stackrel{~}{\sigma }_n^{K^++K^{}(K^0+\overline{K}^0)}}{\stackrel{~}{\sigma }_p^{K^++K^{}(K^0+\overline{K}^0)}+\stackrel{~}{\sigma }_n^{K^++K^{}(K^0+\overline{K}^0)}}}.`$ (42)
### 4.4 The strange quark density $`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}`$ in LO
The approach to $`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}`$ in Section 5.3 and the approach discussed in are unlikely to be reliable. The problem is that for production of pions the strange quark contribution is ”doubly small”, since e.g. one must compare $`\mathrm{\Delta }uD_u^\pi `$ with $`\mathrm{\Delta }sD_s^\pi `$ in which both $`|\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}||\mathrm{\Delta }u|`$ and $`|D_s^\pi ||D_u^\pi |`$. For kaons and $`\mathrm{\Lambda }`$ hyperons it is somewhat better in that $`|D_u^K||D_s^K|`$ and $`|D_u^{\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}}|=|D_d^{\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}}||D_s^{\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}}|`$, but this is similar to the situation in inclusive DIS where we know that the LO determination of $`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}`$ is quite unreliable.
The only possibility we can see for a reasonable determination of $`\mathrm{\Delta }s`$ in LO is via $`\varphi `$ production. For in this case one has, $`|\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}||\mathrm{\Delta }u|`$ but presumably $`|D_s^\varphi ||D_u^\varphi |`$ so that the strange and non-strange quarks are on equal footing.
One has, by charge conjugation invariance $`D_s^\varphi =D_{\overline{s}}^\varphi `$, and it should be quite safe to take $`D_u^\varphi =D_{\overline{u}}^\varphi =D_d^\varphi =D_{\overline{d}}^\varphi `$. One then obtains in LO
$`{\displaystyle \frac{\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}}{s+\overline{s}}}={\displaystyle \frac{3\mathrm{\Delta }\stackrel{~}{\sigma }_{p+n}^\varphi 5\left(\frac{\mathrm{\Delta }q_+}{\mathrm{\Delta }q_3}\right)\mathrm{\Delta }\stackrel{~}{\sigma }_{pn}^\varphi }{3\stackrel{~}{\sigma }_{p+n}^\varphi 5\left(\frac{q_+}{q_3}\right)\stackrel{~}{\sigma }_{pn}^\varphi }}.`$ (43)
Moreover one has expressions for the fragmentation functions as well:
$`D_s^\varphi `$ $`=`$ $`{\displaystyle \frac{3}{2(s+\overline{s})}}\left\{3\stackrel{~}{\sigma }_{p+n}^\varphi 5\left({\displaystyle \frac{q_+}{q_3}}\right)\stackrel{~}{\sigma }_{pn}^\varphi \right\}`$ (44)
$`D_u^\varphi `$ $`=`$ $`{\displaystyle \frac{3\stackrel{~}{\sigma }_{pn}^\varphi }{q_3}}.`$ (45)
As always expressions (43) - (45) must be tested for non-dependence on the relevant passive variable.
Finally we note that the same eqs. (43) - (45) hold for $`K`$, $`\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}`$ and $`\pi `$-production if the superscript $`\varphi `$ is replaced by $`K`$, $`\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}`$ and $`\pi ^++\pi ^{}`$, respectively. And though the production rates are higher, the sensitivity to the strange quarks in the $`K`$, $`\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}`$ and $`\pi `$-productions is lower.
### 4.5 Concerning $`s=\overline{s}`$ and $`\mathrm{\Delta }s=\mathrm{\Delta }\overline{s}`$ in LO <sup>4</sup><sup>4</sup>4We are grateful to M. Anselmino, M. Boglione and U.D’Alesio for drawing our attention to this issue
In the analysis of DIS it is conventional to assume $`s=\overline{s}`$ and $`\mathrm{\Delta }s=\mathrm{\Delta }\overline{s}`$. However there are models and arguments which suggest that these equalities might not hold.
Within the limitations of the LO we can test these relationships via $`(K^+,K^{})`$ and $`(\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$ production. One has for the unpolarized case, assuming $`D_d^{K^+K^{}}=0`$ several different possibilities:
$`(s\overline{s})D_s^{K^+K^{}}`$ $`=`$ $`18\left(F_1^p\right)_{LO}R_p^{K^+K^{}}4u_VD_u^{K^+K^{}}`$ (46)
$`=`$ $`18\left(F_1^n\right)_{LO}R_n^{K^+K^{}}4d_VD_u^{K^+K^{}}`$
$`=`$ $`{\displaystyle \frac{9\left\{u_V\stackrel{~}{\sigma }_n^{K^+K^{}}d_V\stackrel{~}{\sigma }_p^{K^+K^{}}\right\}}{u_Vd_V}}.`$
Then for the polarized case, given that $`\mathrm{\Delta }u_V`$ and $`\mathrm{\Delta }d_V`$ are known from (18) and $`(s\overline{s})D_s^{K^+K^{}}`$ is determined, we can proceed to determine $`(\mathrm{\Delta }s\mathrm{\Delta }\overline{s})D_s^{K^+K^{}}`$ from $`\mathrm{\Delta }A_p^{K^+K^{}}`$ or $`\mathrm{\Delta }A_n^{K^+K^{}}`$:
$`\mathrm{\Delta }A_p^{K^+K^{}}={\displaystyle \frac{4\mathrm{\Delta }u_VD_u^{K^+K^{}}+(\mathrm{\Delta }s\mathrm{\Delta }\overline{s})D_s^{K^+K^{}}}{4u_VD_u^{K^+K^{}}+(s\overline{s})D_s^{K^+K^{}}}}`$
$`\mathrm{\Delta }A_n^{K^+K^{}}={\displaystyle \frac{4\mathrm{\Delta }d_VD_u^{K^+K^{}}+(\mathrm{\Delta }s\mathrm{\Delta }\overline{s})D_s^{K^+K^{}}}{4d_VD_u^{K^+K^{}}+(s\overline{s})D_s^{K^+K^{}}}},`$ (47)
$`D_u^{K^+K^{}}`$ is assumed to be known through (35).
For isoscalar hadrons like $`\mathrm{\Lambda },\overline{\mathrm{\Lambda }}`$ we have the possiblities
$`(s\overline{s})D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$ $`=`$ $`18\left(F_1^p\right)_{LO}R_p^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}(4u_V+d_V)D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$ (48)
$`=`$ $`18\left(F_1^n\right)_{LO}R_n^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}(4d_V+u_V)D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$
$`=`$ $`{\displaystyle \frac{3}{u_Vd_V}}\left\{\left(4u_V+d_V\right)\stackrel{~}{\sigma }_n^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}\left(4d_V+u_V\right)\stackrel{~}{\sigma }_p^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}\right\},`$
where $`D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$ is assumed to be determined in (37). Here and in eq. (46) it is clear that $`\stackrel{~}{\sigma }_N^{K^+K^{}}`$ and $`\stackrel{~}{\sigma }_N^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$ can be replaced by the corresponding experimentally measured quantities $`2\left(F_1^N\right)_{LO}R_N^{K^+K^{}}`$ or $`2\left(F_1^N\right)_{LO}R_N^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$. Then $`(\mathrm{\Delta }s\mathrm{\Delta }\overline{s})D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$ can be determined via $`\mathrm{\Delta }A_p^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$ or $`\mathrm{\Delta }A_n^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$:
$`\mathrm{\Delta }A_p^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}={\displaystyle \frac{(4\mathrm{\Delta }u_V+\mathrm{\Delta }d_V)D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}+(\mathrm{\Delta }s\mathrm{\Delta }\overline{s})D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}}{(4u_V+d_V)D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}+(s\overline{s})D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}}}`$
$`\mathrm{\Delta }A_n^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}={\displaystyle \frac{(4\mathrm{\Delta }d_V+\mathrm{\Delta }u_V)D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}+(\mathrm{\Delta }s\mathrm{\Delta }\overline{s})D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}}{(4d_V+u_V)D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}+(s\overline{s})D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}}}`$ (49)
Although (46), (47), (48) and (49), being LO expressions, cannot be expected to yield accurate values for $`s\overline{s}`$ and $`\mathrm{\Delta }s\mathrm{\Delta }\overline{s}`$, they should nonetheless enable one to say whether they are compatible with zero since $`D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$ should be relatively large. Note that $`s\overline{s}`$ and/or $`\mathrm{\Delta }s\mathrm{\Delta }\overline{s}`$ different from zero would break the independence of the RH sides of (47) and (49) on the passive variable $`z`$. Of particular interest is the speculation that $`\mathrm{\Delta }s\mathrm{\Delta }\overline{s}`$ but $`s\overline{s}`$, the consistency of which could perhaps be tested from (46), (47), (48) and (49).
## 5 Semi-inclusive analysis in NLO QCD
The situation in NLO is much more complicated than in LO, since factorization is replaced by convolution, and it is also more complicated than inclusive DIS in NLO since here one has to contend with double convolutions of the form $`qCD`$ and $`\mathrm{\Delta }q\mathrm{\Delta }CD`$ for the unpolarized and polarized cases respectively, where $`C`$ and $`\mathrm{\Delta }C`$ are Wilson coefficients first derived in and .
The double convolution is defined as
$`qCD={\displaystyle \frac{dx^{}}{x^{}}\frac{dz^{}}{z^{}}q\left(\frac{x}{x^{}}\right)C(x^{},z^{})D\left(\frac{z}{z^{}}\right)}`$ (50)
where the range of integration is given as follows:
* $`𝐈_\mathrm{𝟏}:`$ The range is
$`{\displaystyle \frac{x}{x+(1x)z}}x^{}1withzz^{}1`$ (51)
* $`𝐈_\mathrm{𝟐}:`$ In addition to (51) there is the range
$`xx^{}{\displaystyle \frac{x}{x+(1x)z}}with{\displaystyle \frac{x(1x^{})}{x^{}(1x)}}z^{}1.`$ (52)
Note, that contrary to the case of the usual DIS convolution, the double convolution $`qCD`$ is not commutative.
We shall frequently encounter expressions of the form
$`qD+{\displaystyle \frac{\alpha _s}{2\pi }}qCD`$ (53)
corresponding to the LO plus NLO corrections. The flavour structure of the results becomes much more transparent if we adopt the following symbolic notation:
$`qD+{\displaystyle \frac{\alpha _s}{2\pi }}qCD=q[1+{\displaystyle \frac{\alpha _s}{2\pi }}C]D.`$ (54)
Then the semi-inclusive polarized cross section $`\mathrm{\Delta }\stackrel{~}{\sigma }_p^h`$ defined in (7) is given by
$`\mathrm{\Delta }\stackrel{~}{\sigma }_p^h`$ $`=`$ $`{\displaystyle \underset{i}{}}e_i^2\mathrm{\Delta }q_i[1+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qq}]D_{q_i}^h+`$ (55)
$`+\left({\displaystyle \underset{i}{}}e_i^2\mathrm{\Delta }q_i\right){\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qg}D_G^h+\mathrm{\Delta }G{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{gq}\left({\displaystyle \underset{i}{}}e_i^2D_{q_i}^h\right)`$
where the sum is over quarks and antiquarks of flavour $`i`$ and parton densities and fragmentation functions are to be taken in NLO.
For the unpolarized semi-inclusive cross section in NLO it is not possible to completely factor out the $`y`$-dependence. Consequently $`\stackrel{~}{\sigma }^h`$ defined in (6) will depend upon $`y`$ in NLO, in contrast to the LO situation in (9).
In the notation of the seminal paper of Graudenz the cross section is a sum of what he refers to as “metric” ($`M`$) and “longitudinal” ($`L`$) terms, with corresponding Wilson coefficients $`C^M`$ and $`C^L`$. A further complication is that the Wilson coefficients are different in the two regions of integration $`I_1`$ and $`I_2`$. Thus we have coefficients $`C^{jM}`$, $`C^{jL}`$ with $`j=1,2`$.
We then define the $`y`$-dependent combinations of Wilson coefficients:
$`_{qq}^j=C_{qq}^{jM}+\left[1+4\gamma (y)\right]C_{qq}^{jL}`$ (56)
$`_{qg}^j=C_{qg}^{jM}+\left[1+4\gamma (y)\right]C_{qg}^{jL}`$ (57)
$`_{gq}^j=C_{gq}^{jM}+\left[1+4\gamma (y)\right]C_{gq}^{jL}`$ (58)
where
$`\gamma (y)={\displaystyle \frac{1y}{1+(1y)^2}}.`$ (59)
Then the unpolarized semi-inclusive cross section can be written in a form analogous to the polarized one:
$`\stackrel{~}{\sigma }_p^h`$ $`=`$ $`{\displaystyle \underset{i}{}}e_i^2q_i[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_{q_i}^h+`$ (60)
$`+\left({\displaystyle \underset{i}{}}e_i^2q_i\right){\displaystyle \frac{\alpha _s}{2\pi }}_{qg}D_G^h+G{\displaystyle \frac{\alpha _s}{2\pi }}_{gq}\left({\displaystyle \underset{i}{}}e_i^2D_{q_i}^h\right).`$
In (60) we have used the symbolic notation:
$`q_i{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}D_{q_i}^h={\displaystyle _{I_1}}q_i{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}^1D_{q_i}^h+{\displaystyle _{I_2}}q_i{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}^2D_{q_i}^h,`$ (61)
and analogously for $`_{qg}`$ and $`_{gq}`$.
Note that in NLO the unpolarized inclusive cross section, $`\stackrel{~}{\sigma }^{DIS}`$ (12) is given by
$`\stackrel{~}{\sigma }^{DIS}=2F_1\left[1+2\gamma (y)R\right],`$ (62)
where $`R`$ is the usual DIS ratio of longitudinal to transverse cross sections. Note that strictly speaking here and throughout the rest of this paper $`R`$ should be replaced by $`R\left(R\frac{4M^2x^2}{Q^2}\right)/\left(1+\frac{4M^2x^2}{Q^2}\right)`$, since the correction terms may be important for low values of $`Q^2`$.
As in the LO discussion we doubt the reliability of a global NLO analysis of inclusive and semi-inclusive data, and we suggest that one should feed into the semi-inclusive analysis as much reliable information as one can from other sources.
As we shall see there is the oft found opposition between what is simple theoretically and what is simple experimentally. However, if the systematic errors in detection efficiencies can be brought under control then we can make remarkable theoretical simplifications and we can then extract a vast amount of information from the semi-inclusive data. This is an important experimental challenge as will be seen from the power of the results given below.
### 5.1 Simplification of the semi-inclusive NLO results
Bearing in mind the NLO result (21) for $`g_1^p`$ for polarized DIS and the analogous result for $`F_1`$ in unpolarized DIS, we see that in the second term of (60) we may make the replacement, correct to NLO,
$`{\displaystyle \underset{i}{}}e_i^2q_i2F_1^p`$ (63)
and similarly, in the analogous eq. for $`\mathrm{\Delta }\stackrel{~}{\sigma }_p`$ in (55):
$`{\displaystyle \underset{i}{}}e_i^2\mathrm{\Delta }q_i2g_1^p.`$ (64)
Note that here and throughout the rest of this paper $`F_1^{p,n}`$ and $`g_1^{p,n}`$ are the experimentally measured structure functions.
Further in the reaction $`e^+(p_+)+e^{}(p_{})h+X`$ in the kinematic region when we can neglect $`Z^0`$-exchange effects, we have
$`\sigma ^h(z,\mathrm{cos}\theta ,Q^2)={\displaystyle \frac{3}{8}}(1+\mathrm{cos}^2\theta )\sigma _T^h(z,Q^2)+{\displaystyle \frac{3}{4}}(1\mathrm{cos}^2\theta )\sigma _L^h(z,Q^2)`$ (65)
where $`Q^2=(p_++p_{})^2`$.
In NLO QCD one has
$`\sigma _T^h(z,Q^2)`$ $`=`$ $`3\sigma _0\{{\displaystyle \underset{i}{}}e_i^2D_{q_i}^h(1+{\displaystyle \frac{\alpha _s}{2\pi }}C_q^T)+`$ (66)
$`+{\displaystyle \underset{i}{}}e_i^2D_G^h{\displaystyle \frac{\alpha _s}{2\pi }}C_G^T\}`$
where the sum is over quarks and antiquarks, the $`C^T`$’s are Wilson coefficients, and
$`\sigma _0={\displaystyle \frac{4\pi \alpha ^2}{3Q^2}}.`$ (67)
Then, correct to the required NLO accuracy, in the third term in (60), and in its polarized analogue (55), we may make the replacement
$`{\displaystyle \underset{i}{}}e_i^2D_{q_i}^h={\displaystyle \frac{\sigma _T^h(z,Q^2)}{3\sigma _0}},`$ (68)
where $`\sigma _T^h`$ is the experimentally measured cross section.
Hence, for the unpolarized and polarized semi-inclusive cross sections, in NLO accuracy, we have
$`\stackrel{~}{\sigma }_p^h`$ $`=`$ $`{\displaystyle \underset{i}{}}e_i^2q_i[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_{q_i}^h+`$ (69)
$`+2F_1^p{\displaystyle \frac{\alpha _s}{2\pi }}_{qg}D_G^h+{\displaystyle \frac{1}{3\sigma _0}}G{\displaystyle \frac{\alpha _s}{2\pi }}_{gq}\sigma _T^h`$
$`\mathrm{\Delta }\stackrel{~}{\sigma }_p^h`$ $`=`$ $`{\displaystyle \underset{i}{}}e_i^2\mathrm{\Delta }q_i[1+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qq}]D_{q_i}^h+`$ (70)
$`+2g_1^p{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qg}D_G^h+{\displaystyle \frac{1}{3\sigma _0}}\mathrm{\Delta }G{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{gq}\sigma _T^h.`$
Given that the unpolarized gluon density is reasonably well known, the last term in (69) can be considered as a known quantity. In the following we take as known quantities the NLO values for $`q_+(x,Q^2)`$, $`q_3(x,Q^2)`$, $`u_V(x,Q^2)`$, $`d_V(x,Q^2)`$, $`G(x,Q^2)`$ and $`\mathrm{\Delta }q_3(x,Q^2)`$.
### 5.2 The polarized valence densities in NLO
Using charge conjugation invariance one obtains, for semi-inclusive pion production
$`R_p^{\pi ^+\pi ^{}}={\displaystyle \frac{[4u_Vd_V][1+(\alpha _s/2\pi )_{qq}]D_u^{\pi ^+\pi ^{}}}{18F_1^p[1+2\gamma (y)R^p]}}`$
$`R_n^{\pi ^+\pi ^{}}={\displaystyle \frac{[4d_Vu_V][1+(\alpha _s/2\pi )_{qq}]D_u^{\pi ^+\pi ^{}}}{18F_1^p[1+2\gamma (y)R^p]}}.`$ (71)
The only unknown function in these expressions is $`D_u^{\pi ^+\pi ^{}}(z,Q^2)`$, which evolves as a non-singlet and does not mix with other FF’s. A $`\chi ^2`$ analysis of either or both of the equations (71) should thus determine $`D_u^{\pi ^+\pi ^{}}`$ in NLO without serious ambiguity.
Assuming $`D_u^{\pi ^+\pi ^{}}`$ is now known, one can then determine $`\mathrm{\Delta }u_V`$ and $`\mathrm{\Delta }d_V`$ in NLO via the equations
$`\mathrm{\Delta }A_p^{\pi ^+\pi ^{}}={\displaystyle \frac{(4\mathrm{\Delta }u_V\mathrm{\Delta }d_V)[1+(\alpha _s/2\pi )\mathrm{\Delta }C_{qq}]D_u^{\pi ^+\pi ^{}}}{(4u_Vd_V)[1+(\alpha _s/2\pi )_{qq}]D_u^{\pi ^+\pi ^{}}}}`$ (72)
$`\mathrm{\Delta }A_n^{\pi ^+\pi ^{}}={\displaystyle \frac{(4\mathrm{\Delta }d_V\mathrm{\Delta }u_V)[1+(\alpha _s/2\pi )\mathrm{\Delta }C_{qq}]D_u^{\pi ^+\pi ^{}}}{(4d_Vu_V)[1+(\alpha _s/2\pi )_{qq}]D_u^{\pi ^+\pi ^{}}}}`$ (73)
where, of course, $`\mathrm{\Delta }u_V`$ and $`\mathrm{\Delta }d_V`$ evolve as non-singlets and do not mix with other densities. Eqs. (72) and (73) determine the densities $`\mathrm{\Delta }u_V`$ and $`\mathrm{\Delta }d_V`$ in NLO without any assumptions about the less known polarized gluon and sea densities.
### 5.3 SU(2) symmetry of the sea quark densities in NLO
Once $`\mathrm{\Delta }u_V`$ and $`\mathrm{\Delta }d_V`$ are known in NLO we can calculate
$`[\mathrm{\Delta }\overline{u}(x,Q^2)\mathrm{\Delta }\overline{d}(x,Q^2)]_{NLO}={\displaystyle \frac{1}{2}}[\mathrm{\Delta }q_3(x,Q^2)+\mathrm{\Delta }d_V(x,Q^2)\mathrm{\Delta }u_V(x,Q^2)]_{NLO},`$ (74)
Eq. (74) determines the breaking of SU(2) symmetry for the polarized sea densities in NLO without requiring any knowledge of $`\mathrm{\Delta }\overline{q}`$ and $`\mathrm{\Delta }G`$. It will be interesting to compare the values obtained from (74) with information on $`\mathrm{\Delta }\overline{u}(x)`$ and $`\mathrm{\Delta }\overline{d}(x)`$ which will emerge from Drell-Yan and $`W^\pm `$ production experiments at RHIC .
The separate determination of $`\mathrm{\Delta }\overline{u}`$ and $`\mathrm{\Delta }\overline{d}`$ requires knowledge of $`\mathrm{\Delta }q_+`$, defined in (26), in NLO.
Note that determining, say, $`\mathrm{\Delta }\overline{u}`$ via $`\mathrm{\Delta }\overline{u}=\frac{1}{2}[(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u})\mathrm{\Delta }u_V]`$ is unreliable, since the determination of $`(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u})`$ from inclusive DIS involves knowledge of gluon and strange quark densities.
As in the LO case we can only determine $`\mathrm{\Delta }q_+`$ after obtaining some information about the fragmentation functions.
### 5.4 Fragmentation functions in NLO
We consider the sum for the unpolarized production of $`\pi ^+`$ and $`\pi ^{}`$. We have
$`R_{pn}^{\pi ^++\pi ^{}}={\displaystyle \frac{q_3\{[1+\frac{\alpha _s}{2\pi }_{qq}]D_u^{\pi ^++\pi ^{}}+\frac{\alpha _s}{2\pi }_{qg}D_G^{\pi ^++\pi ^{}}\}}{6\left[F_1^p\left(1+2\gamma (y)R^p\right)F_1^n\left(1+2\gamma (y)R^n\right)\right]}}`$ (75)
$`\mathrm{\Delta }A_{pn}^{\pi ^++\pi ^{}}={\displaystyle \frac{\mathrm{\Delta }q_3\{[1+\frac{\alpha _s}{2\pi }\mathrm{\Delta }C_{qq}]D_u^{\pi ^++\pi ^{}}+\frac{\alpha _s}{2\pi }\mathrm{\Delta }C_{qg}D_G^{\pi ^++\pi ^{}}\}}{q_3\{[1+\frac{\alpha _s}{2\pi }_{qq}]D_u^{\pi ^++\pi ^{}}+\frac{\alpha _s}{2\pi }_{qg}D_G^{\pi ^++\pi ^{}}\}}}`$ (76)
The only unknown functions in these relations are $`D_u^{\pi ^++\pi ^{}}`$ and $`D_G^{\pi ^++\pi ^{}}`$, which will mix with each other under evolution. Thus it should be possible to obtain them from a $`\chi ^2`$ analysis of (75) and (76).
Once we know $`D_u^{\pi ^++\pi ^{}}`$ and utilise $`D_u^{\pi ^+\pi ^{}}`$ from Section 6.2 we clearly have access to both $`D_u^{\pi ^+}`$ and $`D_u^\pi ^{}`$.
Note that if $`K^0`$ can be detected one can obtain information on $`D_u^K`$ and $`D_G^K`$, where $`K=K^++K^{}+K^0+\overline{K}^0`$. One simply replaces the labels $`\pi ^++\pi ^{}`$ by $`K`$ everywhere in (75) and (76). Analogous equations hold also if $`\pi ^++\pi ^{}`$ is replaced by $`\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}`$, if $`\mathrm{\Lambda }`$ is detected.
Returning to the case of $`\pi ^++\pi ^{}`$, the ratio $`R_{p+n}^{\pi ^++\pi ^{}}`$ allows the determination of the only unknown function $`D_s^{\pi ^++\pi ^{}}`$. We have
$`R_{p+n}^{\pi ^++\pi ^{}}=`$
$`=\{(5q_+[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_u^{\pi ^++\pi ^{}}+2(s+\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_s^{\pi ^++\pi ^{}})+`$
$`+18(F_1^p+F_1^n){\displaystyle \frac{\alpha _s}{2\pi }}_{gq}D_G^{\pi ^++\pi ^{}}+{\displaystyle \frac{6}{\sigma _0}}G{\displaystyle \frac{\alpha _s}{2\pi }}_{gq}\sigma _T^{\pi ^++\pi ^{}}\}/`$
$`/6\left[F_1^p\left(1+2\gamma (y)R^p\right)+F_1^n\left(1+2\gamma (y)R^n\right)\right]`$ (77)
Under evolution $`D_u^{\pi ^++\pi ^{}}`$ and $`D_s^{\pi ^++\pi ^{}}`$ mix with $`D_G^{\pi ^++\pi ^{}}`$, but this is not a problem since the latter is supposed to be known. It would be interesting to compare these results with those from $`e^+e^{}hadrons`$, obtained recently in .
Again analogous sets of equations holds for kaon and $`\mathrm{\Lambda },\overline{\mathrm{\Lambda }}`$ production. One simply replaces $`\pi ^++\pi ^{}`$ by $`K`$ or $`\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}`$ and this allows the determination of the only unknown function $`D_s^K`$ or $`D_s^{\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}}`$, respectively.
### 5.5 The sea-quark densities in NLO
With the NLO knowledge of $`D_u^{\pi ^++\pi ^{}}`$, $`D_s^{\pi ^++\pi ^{}}`$ and $`D_G^{\pi ^++\pi ^{}}`$ we are now in a position to try to determine $`\mathrm{\Delta }q_+(x,Q^2)`$, $`\mathrm{\Delta }s(x,Q^2)+\mathrm{\Delta }\overline{s}(x,Q^2)`$ and $`\mathrm{\Delta }G(x,Q^2)`$. We have, in NLO,
$`g_1^p+g_1^n={\displaystyle \frac{1}{18}}\left[5\mathrm{\Delta }q_++2(\mathrm{\Delta }s+\mathrm{\Delta }\overline{s})\right]\left(1+{\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}\delta C_q\right)+{\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}\mathrm{\Delta }G\delta C_G,`$ (78)
and
$`\mathrm{\Delta }A_{p+n}^{\pi ^++\pi ^{}}={\displaystyle \frac{\mathrm{\Delta }\stackrel{~}{\sigma }_p^{\pi ^++\pi ^{}}+\mathrm{\Delta }\stackrel{~}{\sigma }_n^{\pi ^++\pi ^{}}}{\stackrel{~}{\sigma }_p^{\pi ^++\pi ^{}}+\stackrel{~}{\sigma }_n^{\pi ^++\pi ^{}}}}`$ (79)
where
$`\mathrm{\Delta }\stackrel{~}{\sigma }_p^{\pi ^++\pi ^{}}+\mathrm{\Delta }\stackrel{~}{\sigma }_n^{\pi ^++\pi ^{}}=`$
$`=(5\mathrm{\Delta }q_+[1+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qq}]D_u^{\pi ^++\pi ^{}}+2(\mathrm{\Delta }s+\mathrm{\Delta }\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qq}]D_s^{\pi ^++\pi ^{}})+`$
$`+18\left(g_1^p+g_1^n\right){\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}\mathrm{\Delta }C_{qg}D_G^{\pi ^++\pi ^{}}+{\displaystyle \frac{6}{\sigma _0}}\mathrm{\Delta }G{\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}\mathrm{\Delta }C_{gq}\sigma _T^{\pi ^++\pi ^{}}`$ (80)
and
$`\stackrel{~}{\sigma }_p^{\pi ^++\pi ^{}}+\stackrel{~}{\sigma }_n^{\pi ^++\pi ^{}}=`$
$`=(5q_+[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_u^{\pi ^++\pi ^{}}+2(s+\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_s^{\pi ^++\pi ^{}})+`$
$`+18\left(F_1^p+F_1^n\right){\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}_{qg}D_G^{\pi ^++\pi ^{}}+{\displaystyle \frac{6}{\sigma _0}}G{\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}_{gq}\sigma _T^{\pi ^++\pi ^{}}.`$ (81)
Note that an analogous set of equations hold for $`\pi ^++\pi ^{}`$ replaced by $`K`$ or $`\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}`$.
Eqs. (78) to (81) contain three unknown functions $`\mathrm{\Delta }q_+`$, $`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}`$ and $`\mathrm{\Delta }G`$, which nonetheless can all be determined in principle because of their different evolution in $`Q^2`$. However, to be at all efficacious such a determination would require a huge range of $`Q^2`$, far larger than is available in present day polarised DIS.
On the other hand there is a direct and superior method for obtaining $`\mathrm{\Delta }G`$, namely via $`c\overline{c}`$ production. This is one of the major goals of the COMPASS experiment at CERN. We shall thus assume that $`\mathrm{\Delta }G`$ has been determined, so that the last term on the RHS of (70) may be taken to be known.
It should be then straightforward to determine $`\mathrm{\Delta }q_+`$ and $`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}`$ in NLO from a $`\chi ^2`$ analysis of (78) to (81) in which the evolution of $`\mathrm{\Delta }q_+`$ and $`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}`$ would involve mixing with the supposed known $`\mathrm{\Delta }G`$.
An independent determination of $`\mathrm{\Delta }q_+`$ and $`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}`$ could be obtained by combining (78) with $`\mathrm{\Delta }A_{p+n}^K`$ for kaons and $`\mathrm{\Delta }A_{p+n}^{\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}}`$ for $`\mathrm{\Lambda },\overline{\mathrm{\Lambda }}`$ production.
Once $`\mathrm{\Delta }q_+`$ is known in NLO, we can obtain the individual $`\mathrm{\Delta }\overline{u}`$ and $`\mathrm{\Delta }\overline{d}`$ from (74).
### 5.6 Concerning $`s=\overline{s}`$ and $`\mathrm{\Delta }s=\mathrm{\Delta }\overline{s}`$ in NLO
It is possible to get some information on $`s\overline{s}`$ and $`\mathrm{\Delta }s\mathrm{\Delta }\overline{s}`$ in NLO.
We have, assuming $`D_d^{K^+}=D_d^K^{}`$,
$`R_p^{K^+K^{}}=`$
$`=\{4u_V[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_u^{K^+K^{}}+(s\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_s^{K^+K^{}}\}/`$
$`/18F_1^p\left[1+2\gamma (y)R^p\right]`$ (82)
$`R_n^{K^+K^{}}=`$
$`=\{4d_V[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_u^{K^+K^{}}+(s\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_s^{K^+K^{}}\}/`$
$`/18F_1^n\left[1+2\gamma (y)R^n\right]`$ (83)
These two equations, taken together with those for $`\mathrm{\Delta }A_p^{K^+K^{}}`$ and $`\mathrm{\Delta }A_n^{K^+K^{}}`$:
$`\mathrm{\Delta }A_p^{K^+K^{}}=`$
$`=(4\mathrm{\Delta }u_V[1+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qq}]D_u^{K^+K^{}}+(\mathrm{\Delta }s\mathrm{\Delta }\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qq}]D_s^{K^+K^{}})/`$
$`/(4u_V[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_u^{K^+K^{}}+(s\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_s^{K^+K^{}})`$ (84)
$`\mathrm{\Delta }A_n^{K^+K^{}}=`$
$`=(4\mathrm{\Delta }d_V[1+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qq}]D_u^{K^+K^{}}+(\mathrm{\Delta }s\mathrm{\Delta }\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qq}]D_s^{K^+K^{}})/`$
$`/(4d_V[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_u^{K^+K^{}}+(s\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_s^{K^+K^{}})`$ (85)
provide four equations for the three unknown functions $`(s\overline{s})D_s^{K^+K^{}}`$, $`(\mathrm{\Delta }s\mathrm{\Delta }\overline{s})D_s^{K^+K^{}}`$ and $`D_u^{K^+K^{}}`$, so that, in principle, all can be determined via a $`\chi ^2`$ analysis.
In addition one has
$`R_p^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}=`$
$`=\{(4u_V+d_V)[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}+(s\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}\}/`$
$`/18F_1^p\left[1+2\gamma (y)R^p\right]`$ (86)
and
$`R_n^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}=`$
$`=\{(4d_V+u_V)[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}+(s\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}\}/`$
$`/18F_1^n\left[1+2\gamma (y)R^n\right]`$ (87)
together with $`\mathrm{\Delta }A_p^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$ and $`\mathrm{\Delta }A_n^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$:
$`\mathrm{\Delta }A_p^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}=`$
$`=((4\mathrm{\Delta }u_V+\mathrm{\Delta }d_V)[1+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qq}]D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}+(\mathrm{\Delta }s\mathrm{\Delta }\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qq}]D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}})/`$
$`/((4u_V+d_V)[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}+(s\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}})`$ (88)
$`\mathrm{\Delta }A_n^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}=`$ (89)
$`=((4\mathrm{\Delta }d_V+\mathrm{\Delta }d_V)[1+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qq}]D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}+(\mathrm{\Delta }s\mathrm{\Delta }\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }C_{qq}]D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}})/`$
$`/((4d_V+d_V)[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}+(s\overline{s})[1+{\displaystyle \frac{\alpha _s}{2\pi }}_{qq}]D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}})`$ .
These provide four more equations but only two new unknown functions $`D_u^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$ and $`D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$. The system of eight equations (82) - (89) therefore over-constrains the unknown functions and might, hopefully, allow a reasonable determination of the relation between $`s\overline{s}`$ and $`\mathrm{\Delta }s\mathrm{\Delta }\overline{s}`$ and whether or not $`s(x)=\overline{s}(x)`$ and/or $`\mathrm{\Delta }s(x)=\mathrm{\Delta }\overline{s}(x)`$. The actual determination of $`s\overline{s}`$ or $`\mathrm{\Delta }s\mathrm{\Delta }\overline{s}`$ requires knowledge of either $`D_s^{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$ or $`D_s^{K^+K^{}}`$. These could be taken from the study of $`e^+e^{}hX`$.
This completes the determination of all the polarized densities in NLO.
## 6 Conclusions
We have argued that the present LO QCD method of analysing polarized semi-inclusive DIS, using the concept of purity, is quite unjustified. We have also argued against attempts at a global analysis, either in LO or in NLO QCD, based on the combined data on polarized inclusive and semi-inclusive DIS and taking as known exactly the relevant FF’s.
Instead, we have presented a strategy for a step by step evaluation of the polarized parton densities and fragmentation functions from semi-inclusive data using selectively chosen information from inclusive DIS reactions.
In our approach the usually made simplifying assumptions about relations between $`\mathrm{\Delta }\overline{u}`$ and $`\mathrm{\Delta }\overline{d}`$, and between the strange and non-strange polarized sea densities are unnecessary and we have even considered the possibility that $`s(x)\overline{s}(x)`$ and $`\mathrm{\Delta }s(x)\mathrm{\Delta }\overline{s}(x)`$.
Given the simplicity of the LO QCD analysis, we discuss where and when it is likely to be reliable, and stress the need to test the consistency of the LO treatment at each step. In this connection we have introduced the concept of a passive variable in the experimentally measured observables. We have also suggested how one might estimate the errors induced in doing the LO analysis.
In the NLO treatment we have shown how the expressions for the experimental observables can be much simplified by incorporating information from the reaction $`e^+e^{}hX`$ in a novel way. Determination of the polarized valence quark densities $`\mathrm{\Delta }u_V`$ and $`\mathrm{\Delta }d_V`$ is shown to be relatively straight forward, as is the difference $`\mathrm{\Delta }\overline{u}\mathrm{\Delta }\overline{d}`$. However we argue that the determination of $`\mathrm{\Delta }\overline{u}`$, $`\mathrm{\Delta }\overline{d}`$, $`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}`$ and $`\mathrm{\Delta }G`$ separately, from semi-inclusive DIS involving production of $`\pi `$, $`K`$, $`\mathrm{\Lambda }`$ is unlikely to be successful, because of the limited range of $`Q^2`$ available now and in the foreseeable future. It is suggested that the independent determination of $`\mathrm{\Delta }G`$ from charm production is thus an essential element if $`\mathrm{\Delta }\overline{u}`$, $`\mathrm{\Delta }\overline{d}`$, $`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}`$ are to be determined accurately. Finally, motivated by the arguments that possibly $`s(x)\overline{s}(x)`$ and $`\mathrm{\Delta }s(x)\mathrm{\Delta }\overline{s}(x)`$, we have demonstrated how, in principle, one can learn about $`s(x)\overline{s}(x)`$ and $`\mathrm{\Delta }s(x)\mathrm{\Delta }\overline{s}(x)`$.
The procedure we have advocated poses a real challenge to the experimentalists, since it requires a control over the systematic errors involved in hadron detection efficiencies. The price paid in the current practice of considering certain ratios of cross sections in order to limit systematic errors, is enormous, and vast amounts of interesting and theoretically valuable information are thereby lost. We hope this paper will encourage efforts to proceed further.
## 7 Acknowledgements
We thank M. Anselmino, H. Blok, M. Boglione, S. J. Brodsky, U. D’Alesio, D. de Florian, A. Kotzinjan, S. Kretzer, R. Sassot, X. Song, D. Stamenov and C. Weiss for helpful discussions. We are grateful to the UK Royal Society for a Collaborative Grant. E.C. is grateful to the Theory Division of CERN for its hospitality where this work was partly done and finished, her work was also supported by the Bulgarian National Science Foundation. E.L. is grateful to P.J. Mulders for the hospitality of the Department of Physics and Astronomy, the Vrije Universitiet, Amsterdam, where part of this work was carried out supported by the Foundation for Fundamental Research on Matter (FOM) and the Dutch Organisation for Scientific Research (NWO).
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# Matter effects in the 𝐷⁰-(𝐷⁰)̄ system
## I Introduction
The $`D^0\overline{D^0}`$ system is of great current interest. Until recently, the best probes on the mixing resulted only in upper bounds on $`x_D^2+y_D^2`$ of order $`10^2`$ . This changed with the recent results from CLEO and FOCUS . The current situation is beautifully summarized in reference : there is evidence for $`y_D10^2`$ and for a large strong phase, corresponding to large SU(3) breaking effects. The problem with the strong, final-state interaction phases is that they are not known. They are fit in the same experiments that look for $`y_D`$. It would be interesting if matter effects could be used to shed some light on this issue. The idea arises from the kaon system, where the phases due to matter effects are measured experimentally. Thus, we would be dealing with a known CP-even phase.
There is one important difference between the kaon and $`D`$ systems. For a beam of neutral kaons in vacuum, the depletion from the beam is controlled by two exponentials, $`\mathrm{exp}(\mathrm{\Gamma }_St)`$ and $`\mathrm{exp}(\mathrm{\Gamma }_Lt)`$, corresponding to the short-lived ($`K_S`$) and long-lived ($`K_L`$) components. Here $`|\mathrm{\Delta }\mathrm{\Gamma }|2\mathrm{\Gamma }`$ and the two exponential fallouts are clearly separated. As a result, we may wait for the $`K_S`$ component to decay away and use the matter effects on $`K_L`$ to regenerate $`K_S`$; a phenomenon known as ‘regeneration’. In the $`D`$ (and the $`B_d`$) system the situation is very different: $`|\mathrm{\Delta }\mathrm{\Gamma }|\mathrm{\Gamma }`$ and the leading behavior in vacuum is given by $`\mathrm{exp}(\mathrm{\Gamma }t)`$, with the hyperbolic sine and cosine of $`\mathrm{\Delta }\mathrm{\Gamma }t`$ acting as small perturbations. In this system, the depletion from the beam is controlled by $`\mathrm{\Gamma }`$. Because we cannot resolve the two exponentials, the classical ‘regeneration’ experiments cannot be carried out and a new analysis is required. This is the question we address here.
In section II we review the general features of propagation of neutral mesons systems in matter. In section III we use an extrapolation from known kaon results to argue that these effects might, in principle, be sought. Unfortunately, we will also show that, even if visible, the matter effects cannot be used to probe $`x_D`$ or $`y_D`$. This is shown to be a generic feature of the ‘small-time’ approximation in the time-evolution. We draw our conclusions in section IV, where we remark briefly on the $`B_d`$ and $`B_s`$ systems.
## II Propagation in matter
Let us consider a generic neutral meson system $`P^0\overline{P^0}`$. Assuming CPT-conservation in vacuum, this system is characterized by two complex eigenvalues, $`\mu _a=m_ai\mathrm{\Gamma }_a/2`$ and $`\mu _b=m_bi\mathrm{\Gamma }_b/2`$, and by a mixing parameter, $`q/p`$. It is convenient to define the average and the difference of the eigenvalues,
$`\mu ={\displaystyle \frac{\mu _a+\mu _b}{2}}`$ $`=`$ $`m{\displaystyle \frac{i}{2}}\mathrm{\Gamma },`$ (1)
$`\mathrm{\Delta }\mu =\mu _a\mu _b`$ $`=`$ $`\mathrm{\Delta }m{\displaystyle \frac{i}{2}}\mathrm{\Delta }\mathrm{\Gamma }.`$ (2)
These quantities describe the system in vacuum and its time evolution is well known .
The new effects resulting from the interaction with matter can be described by the elastic forward scattering amplitudes of $`P^0`$ and $`\overline{P^0}`$, which we denote by $`f`$ and $`\overline{f}`$, respectively. These enter the effective Hamiltonian through
$$\chi =\frac{2\pi N}{m}f,\text{and}\overline{\chi }=\frac{2\pi N}{m}\overline{f},$$
(3)
where $`N`$ is the density of scattering centers in the medium. As before, we define the average and the difference of these parameters as
$`\chi _{\mathrm{av}}`$ $`=`$ $`{\displaystyle \frac{\chi +\overline{\chi }}{2}},`$ (4)
$`\mathrm{\Delta }\chi `$ $`=`$ $`\chi \overline{\chi }={\displaystyle \frac{2\pi N}{m}}\mathrm{\Delta }f,`$ (5)
with $`\mathrm{\Delta }f=f\overline{f}`$.
In the presence of matter, the time evolution of meson states with a given initial flavor may be written as ,
$`e^{i(\mu +\chi _{\mathrm{av}})t}|P^0(t)`$ $`=`$ $`\left[\mathrm{cos}\left(\mathrm{\Omega }t/2\right)i{\displaystyle \frac{\mathrm{\Delta }\chi }{\mathrm{\Omega }}}\mathrm{sin}\left(\mathrm{\Omega }t/2\right)\right]|P^0+{\displaystyle \frac{q}{p}}\left[i{\displaystyle \frac{\mathrm{\Delta }\mu }{\mathrm{\Omega }}}\mathrm{sin}\left(\mathrm{\Omega }t/2\right)\right]|\overline{P^0},`$ (6)
$`e^{i(\mu +\chi _{\mathrm{av}})t}|\overline{P^0}(t)`$ $`=`$ $`{\displaystyle \frac{p}{q}}\left[i{\displaystyle \frac{\mathrm{\Delta }\mu }{\mathrm{\Omega }}}\mathrm{sin}\left(\mathrm{\Omega }t/2\right)\right]|P^0+\left[\mathrm{cos}\left(\mathrm{\Omega }t/2\right)+i{\displaystyle \frac{\mathrm{\Delta }\chi }{\mathrm{\Omega }}}\mathrm{sin}\left(\mathrm{\Omega }t/2\right)\right]|\overline{P^0},`$ (7)
where $`\mathrm{\Omega }=\sqrt{(\mathrm{\Delta }\mu )^2+(\mathrm{\Delta }\chi )^2}`$ . The evolution in vacuum is reproduced by setting $`\chi _{\mathrm{av}}`$ and $`\mathrm{\Delta }\chi `$ to zero.
As mentioned in the introduction, in the $`D`$ (and in the $`B_d`$) system the depletion from the beam in vacuum is controlled by $`\mathrm{\Gamma }`$. This one exponential is then modulated by trigonometric functions of $`\mathrm{\Delta }\mu t=\frac{\mathrm{\Delta }\mu }{\mathrm{\Gamma }}\tau `$, where $`\tau =\mathrm{\Gamma }t`$ is the time measured in lifetime units. Therefore, one introduces
$$\frac{\mathrm{\Delta }\mu }{\mathrm{\Gamma }}=xiy.$$
(8)
In most cases of interest the matter effects satisfy $`\mathrm{\Gamma }>|\text{Im}\chi _{\mathrm{av}}|`$ and $`|\mathrm{\Delta }\mu |>|\mathrm{\Delta }\chi |`$, and the depletion is still chiefly determined by $`\mathrm{\Gamma }`$. To describe the trigonometric functions in matter we need Eq. (8) and a new complex parameter. This may be choosen to be $`\mathrm{\Delta }\chi /\mathrm{\Gamma }`$, which determines the depletion (statistics available) by the time the effects of $`\mathrm{\Delta }\chi `$ become important. However, since we know the current experimental reach on $`\mathrm{\Delta }\mu `$, it is easiest to use it to access the experimental reach on $`\mathrm{\Delta }\chi `$, throughAlthough it is conventional to use this factor of two in the definition of $`r`$, it seems to confuse rather than enlighten the issue.
$$2r=\frac{\mathrm{\Delta }\chi }{\mathrm{\Delta }\mu }.$$
(9)
Our figure of merit will be $`|\mathrm{\Delta }\chi /\mathrm{\Delta }\mu |`$.
Henceforth we will refer to those effects due to $`\mathrm{\Delta }\mu `$ as ‘mixing effects’, and to those effects due to $`\mathrm{\Delta }\chi `$ as ‘matter effects’. The fundamental question is: when do the matter effects and the mixing effects couple to each other? The answer is easily found by looking at Eqs. (7). They can couple in two ways. The first and simplest way occurs when the production, propagation and decay occur all in the same medium and the decay is flavor tagging. In that case, the matter and mixing effects couple mainly through $`\mathrm{\Omega }`$. The second way occurs when both flavors can decay into the same final state, or when the system crosses several media. For example, for events occuring completely in the medium, with the decays going into CP-eigenstates. Or, when the production and the initial part of the evolution occur in matter, with the final part of the evolution and its decay occuring in vacuum. In these cases, besides $`\mathrm{\Omega }`$, there can also be effects proportional to $`\mathrm{\Delta }\mu \mathrm{\Delta }\chi `$, which arise from the interference between the coefficients of the sine terms multiplying $`|P^0`$ and of the sine terms multiplying $`|\overline{P^0}`$.
## III Matter effects in the $`D^0\overline{D^0}`$ system
### A Interplay between mixing and matter effects
We will now show that, even when visible, the matter effects cannot be used to probe $`x_D`$ or $`y_D`$. The main reason is that the complex trigonometric functions are controled by $`(xiy)\tau `$ and, with current technology, one can only follow $`\tau `$ up to 5 or 10, corresponding to $`5`$ to $`10`$ average lifetimes in the $`D`$ proper frame. This is mostly due to the statistics of perfectly reconstructed events. It is obvious that, given unlimmited event samples, one could follow the full decay curve including times of order $`1/y_D`$. We mention this because the ‘no-go theorem’ we are about to present hinges on the experimental constraint that one can only follow the decay curves for time $`\tau 1/y_D`$. Under these circumstances, we can expand the trigonometric functions on the right-hand side of Eqs. (7) as
$$\mathrm{cos}(\mathrm{\Omega }t/2)\pm i\frac{\mathrm{\Delta }\chi }{\mathrm{\Omega }}\mathrm{sin}(\mathrm{\Omega }t/2)1\pm i\frac{\mathrm{\Delta }\chi }{2}t\frac{(\mathrm{\Delta }\mu )^2+(\mathrm{\Delta }\chi )^2}{8}t^2i\mathrm{\Delta }\chi \frac{(\mathrm{\Delta }\mu )^2+(\mathrm{\Delta }\chi )^2}{48}t^3+(t^4),$$
(10)
and
$$i\frac{\mathrm{\Delta }\mu }{\mathrm{\Omega }}\mathrm{sin}(\mathrm{\Omega }t/2)i\frac{\mathrm{\Delta }\mu }{2}t+i\mathrm{\Delta }\mu \frac{(\mathrm{\Delta }\mu )^2+(\mathrm{\Delta }\chi )^2}{48}t^3+(t^5),$$
(11)
for the right-sign and wrong-sign transitions, respectively.
The old results on semileptonic wrong-sign decays measure the magnitude squared of Eq. (11), given to leading order by
$$\frac{|\mathrm{\Delta }\mu |^2}{4\mathrm{\Gamma }^2}\tau ^2=\frac{x^2+y^2}{4}\tau ^2.$$
(12)
Let us now consider a final state $`f`$ to which both $`P^0`$ and $`\overline{P^0}`$ can decay. We denote the decay amplitudes by $`A_f`$ and $`\overline{A}_f`$, respectively. Such a decay has a term linear in $`t`$ arising from the interference between Eqs. (10) and (11),
$$\text{Im}\left[\mathrm{\Delta }\mu \frac{q}{p}\overline{A}_fA_f^{}\right]t.$$
(13)
CLEO has looked at $`f=K^+\pi ^{}`$ which is a Cabibbo-allowed decay for $`\overline{D^0}`$ and a Doubly-Cabibbo-suppressed decay for $`D^0`$. FOCUS has looked at the CP-eigenstate $`f=K^+K^{}`$.
Could one look at matter effects with semileptonic decays? In principle yes. The time evolution of the right-sign semileptonic decays is proportional to the magnitude squared of Eq. (10), which has a term linear in $`t`$,
$$\text{Im}\left(\frac{\mathrm{\Delta }\chi }{\mathrm{\Gamma }}\right)\tau .$$
(14)
Even if $`|\mathrm{\Delta }\chi |`$ is an order of magnitude smaller than $`|\mathrm{\Delta }\mu |`$, this term should still be easier to detect than the quadratic term in Eq. (12). This idea will be explained in detail below.
For the moment we wish to return to our fundamental question: can matter effects be used to give a new handle on $`\mathrm{\Delta }\mu `$? The answer is no! First consider right-sign semileptonic decays and magnitude square the right-hand side of Eq. (10). We see that only at order $`\tau ^3`$ do we find effects proportional to $`\mathrm{\Delta }\mu \mathrm{\Delta }\chi `$. Even at order $`\tau ^2`$, the effects decouple as $`\text{Re}(\mathrm{\Delta }\mu )^2`$ and $`\text{Re}(\mathrm{\Delta }\chi )^2`$. Next consider wrong-sign semileptonic decays and magnitude square the right-hand side of Eq. (11). This case is even worse, since the effects proportional to $`\mathrm{\Delta }\mu \mathrm{\Delta }\chi `$ only arise at order $`\tau ^4`$. Now consider again a final state $`f`$ to which both $`D^0`$ and $`\overline{D^0}`$ can decay. In the presence of matter Eq. (13) gets changed into
$$\text{Im}\left[\mathrm{\Delta }\mu \frac{q}{p}\overline{A}_fA_f^{}\left(t+i\frac{\mathrm{\Delta }\chi ^{}}{2}t^2\right)\right].$$
(15)
In this case there is an effect proportional to $`\mathrm{\Delta }\mu \mathrm{\Delta }\chi ^{}`$ at order $`t^2`$. But, if $`|\mathrm{\Delta }\chi |`$ is an order of magnitude smaller than $`|\mathrm{\Delta }\mu |`$, then this effect will be an order of magnitude smaller than that in Eq. (12). Notice that, in particular, there is no effect proportional to $`\mathrm{\Omega }t`$.
In this discussion we have implicitly considered only events which occur totally in matter. However, the matter-vacuum transition can only promote the same type of interference effects we have already discussed .
### B Observing matter effects
We have already argued that matter effects cannot be used as a handle on $`x_D`$ and $`y_D`$. We have yet to prove that they are observable at all, even in principle. Let us consider the decay $`D^0K^{}l^+\nu _l`$. Using Eqs. (7) and (10) we find
$$\mathrm{\Gamma }\left[D^0(t)K^{}l^+\nu _l\right]e^{\mathrm{\Gamma }t}e^{2\text{Im}\chi _{\mathrm{av}}t}|A_o|^2\left(1+\text{Im}\mathrm{\Delta }\chi t\right),$$
(16)
where we have denoted $`A_o=A(D^0K^{}l^+\nu _l)`$.To illustrate this effect, we have simplified the discussion by considering a ‘thought-experiment’ with the events taking place (production, evolution, and decay) inside the material. For these purposes, a muon in the final state would be most readily observed. We now need to estimate the linear term.
However, since the elastic forward scattering amplitudes for $`D^0`$ and $`\overline{D}^0`$ ($`f`$ and $`\overline{f}`$) are not known, we have to estimate them somehow. We recall that the matter effects in the kaon system arise because $`K^0`$ interacts quasi-elastically with the nucleons, while $`\overline{K^0}`$ suffers inelastic interactions, such as
$`\overline{K^0}+(p^+,n)`$ $``$ $`\pi ^0+(\mathrm{\Sigma }^+,\mathrm{\Sigma }^0)`$ (17)
$`s\overline{d}+(uud,udd)`$ $``$ $`d\overline{d}+(uus,uds).`$ (18)
The counterpart in the $`D`$ system is
$`D^0+(p^+,n)`$ $``$ $`\pi ^0+(\mathrm{\Sigma }_c^+,\mathrm{\Sigma }_c^0)`$ (19)
$`c\overline{u}+(uud,udd)`$ $``$ $`u\overline{u}+(udc,ddc).`$ (20)
The kinematics are different because $`m(D^0)1865\text{MeV}m(K^0)498\text{MeV}`$ and $`m(\mathrm{\Sigma }_c^+)2454\text{MeV}m(\mathrm{\Sigma }^+)1189\text{MeV}`$. However, on the one hand we will be interested in the high momentum range (larger than 25GeV), where the difference becomes less important. And, on the other hand, we are not interested in exact results but, rather, wish to learn whether the matter effects are in principle observable. Therefore, we will simply scale the results from the kaon system to the $`D`$ system.
We will use an empirical scaling law determined by Gsponer and collaborators from their measurements in C, Al, Cu, Sn, and Pb for kaon momenta between 30 to 150 GeV. They found that
$$|\mathrm{\Delta }f|=1.13\text{fm}\left(\frac{A}{\text{g mol}^1}\right)^{0.758}\left(\frac{p_K}{\text{GeV}c^1}\right)^{0.386},$$
(21)
where $`A`$ the atomic number of the material. This result exhibits a power law momentum dependence in accordance with Regge theory , which also predicts that $`\mathrm{arg}\mathrm{\Delta }f`$ should be constant and given by $`(1+0.386)\pi /2=0.693\pi `$ . The power-law approximation is fairly good down to a few GeV/$`c`$ momentum, where low-energy resonances set in . We neglect these effects since we will be concentrating on the high momentum range. We will also need the imaginary part of $`f+\overline{f}`$. Again we rely on Gsponer and collaborators, who found
$$\text{Im}(f+\overline{f})=1.895\text{fm}\left(\frac{A}{\text{g mol}^1}\right)^{0.840}\left(\frac{p_K}{\text{GeV}c^1}\right)$$
(22)
in the kaon sector.
We start by noticing that the density of scattering centers is a medium is given by
$$N=\frac{N_A\rho }{A},$$
(23)
where $`N_A`$ is Avogadro’s number, $`\rho `$ is the density of the material, and $`A`$ is its atomic number. Using Eqs. (5) and (21), we obtain
$`{\displaystyle \frac{\mathrm{\Delta }\chi _D}{\mathrm{\Gamma }_D}}`$ $`=`$ $`{\displaystyle \frac{1}{m_D\mathrm{\Gamma }_D}}\left({\displaystyle \frac{2\pi N_A\rho }{A}}\right)\mathrm{\Delta }f`$ (24)
$`=`$ $`5.6\times 10^5e^{i0.307\pi }\left({\displaystyle \frac{A}{\text{g mol}^1}}\right)^{0.242}\left({\displaystyle \frac{\rho }{\text{g cm}^3}}\right)^{0.242}\left({\displaystyle \frac{p_D}{\text{GeV}c^1}}\right)^{0.386}.`$ (25)
This result should be compared with that obtained for the neutral kaons, where $`5.6\times 10^5`$ is substituted by $`0.09`$. The sign difference arises from the fact that the inelastic interactions occur for $`D^0`$, not $`\overline{D^0}`$; for example, Eq. (20) involves $`D^0`$, while Eq. (18) involves $`\overline{K^0}`$. Since we are assuming that the forward scattering amplitudes coincide, the difference in magnitudes arises from
$$\frac{|\mathrm{\Delta }\chi _D|/\mathrm{\Gamma }_D}{|\mathrm{\Delta }\chi _K|/\mathrm{\Gamma }_K}=\frac{m_K\mathrm{\Gamma }_K}{m_D\mathrm{\Gamma }_D}=\frac{498\text{MeV}}{1865\text{MeV}}\frac{0.415\times 10^{12}s}{2\times 0.893\times 10^{10}s}=6.2\times 10^4.$$
(26)
The results for carbon and tungsten are listed in table I for a variety of $`D`$ momenta chosen within a range accessible to FOCUS . Tungsten is a preferable material since it leads to values which are roughly $`4.4`$ times larger than those obtained in carbon. For tungsten and momenta above $`100\text{GeV}c^1`$ the effects are of order 0.2%. However, the current experiments already probe $`\mathrm{\Delta }\mu _D/\mathrm{\Gamma }_D`$ of order 1%. Therefore, in principle, the linear term in Eq. (16) would be detected with an experiment of this type.
Of course, we still need to show that the absorption term $`\mathrm{exp}(2\text{Im}\chi _{\mathrm{av}}t)`$ does not affect greatly the number of events. Using Eq. (22), we find
$$2\frac{\text{Im}\chi _{\mathrm{av}D}}{\mathrm{\Gamma }_D}=9.4\times 10^5\left(\frac{A}{\text{g mol}^1}\right)^{0.160}\left(\frac{\rho }{\text{g cm}^3}\right)\left(\frac{p_D}{\text{GeV}c^1}\right).$$
(27)
For tungsten and a momentum of $`100\text{GeV}c^1`$, the ratio becomes $`0.08`$. Looking back at the right-hand side of Eq. (16), this result means that the absorption exponential, $`\mathrm{exp}(2\text{Im}\chi _{\mathrm{av}}t)`$, does not compete with the decay exponential, $`\mathrm{exp}(\mathrm{\Gamma }t)`$.
Although we have been detailed about our estimates, our objective is not to propose a specific experiment, but rather to emphasize that one is not, in principle, out of reach. Therefore, we will not address the obvious list of questions: i) How does one identify the primary and secondary vertices?; ii) How does one instrument the material?; iii) What is the ideal momentum range?; etc. We stress that the results presented here are based on the expectation that we can extrapolate from the kaon to the $`D`$ system. They should be taken as estimates and not as reliable predictions. In particular, the presence of resonances could enhance or suppress these estimates.
## IV Conclusions
The situation in the $`B_d`$ system is similar to the one in the $`D`$ system in that $`\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Gamma }`$. Using the same naive high-momenta, scaling law used in Eq. (26), we obtain
$$\frac{|\mathrm{\Delta }\chi _B|/\mathrm{\Gamma }_B}{|\mathrm{\Delta }\chi _K|/\mathrm{\Gamma }_K}=\frac{m_K\mathrm{\Gamma }_K}{m_B\mathrm{\Gamma }_B}=8.2\times 10^4.$$
(28)
However, in the $`B_d`$ system $`|\mathrm{\Delta }\mu |x_d0.7`$, meaning that $`|\mathrm{\Delta }\chi |`$ cannot compete with $`\mathrm{\Gamma }`$ nor with $`|\mathrm{\Delta }\mu |`$.
The situation in the $`B_s`$ is very different. From the point of view of the lifetime difference it is actually closer in spirit to the kaon system . Indeed, in the $`B_s`$ system $`|\mathrm{\Delta }\mathrm{\Gamma }|/\mathrm{\Gamma }0.15`$ , and one may talk about two exponentials. However, $`B_s`$ contains no quarks of the first family and, hence, it would seem inappropriate to estimate the matter effects as we have done above.
In conclusion, we have analyzed the possibility of detecting matter effects in the $`D^0\overline{D^0}`$ system. We conclude that these effects are, in principle, accessible. Unfortunately, when the time dependence is followed only for a few lifetimes, the matter effects and the mixing effects decouple. As a result, matter effects cannot be used to provide a new handle on $`x_D`$ or $`y_D`$.
###### Acknowledgements.
I am indebted to L. Lavoura and H. R. Quinn for reading the manuscript and for useful suggestions. I would like to thank H. R. Quinn, Th. Schietinger, and A. E. Snyder for discussions and for their collaboration in a related subject. This work is supported in part by the Department of Energy under contract DE-AC03-76SF00515. The work of J. P. S. is supported in part by Fulbright, Instituto Camões, and by the Portuguese FCT, under grant PRAXIS XXI/BPD/20129/99 and contract CERN/S/FIS/1214/98.
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# A framework for perturbations and stability of differentially rotating stars
## 1 Introduction
The study of oscillations of stars is an important and exciting field of current astrophysics. For instance through period-luminosity and period-radius relationships variable stars provide important ‘yardsticks’ for distance measurements in the universe. Their observation yield important information about the interior of stars, like the equation of state of the matter, which is otherwise hard to obtain. Further, Neutron star pulsations may be a source of gravitational radiation detectable for experiments like LIGO, VIRGO and GEO 600 in the near future.
On the other hand, it is probably fair to say that there has not been very much work on the mathematical foundations of the theory of stellar oscillations. In the non relativistic limit there is a well-known framework for the description of oscillations of nonrotating stars , (see for a rigorous version). It is important to note that even on that level it turns out that the governing operator of the spheroidal oscillations belongs to a class of operators which were apparently (apart from special cases considered in ) previously unconsidered in operator theory . Somewhat surprisingly it also turned out that differently to radial oscillations these operators don’t have a compact resolvent. Hence the well developed perturbation theory for such operators cannot be applied and a corresponding theory for the new type of operators still has to be developed. For rotating stars there is little known about the relevant operators apart from abstract properties (like the symmetry, semiboundedness and continuity of associated operators ), indication of a continuous part in the spectrum and instabilities caused by so called ‘r-modes’ (also called ‘quasi-toroidal modes’) . To my knowledge there is no consideration of these operators in sufficient detail. Indeed, a large part of the present paper considers the more modest first step of such an investigation, namely to identify and determine which operators should be considered in that case.
The governing equation for linearized adiabatic oscillations of a stationary differentially rotating perfect-fuid star in an inertial frame $`(t,x)`$ is
$$\frac{^2\xi }{t^2}+B^{}\frac{\xi }{t}+C^{}\xi =0,$$
(1)
where
$$B^{}\frac{\xi }{t}:=2(v)\frac{\xi }{t},$$
(2)
$`C^{}\xi :=`$ $``$ $`{\displaystyle \frac{1}{\rho }}\left(p\mathrm{\Gamma }_1\xi \right)+(v)^2\xi +{\displaystyle \frac{1}{\rho }}\left[\left((p)\xi \right)+(\xi )p\right]`$ (3)
$``$ $`\mathrm{\Phi }_\xi +{\displaystyle \underset{j,k=1}{\overset{3}{}}}\left({\displaystyle \frac{1}{\rho }}{\displaystyle \frac{^2p}{x_jx_k}}+{\displaystyle \frac{^2\psi }{x_jx_k}}\right)\xi _ke_j,`$
$$\mathrm{\Phi }_\xi (t,x):=G_\mathrm{\Omega }\frac{[(\rho \xi )](t,y)}{|xy|}d^{\mathrm{\hspace{0.17em}3}}y,\psi (x):=G_\mathrm{\Omega }\frac{\rho (y)}{|xy|}d^{\mathrm{\hspace{0.17em}3}}y,x\mathrm{\Omega },$$
(4)
$`\xi `$ is the Lagrangian displacement vector field, for $`j\{1,2,3\}`$ the symbol $`e_j`$ denotes the canonical unit vector in the direction of $`x_j`$, $`\mathrm{\Omega }`$ is the (bounded open) volume of the star, and $`v,p,\rho ,\mathrm{\Gamma }_1`$ are the velocity field, pressure, density and the adiabatic index functions of the background star satisfying the equations of momentum and mass conservation
$$(v)v=\left(\frac{1}{\rho }p+\psi \right),(\rho v)=0.$$
(5)
and an equation of state. In addition to (1) the variation $`\delta p`$ of the pressure has to vanish at the boundary $`\mathrm{\Omega }`$ of the star, i.e.,
$$\underset{yx}{lim}\left(p\mathrm{\Gamma }_1\xi \right)(y)=0$$
(6)
for all $`x\mathrm{\Omega }`$.
The remarkable paper of Dyson and Schutz provides a framework for deciding the stability of the solutions of (1). In the following this paper is referred to as DS. Compared to previous frameworks given in , the main step forward in that paper is the fact that it relates the stability of the system directly to the growth properties of the perturbations in time as is usual for nonrotating stars (see, e.g., , , , , , ). Moreover the paper shows that these growth properties are governed by spectral properties of the generator of time evolution. This greatly simplifies the stability discussion. Unfortunately, the approach still has some drawbacks, and in the present paper will be given a varied framework which overcomes those problems. Moreover here are given for the first time sufficient criteria for stability in the form of inequalities which have to be satisfied by the coefficients of an operator polynomial. These criteria show that a negative canonical energy does not necessarily indicate an instability of the star. It is still unclear whether these criteria are strong enough to prove stability for realistic stars.
A rough discussion of the approach of Dyson and Schutz is now given. The paper considers axisymmetric solutions of the form
$$\xi (t,x)=exp(im\phi )\xi _m(t,r,\theta ),$$
(7)
where $`r,\theta ,\phi `$ are spherical coordinates and $`m`$. Inserting this ansatz into (1) leads to an equation of the same structure with induced operators $`B_m^{}`$ $`C_m^{}`$. The index $`m`$ is supressed in the following discussion. A Hilbert space $`H^{}`$ (here $`X`$) for the data is chosen such that, both, $`B^{}`$ becomes continuous and antisymmetric and $`C^{}`$ becomes symmetric. In the nonrotating limit $`H^{}`$ goes over into the usual Hilbert space used in the stability discussion for spherically symmetric stars. A physically reasonable condition on the background model is given which leads to a lower bounded $`C^{}`$. Assuming that condition $`C^{}`$ is substituted by its so called Friedrichs extension. This is an abstractly defined self-adjoint extension which exists for every densely defined linear symmetric and semibounded operator in Hilbert space (see e.g. Vol. II). <sup>1</sup><sup>1</sup>1 The importance of choosing a self-adjoint extension of $`C^{}`$ can be seen in the limit of no rotation. In , it is shown that for polytropic stars with a polytropic index $`n<1`$ there is an infinite number of different self-adjoint extensions which all lead to a well-posed initial value problem for the wave equation. In the standard way the resulting wave equation is written as a first order system in time for $`\xi (r,\theta ,t)`$ and $`(\xi /t)(r,\theta ,t)`$. The initial value problem of the system is studied. The Hilbert space of the data is chosen as $`H^{\mathrm{\hspace{0.17em}2}}`$ with the induced ‘euclidean’ scalar product. However it is noticed that this is physically not meaningful, because the scalar product has no physical interpretation and is not even dimensionally correct. From the first order system the linear operator $`T`$ generating time evolution is read of and it is shown that its spectrum is equal to the spectrum of a quadratic operator polynomial generated by $`B^{}`$ and $`C^{}`$ . Moreover the resolvent of $`T`$ can be given in terms of the inverses of the operator polynomial. That information along with estimates on the resolvent of the operator polynomial are used to give an estimate on the spectrum of $`T`$. In general that estimate is not strong enough to decide the question of stability of the system. From these estimates it is further shown that there is a solution of the initial value problem for the system corresponding to elements of the domain of $`T`$. The uniqueness of the solution is not shown. The authors remark that they could not show that $`T`$ is the generator of a strongly continuous semigroup and as a consequence the results of standard semigroup theory could not be used. Finally, the completeness of normal modes of the system is discussed.
From the description the reader might have noticed that in the derivation of these results only abstract properties like ‘continuity’, ‘symmetry’, ‘semiboundedness’ and ‘self-adjointness’ of $`B^{}`$ and $`C^{}`$ play a role. This is indeed true and is the reason why that approach is called here a ‘framework’. The same also applies to the approach here. As a consequence these frameworks can be used to describe a lot more physical systems than stellar oscillations. The main ingredient for such an application is a system of wave equations which is second order in time (with or without first order time derivatives) and which is not explicitly time dependent. The ‘coefficients’ in that system can be (not necessarily local) linear operators with certain abstract properties. For this reason we abandon in Section 2 any reference to rotating stars and just consider abstract wave equations of type (1). Having this in mind might also provide a better understanding of some of the statements below.
After this digression the discussion of DS is continued. The main problem of the approach comes from the chosen Hilbert space along with a scalar product which is not related to any physical quantity and not dimensionally correct. Of course the latter could be remedied by first introducing a dimensionless time coordinate. But experience tells that this should not be essential at such an early stage. Also it is known that the use of a suitable Hilbert space decides whether semigroup theory can be applied or not. So it is very likely that the use of $`H^{\mathrm{\hspace{0.17em}2}}`$ is responsible for the fact that semigroup theory could not be applied. Indeed a different choice of the Hilbert space will turn out to be the key to the results of this paper.
Another point which was not addressed in DS is the fact that in addition to (1) the boundary condition (6) has to be satisfied that the Lagrangian variation $`\delta p`$ of the pressure vanishes at the surface of the star. Indeed it has been shown in for the limit of no rotation that for a polytropic equation of state with polytropic index $`n<1`$ there is a infinite number of different self-adjoint extensions of $`C^{}`$, which all lead to different initial value formulations for the wave equation. Moreover it has been shown that the condition of a vanishing $`\delta p`$ at the surface of the star picks exactly one of these self-adjoint extensions. Of course the choice of the Friedrichs extension of $`C^{}`$, is equivalent to posing a boundary condition. But because of the abstractness of this extension it is not obvious and has to be investigated whether it is compatible with (6). To my knowledge this has been shown only for the case of radial oscillations of spherically symmetric stars in . This point will not be pursued any further in this paper.
The approach in this paper is similar to that of Dyson and Schutz. The point of departure is in the choice of the Hilbert space for the initial data of the first order system. Here a space $`Y`$ is chosen , which is in general a proper subspace of $`H^{\mathrm{\hspace{0.17em}2}}`$. Moreover a different and dimensionally correct scalar product is chosen. The square of the induced norm of the initial data is a positive definite part of the corresponding canonical energy of the system. For this $`C^{}`$ is split into sum of a strictly positive self-adjoint operator $`A`$ and a ‘rest’ $`C`$. Of course such a decomposition is not unique but it can be shown (see Lemmas 14 and 15 in Section 2) that trivial rescalings all lead to the same set $`Y`$ along with equivalent norms on $`Y`$. In particular such changes lead to theories which are related by a similarity transformation and hence the outcome of the stability discussion is not affected. In general the canonical energy cannot be used as a norm for $`Y`$ because it is not always positive definite. In situations where it is $`C`$ can be chosen as zero. In the limit of no rotation where $`B^{}=0`$ and the operator $`C^{}`$ is semibounded the approach here reduces to the approach in Vol. II (see the proposition at the beginning of chapter X.13) for classical wave equations.
A major consequence of the change is that it allows the use of semigroup theory which is a standard and well developed tool in particular in the theory of partial differential equations (see e.g. ,,, ,, and the cited references therein). This simplifies the reasoning a lot, because it can be and will be built on those results. In particular here the operator $`G_+`$ which corresponds to $`T`$ in DS generates a strongly continuous group of bounded transformations and hence the well-posedness of the initial value problem for the first order system follow from abstract semigroup theory. At the same time a considerable generalization is achieved. The operator $`B^{}`$ (here denoted by $`iB`$) can be unbounded and not antisymmetric. Moreover $`C^{}`$ has not to be assumed symmetric. The restrictions imposed on these operators are the following. The operator $`C^{}`$ has to be of the form $`A+C`$ where $`A`$ is some densely defined and strictly positive self-adjoint operator in $`X`$ and $`C`$ is a relatively bounded perturbation of the positive square root $`A^{1/2}`$ of $`A`$. In addition $`B^{}`$ has to be a relatively bounded perturbation of $`A^{1/2}`$ with relative bound smaller than $`1`$. Finally, $`B^{}`$ has to be antisymmetric or continuous, but not necessarily both. All these conditions are trivially satisfied for the case of axisymmetric solutions of (1) considered by Dyson and Schutz. Whether this generalization is sufficient to provide a framework for (1) and not only for its axisymmetric form is not yet clear. For this it seems necessary that $`C^{}`$ given by (3) is semibounded and this is still open. The reason for considering also more general situations with nonantisymmetric $`B^{}`$ and nonsymmetric $`C^{}`$ is that the framework here will also be used in a future paper in the stability discussion of the Teukolsky equations on a Kerr background where this is the case. A further important advantage of the approach here is that it can be shown (see Theorem 3) that the dominating part of $`G_+`$ (but in general not $`G_+`$ itself) is self-adjoint. Using perturbation theory this gives important information on the spectrum of $`G_+`$ and is also the basis of the proof that $`G_+`$ is the generator of a strongly continuous group of bounded transformations. (See Theorem 7) Further it is the basis for another result (see Corollary 12) having no counterpart in DS namely the conservation of the ‘canonical energy’ $`E`$. On the other hand it turns out that the spectrum of $`G_+`$ is the same as of $`T`$. In particular that spectrum is given by the spectrum of the same operator polynomial $`C^{}\lambda B^{}+\lambda ^2,\lambda `$ (See Theorem)
A plausible definition for the stability of a rotating star is the following. The system is stable if and only if the semigroup $`T_+(t),t[0,\mathrm{})`$ generated by $`G_+`$ is bounded. Note that this definition is invariant to similarity transformations and hence not so sensitive to changes of Hilbert space like one only invariant under unitary transformations.<sup>2</sup><sup>2</sup>2Such a definition would be given for instance by the demand that the semigroup should be contractive, i.e., that the norms of the semigroup elements are $`1`$. From semigroup theory one has then the following.
1. The system is unstable if $`G_+`$ has a spectral value with real part smaller than zero.
2. For a stable system the corresponding spectrum of $`G_+`$ is contained in the closed right half plane of the complex plane.
3. From only the fact that the spectrum of $`G_+`$ is part of the closed right half plane of the complex plane, it cannot be concluded that the system is stable. <sup>3</sup><sup>3</sup>3 For a counterexample compare for instance the note after the proof of Corollary 9.
4. The system is stable if the real part of all ‘expectation values’
$$(\xi |G_+\xi )$$
(8)
is positive ($`0`$) for all elements $`\xi `$ from the domain (or a core) of $`G_+`$.<sup>4</sup><sup>4</sup>4Then $`G_+`$ generates a contraction semigroup.
Point 1 gives a sufficient but not necessary condition for instability. Note that this condition is invariant under similarity transformations. Moreover because of Theorem 13 it is equivalent to the condition that there is complex number $`\lambda `$ with real part smaller than zero such that
$$C^{}\lambda B^{}+\lambda ^2$$
(9)
is not bijective. This reduces in the nonrotating case to the condition that $`C^{}`$ is strictly negative, which is a well known sufficient condition for instability. An important final observation is that from the existence of such a $`\lambda `$ follows the existence of an element $`\xi `$ from the Hilbert space such that the corresponding function of norms $`|T_+(t)\xi |,t`$ grows exponentially for large times.<sup>5</sup><sup>5</sup>5This is easily seen for instance by using Theorem 4.1 in chapter 4 of . Here it is important to remember that in general the spectrum of $`G_+`$ does not only consist of ‘eigenvalues’ (for which this statement is of course trivially satisfied) but also values $`\mu `$ for which the map $`G_+\mu `$ is just not onto. Such values are often from a continuous part of the spectrum. Hence the existence of such a $`\lambda `$ leads to a much stronger kind of instability.
Point 4 gives a sufficient but not necessary condition for stability. It is appealing because it is of the form of an inequality, which is more easily accessible than the spectrum of $`G_+`$. On the other hand it is strong and not invariant to similarity transformations. It turns out to be equivalent to $`C^{}`$ being strictly positive, i.e., that the spectrum of this operator consists only of positive real numbers different from zero. Note that this reduces to a known sufficient condition for stability in the nonrotating case. But for such stars it can be satisfied only for radial oscillations (for instance this is the case for constant $`\mathrm{\Gamma }_1>4/3`$), but not for nonradial oscillations.<sup>6</sup><sup>6</sup>6This is obvious since the spectrum of the trivial toroidal oscillations is $`\{0\}`$. But in it has also been shown that $`0`$ is in the spectrum of spheroidal oscillations. Note that in the limit of no rotation the trivial toroidal oscillations give rise to solutions of (1,6) whose norm increases linear in time for large times. Hence applying the stability definition above to that limit would lead to an ‘unstable star’. Of course, these oscillations can be excluded in that case just by considering a reduced operator.
The following two stability criteria are new. They will be proven in Theorem 17.
1. If $`B^{}`$ and $`C^{}`$ are such that
$$<\xi |C^{}\xi >\frac{1}{4}<\xi |B^{}\xi >^20$$
(10)
for all $`\xi `$ from the domain of $`C^{}`$ such that $`\xi =1`$ then the spectrum of $`G_+`$ is purely imaginary.
2. If the operator
$$C^{}\frac{ib}{2}B^{}\frac{b^2}{4}$$
(11)
is positive for some $`b`$ then the spectrum of $`G_+`$ is purely imaginary and there are $`K0`$ and $`t_00`$ such that
$$|T_+(t)|Kt$$
(12)
for all $`tt_0`$.
Note for the first point that $`(1/4)<\xi |B^{}\xi >^2`$ is positive, because of the antisymmetry of $`B^{}`$. Also note in this connection that in DS it has been shown that $`C^{}(1/4)B^{\mathrm{\hspace{0.17em}2}}`$ is bounded from below uniformly in m. Unfortunately, in general this does not imply (10).
It is still unclear whether these criteria are strong enough to prove stability for realistic stars. On the other hand the second criterium has been sucessfully applied in the stability discussion of the Kerr metric where the master equation governing perturbations is of the form (1), too.
## 2 The framework
This section developes the initial value formulation for abstract differential equations of the form (1). It is self-contained and necessarily very technical. The reader who is not interested in the excessive mathematical details given here is referred to the introduction. The used nomenclature can be found in standard textbooks on Functional analysis. Vol. I,
Before going into the mathematical details it is explained about the meaning of the individual results of this section. The section is based on the assumptions General Assumption 1 and General Assumption 4 on three operators $`A`$, $`B`$ and $`C`$. A different form of General Assumption 1 which is more convenient for applications can be given in the obvious way using Lemma 18. Definition 2 gives the Hilbert space $`Y`$ which is used here instead of the Hilbert space in DS. A rigorous form (55) of (1) along with the existence and uniqueness of the solution corresponding to initial values is given in Theorem 11. Corollary 12 gives the corresponding ‘energy’ along with an identity for its time derivative. The analogue $`G_+`$ of the generator $`T`$ in DS is given in Definition 5. Theorem 3 proves that the ‘dominating parts’ of $`G_+`$ are self-adjoint. In Theorem 7 it is proved that under General Assumption 1 and General Assumption 4, both, $`G_+`$ and $`G_+`$ are generators of strongly continuous semigroups $`T_+`$ and $`T_{}`$, resp. Theorem 13 shows the identity of the spectrum of $`G_+`$ with the spectrum of an operator polynomial generated by the operators $`B`$ and $`A+C`$. Lemmas 14 and 15 show that certain simple rescalings of $`A`$ and $`C`$ which formally leave invariant (55) lead to theories which are related by a similarity transformation. Theorem 16 shows for a special case how these rescalings can be used to derive a better estimate for the growth of $`T_+`$ and $`T_{}`$ than the one induced by $`(\text{30})`$ in Lemma 6. Theorem 17 gives sufficient criteria for stability in the form of inequalities which have to be satisfied by the coefficients of the operator polynomial. Part (ii) of this Theorem has been sucessfully applied in the discussion of the stability of the Kerr metric.
The rest of this section contains the mathematical details.
###### Assumption 1
In the following let $`(X,<|>)`$ be a non trivial complex Hilbert space. Denote by $``$ the norm induced on $`X`$ by $`<|>`$. Further let $`A:D(A)X`$ be a densely defined linear self-adjoint operator in $`X`$ for which there is an $`\epsilon (0,\mathrm{})`$ such that
$$<\xi |A\xi >\epsilon <\xi |\xi >$$
(13)
for all $`\xi D(A)`$. Denote by $`A^{1/2}`$ the square root of $`A`$ with domain $`D(A^{1/2})`$. Further let be $`B:D(A^{1/2})X`$ a linear operator in $`X`$ such that for some $`a[0,1)`$ and $`b`$
$$B\xi ^2a^2A^{1/2}\xi ^2+b^2\xi ^2$$
(14)
for all $`\xi D(A^{1/2})`$. Finally, let $`C:D(A^{1/2})X`$ be linear and such that for some real numbers $`c`$ and $`d`$
$$C\xi ^2c^2A^{1/2}\xi ^2+d^2\xi ^2$$
(15)
for all $`\xi D(A^{1/2})`$.
Note that as a consequence of (13) the spectrum of $`A`$ is contained in the interval $`[\epsilon ,\mathrm{})`$. Hence $`A`$ is in particular positive and bijective and there is a uniquely defined linear and positive selfadjoint operator $`A^{1/2}:D(A^{1/2})X`$ such that $`(A^{1/2})^2=A`$. That operator is the so called square root of $`A`$. Further note that from its definition and the bijectivity of $`A`$ follows that $`A^{1/2}`$ is in particular bijective. This can be concluded for instance as follows. By using the fact that $`A^{1/2}`$ commutes with $`A`$ it easy to see that for every $`\lambda [0,\epsilon ^{1/2})`$ by $`(A^{1/2}+\lambda )(A\lambda ^2)^1`$ there is given the inverse to $`A^{1/2}\lambda `$. Hence the spectrum of $`A^{1/2}`$ is contained in the interval $`[\epsilon ^{1/2},\mathrm{})`$. All these facts will be used later on.
###### Definition 2
We define
$$Y:=D(A^{1/2})\times X$$
(16)
and $`(|):Y^2`$ by
$$(\xi |\eta ):=<A^{1/2}\xi _1|A^{1/2}\eta _1>+<\xi _2|\eta _2>$$
(17)
for all $`\xi =(\xi _1,\xi _2),\eta =(\eta _1,\eta _2)Y`$.
Then we have the following
###### Theorem 3
$`(Y,(|))`$ is a complex Hilbert space.
The operator $`H:D(A)\times D(A^{1/2})Y`$ in $`Y`$ defined by
$$H\xi :=(i\xi _2,iA\xi _1)$$
(18)
for all $`\xi =(\xi _1,\xi _2)D(A)\times D(A^{1/2})`$ is densely-defined, linear and self-adjoint.
The operator $`\widehat{B}:D(H)Y`$ defined by
$$\widehat{B}\xi :=(0,B\xi _2)$$
(19)
for all $`\xi =(\xi _1,\xi _2)D(H)`$ is linear. If $`B`$ is symmetric then $`\widehat{B}`$ is symmetric, too. If $`B`$ is bounded then $`\widehat{B}`$ is bounded, too, and the corresponding operator norms $`B`$ and $`|\widehat{B}|`$ satisfy
$$|\widehat{B}|B.$$
(20)
The sum $`H+\widehat{B}`$ is closed. If $`B`$ is symmetric then $`H+\widehat{B}`$ is self-adjoint.
The operator $`V:YY`$ defined by
$$V\xi :=(0,iC\xi _1)$$
(21)
for all $`\xi =(\xi _1,\xi _2)Y`$ is linear and bounded. The operator norm $`|V|`$ of $`V`$ satisfies
$$|V|(c^2+d^2/ϵ)^{1/2}.$$
(22)
Proof: (i): Obviously, $`(|)`$ defines a hermitean sesquilinear form on $`Y^2`$. That $`(|)`$ is further positive definite follows from the positive definiteness of $`<|>`$ and the injectivity of $`A^{1/2}`$. Finally, the completeness of $`(Y,||)`$, where $`||`$ denotes the norm on $`Y`$ induced by $`(|)`$, follows from the completeness of $`(X,)`$ together with the fact that $`A^{1/2}`$ has a bounded inverse. Here it is essentially used that $`0`$ is not contained in the spectrum of $`A`$. (ii): That $`D(A)\times D(A^{1/2})`$ is dense in $`Y`$ is an obvious consequence of the facts that $`D(A)`$ is a core for $`A^{1/2}`$ (see e.g. Theorem 3.24 in chapter V.3 of ) and that $`D(A^{1/2})`$ is dense in $`X`$. The linearity of $`H`$ is obvious. Also the symmetry of $`H`$ follows straighforwardly from the symmetry of $`A^{1/2}`$. By that symmetry one gets further for any $`\xi =(\xi _1,\xi _2)D(H^{})`$ and any $`\eta =(\eta _1,\eta _2)D(H)`$:
$`(H^{}\xi |\eta )`$ $`=`$ $`<\left(H^{}\xi \right)_1|A\eta _1>+<\left(H^{}\xi \right)_2|\eta _2>`$
$`=(\xi |H\eta )`$ $`=`$ $`<i\xi _2|A\eta _1>+<iA^{1/2}\xi _1|A^{1/2}\eta _2>`$ (23)
and from this by using that $`A`$ is bijective and $`A^{1/2}`$ is self-adjoint that $`\xi _1D(A)`$ and
$$\left(H^{}\xi \right)_1=i\xi _2,\left(H^{}\xi \right)_2=iA\xi _1.$$
(24)
Hence $`H`$ is an extension of $`H^{}`$ and thus $`H=H^{}`$. (iii): The linearity of $`\widehat{B}`$ is obvious. Also it is straightforward to see that $`\widehat{B}`$ is symmetric if $`B`$ is symmetric. If $`B`$ is bounded then
$$|\widehat{B}\xi |^2=B\xi _2^2B^2\xi _2^2B^2|\xi |^2$$
(25)
for all $`\xi =(\xi _1,\xi _2)D(H)`$. Hence $`\widehat{B}`$ is also bounded and $`|\widehat{B}|`$, $`B`$ satisfy the claimed inequality. (iv): Obviously, (14) implies
$$|\widehat{B}\xi |^2a^2|H\xi |^2+b^2|\xi |^2$$
(26)
for all $`\xi D(H)`$. From this it is easily seen that $`H+\widehat{B}`$ is closed (see, e.g., , Lemma V.3.5). Moreover in the case that $`B`$ (and hence by (iii) also $`\widehat{B}`$) is symmetric (26) implies according to the Kato-Rellich Theorem (see, e. g., Theorem X.12 in Vol. II) that $`H+\widehat{B}`$ is self-adjoint. For the application of these theorems the assumption $`a<1`$ made above is essential. (v) The linearity of $`V`$ is obvious. For any $`\xi =(\xi _1,\xi _2)Y`$ one has
$`|V\xi |^2`$ $`=`$ $`iC\xi _1^2c^2A^{1/2}\xi _1^2+d^2\xi _1^2`$ (27)
$`=`$ $`c^2A^{1/2}\xi _1^2+d^2(A^{1/2})^1A^{1/2}\xi _1^2(c^2+d^2/ϵ)|\xi |^2.`$
In the last step it has been used that
$$(A^{1/2})^11/\sqrt{\epsilon }.$$
(28)
This follows by an application of the spectral theorem (see, e.g. Theorem VIII.5 in Vol. I) to $`A^{1/2}`$. Since $`\xi `$ is otherwise arbitrary from (27) follows the boundedness of $`V`$ and the claimed inequality.
###### Assumption 4
In the following we assume in addition that $`B`$ is symmetric or bounded.
Note that condition (14) is trivially satisfied if $`B`$ is bounded. We define:
###### Definition 5
$$G_+:=i(H+\widehat{B}+V),G_{}:=i(H+\widehat{B}+V).$$
(29)
then
###### Lemma 6
The operators $`G_+`$ and $`G_{}`$ are closed and quasi-accretive. In particular
$$Re(\xi |G\xi )(\mu _B+|V|)(\xi |\xi )$$
(30)
for $`G\{G_+,G_{}\}`$ and all $`\xi D(H)`$. Here $`Re`$ denotes the real part and
$$\mu _B:=\{\begin{array}{cc}0& \text{ if }B\text{ is symmetric}\hfill \\ B& \text{if }B\text{ is bounded}\hfill \end{array}.$$
(31)
Proof: That $`G_+`$ and $`G_{}`$ are closed is an obvious consequence of (iv) and (v) of the previous theorem. Further if $`B`$ is symmetric one has because of (iv) and (v) of the preceeding theorem
$$Re(\xi |G_\pm \xi )=Re(\xi |iV\xi )|(\xi |iV\xi )||V|(\xi |\xi )$$
(32)
for all $`\xi D(H)`$. Similarly, if $`B`$ is bounded one has because of (ii),(iii),(iv), (20)
$$Re(\xi |G_\pm \xi )=Re(\xi |i(\widehat{B}+V)\xi )|(\xi |i(\widehat{B}+V)\xi )|(B+|V|)(\xi |\xi )$$
(33)
for all $`\xi D(H)`$. Hence in both cases $`G_+`$ and $`G_{}`$ are quasi-accretive.
###### Theorem 7
The operators $`G_+`$ and $`G_{}`$ are infinitesimal generators of strongly continuous semigroups $`T_+:[0,\mathrm{})L(Y,Y)`$ and $`T_{}:[0,\mathrm{})L(Y,Y)`$, respectively. If $`\mu _\pm `$ are such that
$$Re(\xi |G_\pm \xi )\mu _\pm (\xi |\xi )$$
(34)
for all $`\xi D(H)`$ the spectra of $`G_+`$ and $`G_{}`$ are contained in the half-plane $`[\mu _+,\mathrm{})\times `$ and $`[\mu _{},\mathrm{})\times `$, respectively, and
$$|T_+(t)|\mathrm{exp}(\mu _+t),|T_{}(t)|\mathrm{exp}(\mu _{}t)$$
(35)
for all $`t[0,\mathrm{})`$.
Proof: Obviously, by the Lumer-Phillips theorem (see, e.g., Theorem X.48 in Vol. II of ) and the preceeding lemma the theorem follows if we can show that there is a real number $`\lambda <\mathrm{min}\{\mu _+,\mu _{}\}`$ such that $`G_\pm \lambda `$ has a dense range in $`Y`$. For that proof let be $`\xi `$ some element of $`D(H)`$ and $`\lambda `$ any real number such that $`|\lambda ||V|^2`$. Then we get from the symmetry of $`H`$
$$|(Hi\lambda )\xi |^2=|H\xi |^2+\lambda ^2|\xi |^2$$
(36)
and
$$|(Hi\lambda )\xi |\mathrm{max}\{|H\xi |,|\lambda |^{1/2}|V\xi |\}.$$
(37)
Using these identities together with (14)
$`|(\widehat{B}+V)\xi |^2`$ $``$ $`|\widehat{B}\xi |^2+2|\widehat{B}\xi ||V\xi |+|V\xi |^2`$ (38)
$``$ $`a^2|H\xi |^2+2|\widehat{B}\xi ||V\xi |+(b^2+|V|^2)|\xi |^2`$
$``$ $`a^2|H\xi |^2+2a|H\xi ||V\xi |+(b+|V|)^2|\xi |^2`$
$``$ $`a^2|(Hi\lambda )\xi |^2+2a|(Hi\lambda )\xi ||V\xi |+[(b+|V|)^2a^2\lambda ^2]|\xi |^2`$
$``$ $`a(a+2|\lambda |^{1/2})|(Hi\lambda )\xi |^2+[(b+|V|)^2a^2\lambda ^2]|\xi |^2.`$
Hence for any real $`\lambda `$ with
$$|\lambda |>\mathrm{max}\{|V|^2,4(1a)^2,(b+|V|)/a,|\mu _+|,|\mu _{}|\},$$
(39)
where we assume without restriction that $`a>0`$, we get
$$|(\widehat{B}+V)\xi |a^{}|(Hi\lambda )\xi |$$
(40)
where $`a^{}`$ is some real number from $`[0,1)`$. Since $`\xi D(H)`$ is otherwise arbitrary, we conclude that
$$(\widehat{B}+V)(Hi\lambda )^1$$
(41)
defines a bounded linear operator on $`Y`$ with operator norm smaller than $`1`$. Since
$$H+\widehat{B}+Vi\lambda =\left(1+(\widehat{B}+V)(Hi\lambda )^1\right)(Hi\lambda )$$
(42)
we conclude that $`H+\widehat{B}+Vi\lambda `$ is bijectice and hence also that $`G_+\lambda `$ and $`G_{}\lambda `$ are both bijective. Hence the theorem follows.
We note that General Assumption 4 has been used only to conclude that $`G_+`$ and $`G_{}`$ are both quasi-accretive. Now it is easy to see that if $`B`$ is in addition such that $`iB`$ is quasi-accretive (but not necessarily bounded or antisymmetric) then $`i\widehat{B}`$ and hence also $`G_+`$ are quasi-accretive, too. As a consequence we have the following
###### Corollary 8
Instead of General Assumption 4 let $`B`$ be such that $`iB`$ is quasi-accretive. Then $`G_+`$ is the infinitesimal generator of a strongly continuous semigroup $`T_+:[0,\mathrm{})L(Y,Y)`$. If $`\mu _+`$ is such that
$$Re(\xi |G_+\xi )\mu _+(\xi |\xi )$$
(43)
for all $`\xi D(H)`$ the spectrum of $`G_+`$ is contained in the half-plane $`[\mu _+,\mathrm{})\times `$ and
$$|T_+(t)|\mathrm{exp}(\mu _+t)$$
(44)
for all $`t[0,\mathrm{})`$.
Theorem 7 has the following
###### Corollary 9
By
$$T(t):=\{\begin{array}{cc}T_+(t)& \text{for }t0\hfill \\ T_{}(t)& \text{for }t<0\hfill \end{array}$$
(45)
for all $`t`$ there is defined a strongly continuous group $`T:L(Y,Y)`$.
For every $`t_0`$ and every $`\xi D(G_+)`$ there is a uniquely determined differentiable map $`u:Y`$ such that
$$u(t_0)=\xi $$
(46)
and
$$u^{}(t)=G_+u(t)$$
(47)
for all $`t`$. Here $`^{}`$ denotes differentiation of functions assuming values in $`Y`$.
The function $`(u|u):`$ defined by
$$(u|u)(t):=(u(t)|u(t)),t$$
(48)
is differentiable and
$$(u|u)^{}(t)=2Re(u(t)|G_+u(t))$$
(49)
for all $`t`$.
Proof: The corollary follows from Theorem 7 by standard results of semigroup theory. For instance, see section 1.6 in for (i) and section IX.3 in for (ii). (iii) is an obvious consequence of (ii).
Note in particular the the special case <sup>7</sup><sup>7</sup>7Such cases are easy to construct. that there is a non trivial element $`\eta `$ in the kernel of $`A+C`$ for which there is $`\xi D(A)`$ such that
$$(A+C)\xi =iB\eta .$$
(50)
Then by
$$u(t):=(\xi +t\eta ,\eta ),t$$
(51)
there is given a growing solution of (47).
The following lemma is needed in the formulation of the subsequent theorem.
###### Lemma 10
By
$$\xi _{A^{1/2}}:=A^{1/2}\xi ,\xi D(A^{1/2})$$
(52)
there is defined a norm $`_{A^{1/2}}`$ on $`D(A^{1/2})`$. Moreover
$$W_1:=(D(A^{1/2}),_{A^{1/2}})$$
(53)
is complete.
Proof: The lemma is a trivial consequence of the completeness of $`X`$ and the bijectivity of $`A^{1/2}`$.
###### Theorem 11
Let be $`t_0`$, $`\xi D(A)`$ and $`\eta D(A^{1/2})`$. Then there is a uniquely determined differentiable map $`u:W_1`$ with
$$u(t_0)=\xi \text{and}u^{}(t_0)=\eta $$
(54)
and such that $`u^{}:X`$ is differentiable with
$$(u^{})^{}(t)+iBu^{}(t)+(A+C)u(t)=0$$
(55)
for all $`t`$.
Proof: For this let be $`v=(v_1,v_2):Y`$ be such that
$$v(t_0)=(\xi ,\eta )$$
(56)
and
$$v^{}(t)=G_+v(t),t.$$
(57)
Such $`v`$ exists according to Corollary 9 (ii). Using the continuity of the canonical projections of $`Y`$ onto $`W_1`$ and $`X`$ it is easy to see that $`u:=v_1`$ is a differentiable map into $`W_1`$ such that $`u^{}:X`$ is differentiable and such that (54), (55) are both satisfied. On the other hand if $`u:W_1`$ has the properties stated in the corollary it follows by the continuity of the canonical imbeddings of $`W_1`$, $`X`$ into $`Y`$ that $`w:=(u,u^{})`$ satisfies both equations (56) and (57). Then $`u=v_1`$ follows by Corollary 9 (ii).
###### Corollary 12
In addition to the assumptions made let $`C`$ be in particular bounded. <sup>8</sup><sup>8</sup>8Note that in this case (15) is trivially satisfied. Further let $`u:W_1`$ be differentiable with a differentiable derivative $`u^{}:X`$ and such that (55) holds. Finally, define $`E_u:`$ by
$$E_u(t):=\frac{1}{2}\left(<u^{}(t)|u^{}(t)>+<u(t)|(A+Re(C))u(t)>\right).$$
(58)
Then $`E_u`$ is differentiable and
$`E_u^{}(t)=\{\begin{array}{cc}Im<u(t)|Im(C)u^{}(t)>& \text{for symmetric }B\hfill \\ \frac{1}{2}<u^{}(t)|Im(B)u^{}(t)>Im<u(t)|Im(C)u^{}(t)>& \text{for bounded }B\hfill \end{array}`$ (61)
for all $`t`$, where for any bounded linear operator $`F`$ on $`X`$:
$$Re(F):=\frac{1}{2}\left(F+F^{}\right),Im(F):=\frac{1}{2i}\left(FF^{}\right).$$
(62)
Proof: For this define $`v:=(u,u^{})`$. Then according to the preceeding proof $`v`$ satisfies (57). For a symmetric $`B`$ it follows by Corollary 9 and Theorem 3 (iv) that
$`(v|v)^{}(t)`$ $`=`$ $`2Re(v(t)|iVv(t))`$
$`=`$ $`<u^{}(t)|Cu(t)><Cu(t)|u^{}(t)>`$
$`=`$ $`<u|Re(C)u>^{}(t)2Im<u(t)|Im(C)u^{}(t)>`$
for all $`t`$. In the last step it has been used that $`u`$ is also differentiable with the same derivative viewed as map with values in $`X`$. This follows from the fact the canonical imbedding of $`W_1`$ into $`X`$ is continuous since $`A^{1/2}`$ is bijective. Further the definition
$$<u|Re(C)u>(t):=<u(t)|Re(C)u(t)>,t$$
(64)
for the map $`<u|Re(C)u>:`$ has been used. Obviously, (61) follows from (2) by using definition (58). In this step also the symmetry of $`A^{1/2}`$ is used together with the fact that $`u`$ assumes values in $`D(A)`$. For a bounded $`B`$ by Corollary 9 and Theorem 3 (ii) follows that
$`(v|v)^{}(t)=`$ $`2Re(v(t)|i(\widehat{B}+V)v(t))`$
$`=`$ $`2Im<u^{}(t)|Bu^{}(t)>`$ (65)
$`<u^{}(t)|Cu(t)><Cu(t)|u^{}(t)>`$
$`=`$ $`2<u^{}(t)|Im(B)u^{}(t)>`$
$`<u|Re(C)u>^{}(t)2Im<u(t)|Im(C)u^{}(t)>`$
for all $`t`$. Obviously, (61) follows from (2) by using definition (58).
The next theorem relates the spectrum of $`G_+`$ to the spectrum of the so called operator polynomial $`A+C\lambda B\lambda ^2`$, where $`\lambda `$ runs through the complex numbers.
###### Theorem 13
Let $`\lambda `$ be some complex number.
Then $`H+\widehat{B}+V\lambda `$ is not injective if and only if $`A+C\lambda B\lambda ^2`$ is not injective. If $`H+\widehat{B}+V\lambda `$ is not injective then
$$\mathrm{ker}(H+\widehat{B}+V\lambda )=\{(\xi ,i\lambda \xi ):\xi \mathrm{ker}(A+C\lambda B\lambda ^2)\}.$$
(66)
Further $`H+\widehat{B}+V\lambda `$ is bijective if and only if $`A+C\lambda B\lambda ^2`$ is bijective. If $`H+\widehat{B}+V\lambda `$ is bijective then for all $`\eta =(\eta _1,\eta _2)Y`$:
$$(H+\widehat{B}+V\lambda )^1\eta =(\xi ,i(\lambda \xi +\eta _1)),$$
(67)
where
$$\xi =(A+C\lambda B\lambda ^2)^1[(B+\lambda )\eta _1i\eta _2].$$
(68)
Proof: (i) If $`H+\widehat{B}+V\lambda `$ is not injective and $`\xi =(\xi _1,\xi _2)\mathrm{ker}(H+\widehat{B}+V\lambda )`$ it follows from the definitions in theorem 3 that
$$\xi _2=i\lambda \xi _1,(A+C\lambda B\lambda ^2)\xi _1=0$$
(69)
and hence also that $`A+C\lambda B\lambda ^2`$ is not injective. If $`A+C\lambda B\lambda ^2`$ is not injective it follows again from the definitions in theorem 3 that
$$(H+\widehat{B}+V\lambda )(\xi ,i\lambda \xi )=0$$
(70)
and hence also that $`H+\widehat{B}+V\lambda `$ is not injective. (ii) If $`H+\widehat{B}+V\lambda `$ is bijective it follows by (i) that $`A+C\lambda B\lambda ^2`$ is injective. For $`\eta X`$ and $`\xi =(\xi _1,\xi _2):=(H+\widehat{B}+V\lambda )^1(0,i\eta )`$ it follows from the definitions in theorem 3 that
$$(A+C\lambda B\lambda ^2)\xi _1=\eta $$
(71)
and hence that $`A+C\lambda B\lambda ^2`$ is also surjective. If $`A+C\lambda B\lambda ^2`$ is bijective it follows by (i) that $`H+\widehat{B}+V\lambda `$ is injective. Further if $`\eta =(\eta _1,\eta _2)Y`$ and $`\xi `$ is defined by (68) it follows from the definitions in theorem 3 that
$$(H+\widehat{B}+V\lambda )(\xi ,i(\lambda \xi +\eta _1))=\eta $$
(72)
and hence that $`H+\widehat{B}+V\lambda `$ is also surjective.
###### Lemma 14
Let be $`\epsilon ^{}<\epsilon `$ and
$$A^{}:=A\epsilon ^{},C^{}:=C+\epsilon ^{}.$$
(73)
Then
$$D(A^{1/2})=D(A^{1/2})$$
(74)
and for all $`\xi D(A^{1/2})`$
$$A^{1/2}\xi ^2=A^{1/2}\xi ^2+\epsilon ^{}\xi ^2.$$
(75)
The operators $`A^{},B`$ and $`C^{}`$ satisfy
$`<\xi |A^{}\xi >`$ $`(\epsilon \epsilon ^{})<\xi |\xi >`$
$`B\xi ^2`$ $`a^2A^{1/2}\xi ^2+(a^2\epsilon ^{}+b^2)\xi ^2`$
$`C^{}\xi ^2`$ $`|c|\left[|c|+2|\epsilon ^{}|(\epsilon \epsilon ^{})^{1/2}\right]A^{\mathrm{\hspace{0.17em}1}/2}\xi ^2+`$ (76)
$`\left[|\epsilon ^{}|+(c^2|\epsilon ^{}|+d^2)^{1/2}\right]^2\xi ^2`$
for all $`\xi D(A^{1/2})`$.
Proof: (i) First, since $`\epsilon ^{}<\epsilon `$ by (73) there is defined a linear self-adjoint and positive operator $`A^{}`$ in $`X`$. Obviously, using the symmetry of $`A^{1/2}`$ and $`A^{1/2}`$ (75) follows for all elements of $`D(A)`$. From this (74) and (75) follow straightforwardly by using the facts that $`D(A)`$ is a core for both, $`A^{1/2}`$ and $`A^{1/2}`$ (see e.g. Theorem 3.24 in chapter V.3 of ), that $`X`$ is complete and that both operators, $`A^{1/2}`$ and $`A^{1/2}`$ are closed. (ii) The first two inequalities are obvious consequences of the corresponding ones in General Assumption 1, the definition (73) and of (75). For the proof of the third we notice that from the first inequality along with an application of the spectral theorem (see, e.g. Theorem VIII.5 in Vol. I) to $`A^{\mathrm{\hspace{0.17em}1}/2}`$ follows that
$$(A^{\mathrm{\hspace{0.17em}1}/2})^11/\sqrt{\epsilon \epsilon ^{}}.$$
(77)
Further from General Assumption 1 and (75) one gets
$$C\xi ^2c^2A^{\mathrm{\hspace{0.17em}1}/2}\xi ^2+(c^2|\epsilon ^{}|+d^2)\xi ^2.$$
(78)
for all $`\xi D(A^{1/2})`$. From these inequalities we get
$`C^{}\xi ^2`$ $``$ $`C\xi ^2+2|\epsilon ^{}|C\xi \xi +\epsilon ^{\mathrm{\hspace{0.17em}2}}\xi ^2`$
$``$ $`c^2A^{\mathrm{\hspace{0.17em}1}/2}\xi ^2+(\epsilon ^{\mathrm{\hspace{0.17em}2}}+c^2|\epsilon ^{}|+d^2)\xi ^2+2|\epsilon ^{}|C\xi \xi `$
$``$ $`c^2A^{\mathrm{\hspace{0.17em}1}/2}\xi ^2+\left[|\epsilon ^{}|+(c^2|\epsilon ^{}|+d^2)^{1/2}\right]^2\xi ^2+2|\epsilon ^{}||c|A^{\mathrm{\hspace{0.17em}1}/2}\xi \xi `$
$``$ $`|c|\left[|c|+2|\epsilon ^{}|(\epsilon \epsilon ^{})^{1/2}\right]A^{\mathrm{\hspace{0.17em}1}/2}\xi ^2+\left[|\epsilon ^{}|+(c^2|\epsilon ^{}|+d^2)^{1/2}\right]^2\xi ^2`$
for all $`\xi D(A^{1/2})`$ and hence the third inequality.
As a consequence of (ii) the sequence $`X,A^{},B,C^{}`$ satisfies General Assumption 1. The corresponding $`Y`$ given by Definition 2 is because of (i) again given by (16). Moreover the corresponding norm $`||^{}`$ on $`Y`$ turns out to be equivalent to $`||`$. More precisely one has for every $`\epsilon ^{}0`$
###### Lemma 15
$$||^{}||\epsilon ^{1/2}(\epsilon \epsilon ^{})^{1/2}||^{}$$
(80)
and for every bounded linear operator $`F`$ on $`Y`$:
$$\epsilon ^{1/2}(\epsilon \epsilon ^{})^{1/2}|F|^{}|F|\epsilon ^{1/2}(\epsilon \epsilon ^{})^{1/2}|F|^{}.$$
(81)
Proof: The first inequality is a straightforward consequence of (75) and (77). The second inequality is a straightforward implication of the first. .
Note that the $`G_\pm `$ corresponding to the the sequence $`X,A^{},B,C^{}`$ are the same for all $`\epsilon ^{}`$ ($`\epsilon ^{}`$ drops out of the definition). Moreover as a consequence of the preceding lemma the the topologies induced on $`Y`$ are equivalent. Hence the generated groups are the same, too. This will be used in the following important special case.
###### Theorem 16
Let be $`A=A_0+\epsilon `$, where $`A_0`$ is a densely defined linear positive self-adjoint operator and let be $`C=\epsilon `$. Then
$$|T_\pm (t)|e\epsilon ^{1/2}t\mathrm{exp}(\mu _Bt)$$
(82)
for all $`t\epsilon ^{1/2}`$.
Proof: For this let be $`\epsilon ^{}[0,\epsilon )`$ and define $`A^{}`$ and $`C^{}`$ as in Lemma 14. Hence
$$A^{}=A_0+\epsilon \epsilon ^{},C^{}=(\epsilon \epsilon ^{}).$$
(83)
Then from Theorem 3(v), Lemma 6 and Theorem 7 we conclude that
$$|T_\pm (t)|^{}\mathrm{exp}\left([\mu _B+(\epsilon \epsilon ^{})^{1/2}]t\right)$$
(84)
for all $`t`$ and hence by Lemma 15 that
$$|T_\pm (t)|\epsilon ^{1/2}(\epsilon \epsilon ^{})^{1/2}\mathrm{exp}\left([\mu _B+(\epsilon \epsilon ^{})^{1/2}]t\right).$$
(85)
For $`t\epsilon ^{1/2}`$ we get from this (82) by choosing
$$\epsilon ^{}:=\epsilon t^2._{\mathrm{}}$$
(86)
Note that in this special case (58) is conserved and positive.
We are now giving stability criteria.
###### Theorem 17
In addition let $`B`$ and $`C`$ be both symmetric.
Let $`A`$, $`B`$ and $`C`$ be such that
$$<\xi |(A+C)\xi >+\frac{1}{4}<\xi |B\xi >^20$$
(87)
for all $`\xi D(A)`$ with $`\xi =1`$. Then the spectrum of $`iG_+`$ is real.
In addition let $`B`$ and $`C`$ be both bounded and let $`A+C+(b/2)B(b^2/4)`$ be positive for some $`b`$. Then the spectrum of $`iG_+`$ is real and there are $`K0`$ and $`t_00`$ such that
$$|T(t)|K|t|$$
(88)
for all $`|t|t_0`$.
Proof: (i): First from General Assumption 1 and the assumed symmetry of $`B`$ and $`C`$ follows that, both, by $`A^{1/2}BA^{1/2}`$ and $`A^{1/2}CA^{1/2}`$ there is given a bounded symmetric and hence (by the theorem of Hellinger and Toplitz) also self-adjoint operator on $`X`$. Hence
$$A(\lambda ):=\lambda ^2A^1+\lambda A^{1/2}BA^{1/2}\left(1+A^{1/2}CA^{1/2}\right),\lambda $$
(89)
defines a self-adjoint operator polynomial in $`L(X,X)`$. In addition one has $`A^11/\epsilon `$ . Further for every $`\xi D(A^{1/2})`$ and $`\lambda `$
$$<\xi |A(\lambda )\xi >=<\eta |A^{1/2}A(\lambda )A^{1/2}\eta >=<\eta |(A+C\lambda B\lambda ^2)\eta >$$
(90)
where $`\eta :=A^{1/2}\xi D(A)`$. Now (87) implies that the roots of the polynomial $`<\eta |(A+C\lambda B\lambda ^2)\eta >,\lambda `$ are real. Hence by (90) the roots of $`<\xi |A(\lambda )\xi >,\lambda `$ are real, too. Since $`\xi D(A^{1/2})`$ is otherwise arbitrary and $`D(A^{1/2})`$ is dense in $`X`$ this implies also that $`<\xi |A(\lambda )\xi >`$ has only real roots for all $`\xi X`$. Hence (see , Lemma 31.1) the polynomial $`A(\lambda ),\lambda `$ is weakly hyperbolic and has therefore a real spectrum. As a consequence $`A(\lambda )`$ is bijective for all non real $`\lambda `$. Now for any such $`\lambda `$
$$A+C\lambda B\lambda ^2=A^{1/2}A_|(\lambda )A_|^{1/2},$$
(91)
where $`A_|^{1/2}`$ denotes the restriction of $`A^{1/2}`$, both, to $`D(A)`$ in domain and $`D(A^{1/2})`$ in range and $`A_|(\lambda )`$ denotes the restriction of $`A(\lambda )`$ to $`D(A^{1/2})`$, both, in domain and in range. For this note that $`A(\lambda )`$ leaves $`D(A^{1/2})`$ invariant. Further from the bijectivity of $`A^{1/2}`$, $`A(\lambda )`$ and (89) follows the bijectivity of $`A_|^{1/2}`$ and $`A_|(\lambda )`$, respectively and hence by (91) that $`A+C\lambda B\lambda ^2`$ is bijective. This is true for all non real $`\lambda `$ and hence it follows by Theorem 13 that the spectrum of $`iG_+`$ is real. (ii) So let $`B`$ and $`C`$ be both bounded and let $`A+C+(b/2)B(b^2/4)`$ be positive for some $`b`$. In addition let be $`ϵ`$ some real number greater than zero and define
$$A^{}:=A+C+(b/2)B(b^2/4)+\epsilon ,C^{}:=\epsilon ,B^{}:=Bb.$$
(92)
First it is observed that
$$D(A^{\mathrm{\hspace{0.17em}1}/2})=D(A^{1/2})$$
(93)
and that there exist nonvanishing real constants $`K_1`$ and $`K_2`$ such that
$$K_1^2A^{1/2}\xi ^2A^{\mathrm{\hspace{0.17em}1}/2}\xi ^2K_2^2A^{1/2}\xi ^2$$
(94)
for every $`\xi D(A^{1/2})`$. This can be proved as follows. Obviously, by the symmetry of $`A^{1/2}`$ and $`A^{\mathrm{\hspace{0.17em}1}/2}`$, the Cauchy-Schwarz inequality, the boundedness of $`B`$, $`C`$, $`A^{1/2}`$ and $`A^{1/2}`$ follows the existence of nonvanishing real constants $`K_1`$ and $`K_2`$ such that (94) is valid for all $`\xi D(A)`$. Since $`D(A)`$ is a core for both, $`A^{1/2}`$ and $`A^{\mathrm{\hspace{0.17em}1}/2}`$ (see e.g. Theorem 3.24 in chapter V.3 of ) from that inequality follows (93) and (94) for all $`\xi D(A^{1/2})`$. Note that to conclude this it is used that $`X`$ is complete and that, both, $`A^{1/2}`$ and $`A^{\mathrm{\hspace{0.17em}1}/2}`$ are closed.
Obviously, from the assumptions made follows that also $`A^{},B^{}`$ and $`C^{}`$ instead of $`A`$, $`B`$ and $`C`$, respectively, satisfy General Assumption 1 and General Assumption 4. Hence by Theorem 16 follows that
$$|T_\pm ^{}(t)|^{}e\epsilon ^{1/2}t$$
(95)
for all $`t\epsilon ^{1/2}`$, where primes indicate quantities whose definition uses one or more of the operators $`A^{},B^{}`$ and $`C^{}`$ instead of $`A,B`$ and $`C`$. In addition (93) and (94) imply $`Y=Y^{}`$ as well as the equivalence of the norms $`||`$ and $`||^{}`$. Now define the auxiliar transformation $`S_0:Y^{}Y`$ by
$$S_0\xi :=(\xi _1,\xi _2i(b/2)\xi _1)$$
(96)
for all $`\xi =(\xi _1,\xi _2)Y^{}`$. Obviously, $`S_0`$ is bijective and bounded with the bounded inverse $`S_0^1`$ given by $`S_0^1\xi :=(\xi _1,\xi _2+i(b/2)\xi _1)`$ for all $`\xi =(\xi _1,\xi _2)Y`$. In addition define $`S_\pm :[0,\mathrm{})L(Y,Y)`$ by
$$S_\pm (t):=\mathrm{exp}(ibt/2)S_0T_\pm ^{}(t)S_0^1,$$
(97)
for all $`t[0,\mathrm{})`$. Obviously, $`S_\pm `$ defines a strongly continuous semigroup with the corresponding generator
$$S_0G_\pm ^{}S_0^1\pm i\frac{b}{2}=G_\pm .$$
(98)
This implies $`S_\pm =T_\pm `$ and by (95) and (97) the existence of $`K0`$ and $`t_00`$ such that (88) is valid for all $`|t|t_0`$. Finally, from this follows by the Theorem of Hille-Yosida-Phillips that the spectrum of $`iG_+`$ is real.
###### Lemma 18
Let $`D`$ be a core for $`A`$. Further let be $`B_0:DX`$ a linear operator in $`X`$ such that for some real numbers $`a_0`$ and $`b_0`$
$$B_0\xi ^2a_0^2<\xi |A\xi >+b_0^2\xi ^2$$
(99)
for all $`\xi D`$. Then there is a uniquely determined linear extension $`\overline{B}_0:D(A^{1/2})X`$ of $`B_0`$ such that
$$\overline{B}_0\xi ^2a_0^2A^{1/2}\xi ^2+b_0^2\xi ^2$$
(100)
for all $`\xi D(A^{1/2})`$. If $`B_0`$ is in addition symmetric $`\overline{B}_0`$ is symmetric, too.
Proof: First we notice that $`D`$ is a core for $`A^{1/2}`$, too. Obviously, since $`D(A)`$ is a core for $`A^{1/2}`$ (see e.g. Theorem 3.24 in chapter V.3 of ) this follows if we can show that the closure of the restriction of $`A^{1/2}`$ to $`D`$ extends the restriction of $`A^{1/2}`$ to $`D(A)`$. To prove this let $`\xi `$ be some element of $`D(A)`$. Since $`D`$ is a core for $`A`$ there is a sequence $`\xi _0,\xi _1\mathrm{}`$ of elements of $`D`$ converging to $`\xi `$ and at the same time such that $`A\xi _0,A\xi _1\mathrm{}`$ converges to $`A\xi `$. Since $`A^{1/2}`$ has a bounded inverse it follows from this that $`A^{1/2}\xi _0,A^{1/2}\xi _1\mathrm{}`$ vonverges to $`A^{1/2}\xi `$. Since $`\xi `$ can be chosen otherwise arbitrarily it follows that the closure of the restriction of $`A^{1/2}`$ to $`D`$ extends the restriction of $`A^{1/2}`$ to $`D(A)`$ and hence that $`D`$ is a core for $`A^{1/2}`$. Hence for any $`\xi D(A^{1/2})`$ there is a a sequence $`\xi _0,\xi _1\mathrm{}`$ in $`D`$ converging to $`\xi `$ and at the same time such that $`A^{1/2}\xi _0,A^{1/2}\xi _1\mathrm{}`$ is converging to $`A^{1/2}\xi `$. Hence by (99) along with the completeness of $`X`$ follows the convergence of the sequence $`B_0\xi _1,B_0\xi _2\mathrm{}`$ to some element $`\overline{B}\xi `$ of $`X`$ and
$$\overline{B}\xi ^2a_0^2A^{1/2}\xi ^2+b_0^2\xi ^2.$$
(101)
Moreover if $`\xi _0^{},\xi _1^{}\mathrm{}`$ is another sequence having the same properties as $`\xi _0,\xi _1\mathrm{}`$ by (99) follows that
$$\overline{B}\xi =\underset{n\mathrm{}}{lim}B_0\xi _n=\underset{n\mathrm{}}{lim}B_0\xi _n^{}.$$
(102)
From this it easily seen that by defining
$$\overline{B}:=(D(A^{1/2})X,\xi \overline{B}\xi )$$
(103)
there is also given a linear map. Hence the existence of a linear extension of $`B_0`$ satisfying (100) is shown. Moreover from the definition it is obvious that $`\overline{B}`$ is symmetric if $`B_0`$ is in addition symmetric. If on the other hand $`\overline{B}_0`$ is a linear extension of $`B_0`$ satisfying (100) and $`\xi `$ and $`\xi _1,\xi _2`$ are as above from (100) follows that
$$\overline{B}_0\xi =\underset{n\mathrm{}}{lim}B_0\xi _n.$$
(104)
Finally, since $`\xi `$ can be chosen otherwise arbitrarily from this follows $`\widehat{B}_0=\widehat{B}`$.
## 3 Discussion and results
This paper provides a rigorous framework for the description of linearized adiabatic lagrangian perturbations and stability of differentially rotating newtonian stars using semigroup theory. Problems of a previous framework by Dyson and Schutz are overcome and a basis for a rigorous analysis of the stability of such stars is provided. The spectrum of the oscillations is shown to coincide with the spectrum of an operator polynomial whose coefficients can be read off from the equation governing the oscillations about the equilibrium configuration. Moreover, for the first time sufficient criteria for stability are given in form of inequalities for the coefficients of that polynomial. These show that a negative canonical energy of the star does not necessarily indicate instability.
It is still unclear whether these criteria are strong enough to prove stability for realistic stars. On the other hand the second criterium has been sucessfully applied in the (on first sight seemingly unrelated case of the) stability discussion of the Kerr metric where the master equation governing perturbations is of the form (1), too. Another similarity of that case to the cases considered here is the fact that the corresponding operators $`C^{}`$ and $`B^{\mathrm{\hspace{0.17em}2}}`$ there are such that $`C^{}(1/4)B^{\mathrm{\hspace{0.17em}2}}`$ is positive whereas here this combination is semibounded as has been shown in DS.
Also the determination of the spectrum of the operator polynomial $`C^{}\lambda B^{}+\lambda ^2,\lambda `$ for some special case would be very useful. It is likely that this cannot be done for a physically relevant case. But it is also likely that the outcome to qualitative questions like
* Does one have uniform stability in $`m`$?
* Does a continuous part occur in the oscillation spectrum?
only depends on structural properties of the operators $`C^{}`$ and $`B^{}`$. So from $`C^{}`$ probably only the highest order derivatives are relevant and details of the equation of state should be unimportant. From this point of view even the highly idealized case of a spherical background model with a truncated $`C^{}`$ along with a non constant velocity field $`v`$ would be interesting to consider.
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# Untitled Document
Space-Time Structure as Hidden Variable
Bart Jongejan
Cæciliavej 31, 2500 Valby, Denmark
Electronic address: bart@cst.ku.dk
EPR correlations exist and can be observed independently of any a priori given frame of reference. We can even construct a frame of reference that is based on these correlations. This observation-based frame of reference is equivalent to the customary a priori given frame of reference of the laboratory when describing real EPR experiments.
J.S. Bell has argued that local hidden parameter theories that reproduce the predictions of Quantum Mechanics cannot exist, but the counterfactual reasoning leading to Bell's conclusion is physically meaningless if the frame of reference that is based on EPR-correlations is accepted as the backdrop for EPR-type experiments.
The refutal of Bell's proof opens up for the construction of a viable hidden parameter theory. A model of a spin $`\mathrm{}`$/2 particle in terms of a non-flat metric of space-time is shown to be able to reproduce the predictions of quantum mechanics in the Bohm-Aharonov version of the EPR experiment, without introducing non-locality.
PACS numbers: 03.65.Bz, 02.40.Ky, 04.20.Gz
I.HIDDEN VARIABLES
A.The Einstein, Podolsky and Rosen Gedanken experiment
Albert Einstein was convinced that quantum mechanics is an incomplete theory, which was a position opposite to that of Niels Bohr. Their discussion culminated in a paper by Einstein, Podolsky and Rosen (EPR) , in which they showed how one can measure two non-commuting physical quantities to any degree of accuracy. Bohr's reply was soon to follow. The issue is still debated.
The set-up of the EPR experiment consists of two observation posts doing measurements of momentum or position on the flown-apart members of particle pairs that have been carefully prepared to have no net momentum relative to the laboratory frame.
The particles in a pair can be regarded as exact copies of each other, apart from being mirrored: if the same measurement is performed at both observation posts, then the outcomes are each other’s exact opposites. Given a quantity to be determined for both particles, we can suffice with doing only one of the measurements.
EPR’s idea was that two non-commuting quantities, such as momentum and position, can be determined by measuring one of the quantities directly and by deriving the value of the other quantity from the outcome of the measurement of that quantity on the other particle in the pair.
Bohm and Aharonov devised a version of the EPR Gedanken experiment that has been the focus of much theoretical and experimental work. In their experiment, the measurements are done on pairs of spin $`\mathrm{}`$/2 particles that are prepared in the singlet state, which is a quantum state that does not hold any information about the directions of the spins of the individual particles.
In the Bohm-Aharonov experiment the particles fly apart toward two widely separated observation posts, where they traverse the gaps of Stern-Gerlach magnets. During such a traversal, due to a coupling between the particle’s intrinsic spin and the longitudinal gradient of the magnetic field, the particle’s path is bent either away from or toward the pole where the magnetic field is strongest. The particle finally hits one of two detectors, depending on which route it took. One of the detectors only records particles that had spin up $`()`$, while the other only records those with spin down $`()`$, “up” and “down” being defined relative to the Stern-Gerlach magnet. The detectors are fixed to their respective Stern-Gerlach magnets, so that the directions that are “up” and “down” rotate together with the freely orientable Stern-Gerlach magnets.
B.Can counterfactual considerations complete the description of physical reality?
Employing propositions of the type that EPR used to show that quantum mechanics can not be complete, J.S. Bell showed that any theory that reproduces the predictions made by quantum mechanics and yet is more complete than quantum mechanics necessarily postulates instantaneous action at a distance. In other words, the kind of theories that Einstein envisaged as successors to quantum mechanics would be difficult to reconcile with relativity theory, which champions locality and does not allow any signal to travel faster than light.
Bell’s proof is dependent on counterfactual propositions. A counterfactual proposition assigns a determinate value to a quantity that could have been directly observable, but is not, typically because another, incommensurable quantity is measured. The experimental basis - if you can call it that - for this assignment is the measurement of the same quantity on the far away twin particle. The reasoning is that the measurement on the twin is as good as the measurement on the particle itself, because the inner states of the particles must be fully correlated in order to preserve the isotropy of the quantum state of the pair.
To give teeth to such a counterfactual proposition, not only do we have to ascribe reality to the result of the measurement, but also to the angle between the counterfactual orientation of the instrument and the (factual or counterfactual) orientation of the other instrument.
Let us hold EPR’s and Bell’s counterfactual reasoning against the background of Riemannian geometry, or, more specifically, general relativity. Only geometric relations, such as angles and distances, between things that are local to each other in space and in time bear physical meaning according to the general relativity theory. Geometric relations over long spatio-temporal distances, such as the angle between the directions of observation in the Bohm-Aharonov experiment, are a different matter.
If it is not assumed that space-time is flat, then only an operational definition can lend physical meaning to these relations. In general, different operational definitions can give physical meaning to the same relation, but in theory the methods need not agree on the outcomes of the measurements. This is clearly exemplified by the multitude of theory laden operational definitions for distances on cosmological scales.
What then is the angel between a factual set-up of a Stern-Gerlach magnet at one of the two observation posts and a counterfactual set-up of the Stern-Gerlach magnet at the other observation post? Of course, because one of the two set-ups is not effectuated, there is no direct means of measuring this angle. We can only base our answer on interpolation, together with a smoothness assumption that coordinates the interpolated but counterfactual observation with real ones. Normally, interpolation depends on a smoothness assumption that is innocuous, because the missing, interpolated values may some day be replaced by outcomes of real experiments. In such cases, interpolation is a falsifiable theory and therefore acceptable. In Bell’s proof, the interpolation is not falsifiable, because the measuring apparatus is sitting itself in the way. That weakens Bell’s conclusion. We can only accept Bell’s proof if we assume that the space time backdrop is smooth and constant enough to allow us to interpolate between measurements, but this assumption excludes from consideration any theory that denies that space-time is like that.
C.The Bohm-Aharonov experiment without flat-space preconception.
Bell's proof hinges on the postulate that space time is flat, but this postulate may be false. This is the main theme of this paper and we will dwell on it a little more, because understanding the epistemological restrictions that relativity imposes is essential for appreciating the approach towards hidden variables that is presented later.
First think of the Bohm-Aharonov experiment as a set-up consisting of two observation posts connected by the floor of a laboratory or something else that we may regard as rigid. Each observation post consists of a Stern-Gerlach magnet that is freely orientable in its mounting into any of a large number of directions, or lines of observation. Each such orientation is identifiable, for example by reading off the color of a mark on the mounting that a pointer aligned with the axis of the magnet is pointing at. There may be many differently colored marks, each identifying a unique orientation of the magnet. If one wishes so one can define the mountings of the instrument to include far away stars, which then can be used as the marks for that post. We also have calibrated scales on the mounting so that we have the option to read off the angular coordinates of the direction vector of the Stern-Gerlach magnet. A post's electron detectors are fixed to the Stern-Gerlach magnet inside that post. The rigid connection between the two posts (or rather: between the mountings), together with conventional means of doing geodesy (measuring rods, light beams, gyroscopes) provides us with a reference frame in which both measuring instruments have definite positions and orientations. Even counterfactual orientations can be tracked, because the rigid frame "fixes" all thinkable orientations. This is the normal, "robust" experimental context of Bell's proof.
Now remove all unnecessary equipment: the measuring rods, light beams, gyroscopes, scales and even the laboratory floor. Would that make any difference? We did not do away the mountings of the instruments and are therefore still able to identify each orientation by reading off the color of the mark that a Stern-Gerlach magnet is pointing at. Can we reconstruct the experiment with this basic equipment?
From each pair of measurements we obtain a data-triplet: the color of the mark that the left magnet was pointing at, idem for the right magnet and finally the outcome of the detectors, which arbitrarily may be defined to be "S" (for "same") if both "up" or both "down" detectors were hit and "N" (for "not same") if one "up" and one "down" detector was triggered. Our logbook will have just three columns, the experiment has only three degrees of freedom.
After doing a long series of such experiments, with the magnets having been oriented in all possible combinations of directions many times, the log of outcomes will enable us to assign statistical probabilities for measuring the same spin component to each pair of orientations of the Stern-Gerlach magnets. We might for example observe that "yellow" left and "blue" right have a 73% chance of resulting in the value "S". There is a conceptual difference with the full blown experimental set-up, though: we have no prior knowledge of the geometric relations between the Stern-Gerlach magnets; we can only pairwise statistically relate marks on the mountings with each other. We do not even know the angles between two different orientations of the same magnet!
Now we will try to organize the three columns of obtained data by mapping the colored marks that identify directions of observations onto points on a sphere in such a way that exactly one point is assigned to each mark. This mapping must map the marks on both surroundings, however far separated from each other, onto a single sphere that does not exist physically, but only as a mathematical tool. Underlying the mapping is the working hypothesis that there is a functional relation between the statistical probability to measure the same component and the angle, as measured on this sphere, between the orientations of the Stern-Gerlach magnets. We can start to assume a linear dependency: a 100% probability and a 0 % probability correspond to angles of $`180^{}`$ and of $`0^{}`$ between the magnets, a 10% probability would correspond to $`18^{}`$, and so on. Such a relation would be appropriate if the spinning particles were macroscopic objects with observationally well defined "north" and "south" hemispheres. However, if we apply this relation to our data then the obtained angles force us to map the same mark onto several points, which is not what we wanted.
Eventually, we would find that
$$P\left(S\right)=probabilitytomeasuresamevalue=$$
$$=\mathrm{sin}^2\left(\frac{anglebetweenorientations}{2}\right)$$
$`(1a)`$
establishes a 1-1 mapping, as does its mirroring twin
$$P\left(S\right)=probabilitytomeasuresamevalue=$$
$$=\mathrm{cos}^2\left(\frac{anglebetweenorientations}{2}\right).$$
$`(1b)`$
Now, not only would we have learned how to define angles between the two measuring instruments when directed towards two given colored marks, we would also have a consistent way to figure out the angles between the colored marks at one and the same post. That means that we would have regained the spherical geometry relations that we did not assume as given a priori. That, in turn, would enable us to deliver exactly the same kind of experimental report as someone who had a rigid frame to connect the measuring posts and rods, light beams, gyroscopes and scales to measure the geometry. We could also verify the contingent fact that this statistical way of determining geometric relations is perfectly consistent with other, more conventional means, such as with the help of rods, light beams, gyroscopes and scales. But of course, as little physical sense it makes to ascribe temperature and pressure (or the mean kinetic energy and momentum per particle) to a single particle in a gas, as little sense would it make to ascribe the statistically defined angle to a pair of measurements that contributed to the very determination of the angle. A fortiori, we can not draw firm conclusions from an argument that hinges on a counterfactual set-up controlled by the statistically defined angle. There is no observational basis for assigning a value to such an angle, because the value, as defined operationally in the above way, is of statistical character and therefore not applicable to any pair of orientations in any specific pair of measurements, but only in the long run of many measurement event pairs. In addition, the assignment of a value of the angle, operationally defined as above, to pairs of orientations that are not both effectuated, can only be based on an arbitrary convention and therefore renders any argument that is based on such angles unconvincing.
If the presented way of constructing a frame of reference is so feeble that it does not allow us to assign values to angles between counterfactual setups, why then should we not stick with the conventional means using rods, gyroscopes and so on? The reason is that our method is not any more feeble than conventional means! Unless conventional methods mysteriously gain strength somewhere in the transition from the quantum to the classical regime, any angle that can be measured by conventional means can also be measured using a sufficiently large ensemble of spin component observations on an equally large number of pairs of spin $`\mathrm{}`$/2 particles, to any desired accuracy. However, the proposed method, which is obviously based on quantum phenomena, has the additional advantage that it very clearly delineates the domain of applicability of the method; a domain of applicability that can not be surpassed by conventional methods, unless the aforementioned mysterious powers opened a back door for the conventional methods to do measurements of angles that are out of reach for the proposed method.
Whereas proofs like Bell's are hard pressed because of the lack of an observational basis for the assumed frame of reference, a realist point of view is not obstructed in the same way. It is sensible to imagine a counterfactual set-up of a measuring apparatus that is oriented towards a particular colored mark and with a definite value of the spin component in that direction, but we must keep in mind that the whereabouts of the instrument's orientation relative to other (factual or counterfactual) orientations is unknown. Underlying Bell's and similar proofs is a concept of realism that is far more encompassing than necessary. Things that derive the status of being part of reality by force of real observations, such as the angle between directions of observation and even systems of coordinates in general, can not be idealized to an existence detached from these observations without introducing trouble in some corners of our theoretical picture of the world.
II.MATHEMATICAL CONSTRAINTS ON HIDDEN VARIABLE THEORIES OF SPIN $`\mathrm{}`$/2 PARTICLE
A.What makes an aspirant HV theory?
We have seen that the Bohm-Aharonov experiment has just three relevant degrees of freedom: color of left mark, color of right mark and the combination of the outcomes. We will now discuss hidden variable theories that also exhibit three degrees of freedom and hope to find one that can be made to correspond to the Bohm-Aharonov experiment and that explains the statistical correlations found in the Bohm-Aharonov experiment (which are assumed to be accurately predicted by quantum mechanics).
The contemplated hidden variable theories all have one aspect in common: not only the Stern-Gerlach magnets have definite orientations, also the particle itself has one, which is the axis of rotational symmetry or "spin axis". The three variables that specify any configuration of the three directions are (a) the angle between the left measuring apparatus and the spin axis, (b) the angle between the right measuring apparatus and the spin axis and (c) the angle between the left measuring apparatus and the right measuring apparatus. None of these variables are precisely measurable, but each corresponds to observed data: if the left "up" detector is triggered, then the angle between the left measuring apparatus and the spin axis is less than $`90^{}`$. If the "down" detector is triggered, the angle is somewhere between $`90^{}`$ and $`180^{}`$. Likewise for the right detector. The angle between the measuring apparatuses is taken to be the angle that was statistically derived from the outcomes of a long series of Bohm-Aharonov measurements, using Eq. (1b).
The main point made by Bell was that no local hidden variable theory is able to reproduce the predictions of quantum mechanics. The requirement that our aspirant hidden variable theories are local puts three constraints on the statistical distribution of the angles between the measuring apparatuses and between each measuring apparatus and the particle's spin axis. These will be discussed now.
Constraint 1. The orientations of the measuring instruments are unrelated.
Each measuring apparatus is oriented in a way that does not depend on the orientation of the other instrument, not even statistically. If the movements of the measuring apparatuses A and B were restricted to a plane, then this constraint would translate to a uniform distribution of angles $`\mathrm{}AB`$ in the range $`0\mathrm{}AB\pi `$. We do, however, assume that the instruments are freely orientable in space. In that case, the probability that the angle between A and B is $`\mathrm{}AB`$ is proportional to $`\mathrm{sin}\mathrm{}AB`$.
Whereas the angle is non-uniformly distributed, its cosine is not. So we require a uniform distribution $`\rho \left(Z_{AB}\right)`$ of the inner product $`Z_{AB}=\mathrm{cos}\mathrm{}AB=𝚊.𝚋`$ in the range $`1Z_{AB}1`$. The minus sign is arbitrarily introduced to compensate for the circumstance that the particles' spins are anti-parallel.
Normalization requires $`_1^1\rho \left(Z_{AB}\right)𝑑Z_{AB}=1`$, so that
$$\rho \left(Z_{AB}\right)=1/_2$$
$`(2)`$
It is worthwhile to indicate which role the set-up plays. We ask that during each measurement the measuring apparatus points at a randomly chosen point of its own surroundings. This does not automatically imply an isotropic distribution of the orientations with respect to each other: one could imagine that each observation post's orientations, taken separately, would survive a "randomness test", but that the orientations were not distributed isotropically with respect to each other. That situation could arise if the observers used the same sequence of random numbers to prepare the instruments for each pair of measurements. We assume that such correlation does not occur, because that seems to be the only assumption that is compatible with the principles of locality, causality and free will.
Constraint 2. (Locality condition.) The orientation of one magnet does not influence the result obtained with the other.
Suppose that someone came up with a HV theory of a spin $`\mathrm{}`$/2 particle. In order to test the claim that it reproduces the predictions of QM in Bohm-Aharonov experiments, we had to subject the theory to a Gedanken experiment in which a great number of spin-component measurements were randomly chosen. How would the randomly chosen orientations of a measuring instrument be distributed with respect to the particle's axis of rotational symmetry? As we have no means of observing this distribution, we postulate one. In the absence of any reason to assume a non-isotropic distribution, we assume the isotropic distribution. Call the inner product of the orientation of the instrument and the direction of the particle (denoted by unit vectors) $`Z_A`$ and $`Z_B`$ (for instrument A and instrument B). Require that $`\rho \left(Z_A\right)=\rho \left(Z_B\right)=1/_2`$.
The sign of $`Z_A`$ determines the outcome of the measurement made with the magnet at observation post A. The locality condition says that $`Z_A`$ is independent of the orientation of instrument B. The angle between the instruments can be taken to represent this orientation, in which case we locality condition translates to
$$\rho (Z_A,Z_{AB})=\rho \left(Z_A\right)\rho \left(Z_{AB}\right)=1/4.$$
$`(3a)`$
Alternatively, we can take the (hidden) angle between instrument B and the spin axis as representing the orientation, so we also require that
$$\rho (Z_A,Z_B)=\rho \left(Z_A\right)\rho \left(Z_B\right)=1/4.$$
$`(3b)`$
As we are used to specify orientations with two mutually independent angles, it is tempting to require that $`Z_A`$ is independent of both of $`Z_{AB}`$ and $`Z_B`$:
$$\rho (Z_A,Z_B,Z_{AB})=\rho \left(Z_A\right)\rho \left(Z_B\right)\rho \left(Z_{AB}\right)=1/8.$$
$`(4)`$
However, below we will see that this conflicts with the next constraint and even with the statistics of classical spinning particles.
Constraint 3. The theory reproduces the predictions of quantum mechanics
The correlation between the outcomes of measurements on flown-apart particles with counter parallel spin axis conforms exactly to the predictions of QM (see also Eq. (1b)):
$$P\left(S\right)P()+P()=\mathrm{sin}^2\frac{\mathrm{}AB}{2}$$
$$=\frac{1+Z_{AB}}{2}$$
$`(5a)`$
$$P\left(N\right)P()+P()=\mathrm{cos}^2\frac{\mathrm{}AB}{2}$$
$$=\frac{1Z_{AB}}{2}.$$
$`(5b)`$
B.The joint distribution $`\rho (𝚉_𝙰,𝚉_𝙱,𝚉_{\mathrm{𝙰𝙱}})`$ that explains the quantum mechanical predictions.
We assume that in a HV theory, a full specification of a measurement of two spin components in the Bohm-Aharonov experiment requires three angles, or there cosines. We must now investigate whether there are distributions $`\rho (Z_A,Z_B,Z_{AB})`$ of these three quantities that fulfill all three constraints. For example, the first constraint requires that
$$\rho \left(Z_{AB}\right)=_1^1_1^1\rho (Z_A,Z_B,Z_{AB})𝑑Z_A𝑑Z_B=1/2,$$
$`(6)`$
and second constraint, the locality condition, requires that
$$\rho (Z_A,Z_{AB})=_1^1\rho (Z_A,Z_B,Z_{AB})𝑑Z_B=1/4.$$
$`(7)`$
Finally, according to quantum mechanics we must find that
$$P\left(\right)=\rho \left(Z_{AB}\right)^1_0^1_0^1\rho (Z_A,Z_B,Z_{AB})𝑑Z_A𝑑Z_B$$
$$=\frac{1+Z_{AB}}{4}.$$
$`(8)`$
As it did not seem a trivial task to solve the set of equations constituting the three constraints, a computer aided approach was chosen. It was not difficult to find a distribution that fulfills constraints 1 an 2, Eq. (4) is such a distribution. Then, repeatedly applying an algorithm that transforms the distribution to a new distribution that also fulfills constraints 1 and 2, distributions were found that fulfill the third constraint as well. The algorithm is as follows: Choose two values for each of $`Z_A,Z_B`$ and $`Z_{AB}`$. These values are the coordinates of eight cells, the probability density of which we are going to redistribute. Choose an amount $`\mathrm{}\rho `$ and add this amount to four cells spanning a tetrahedron and subtract the same amount from the remaining cells. By the right choice of $`\mathrm{}\rho `$ we can empty at least one of the eight cells. See Fig. 1. Some of the results of this discrete approximation can be seen in Figs. 2-4.
The distribution in Eq. (4) fully acknowledges the freedom of the experimenters to vary the angle between the instruments ($`Z_{AB}`$) and it guarantees the isotropic distribution of the axis of the model relative to the lines of observation ($`Z_A`$ and $`Z_B`$). Yet this distribution is not realistic at all, because it does not restrict the angles between the instruments ($`\mathrm{arccos}Z_{AB}`$) and the angles between the measurement instruments and the axis of the model ($`\mathrm{arccos}Z_A`$ and $`\mathrm{arccos}Z_B`$). These three angles can not be completely independent: for example can the angle between the measuring instruments not exceed the sum of the angles between the instruments and the axis of the model.
If two of the angles already are fixed to any values between $`0`$ and $`\pi `$ (any two of $`\{\mathrm{arccos}Z_{AB}`$,$`\mathrm{arccos}Z_A`$,$`\mathrm{arccos}Z_B\}`$, call them $`\alpha _1`$ and $`\alpha _2`$), then we have the following constraint on the third angle $`\alpha _3`$:
$$\left|\alpha _1\alpha _2\right|\alpha _3\alpha _1+\alpha _2.$$
$`(9)`$
That means that whereas any two of the three angles are independent of each other, there exists a mutual dependency between the three angles.
The uniform distribution of three vectors $`𝚊,𝚋`$ and $`𝚌`$ over all directions illustrates this dependency. In an Euclidean frame of reference, the joint density of the three inner products $`s=Z_A=𝚊.𝚌,t=Z_B=𝚋.𝚌,`$$`u=Z_{AB}=𝚊.𝚋`$ is
$$\rho (s,t,u)=\left(8\pi \sqrt{1+2stus^2t^2u^2}\right)^1$$
$$\left(1+2stus^2t^2u^2>0\right)$$
$$=0(1+2stus^2t^2u^20).$$
$`(10)`$
Figure 3 illustrates the mutual dependency between the three directions and the remarkable discontinuity of the density at the border between possible and impossible configurations,
The next step is the fulfillment of constraint 3. Using the transformation algorithm again, we eventually approach a distribution that reproduces the predictions of quantum mechanics (Eq. 5a/b). There are many distributions that come very close, but they are all characterized by regions of almost emptiness and steep climbs to high values of probability density. The following distribution, which is the limiting distribution as the number of subdivisions in the discrete approximation goes to infinity, and which uses Dirac delta functions, fulfills all three constraints (see Fig. 4 for a discrete approximation):
$$\rho (Z_A,Z_B,Z_{AB})=1/8[\delta (Z_A+Z_B+Z_{AB}+1)$$
$$+\delta \left(Z_AZ_BZ_{AB}+1\right)$$
$$+\delta \left(Z_A+Z_BZ_{AB}1\right)$$
$$+\delta (Z_AZ_B+Z_{AB}1)]$$
$`(11)`$
This distribution indicates that one continuous degree of freedom is replaced by a discrete one: the density is only non-zero on the surface of a tetrahedron spanned by four of the eight corners of the configuration cube.
A classical configuration of three independent vectors is specified by three numbers, which are the lengths of the sides of a triangle on the unit sphere. They can not live within less than the two dimensions of this sphere. On the other hand, configurations that are compatible with QM require only two numbers and a sign, the third number being a function of either the sum or the difference of the other two, depending on the sign.
Perhaps somewhat unexpectedly, our search for a probability distribution for the angle between the instruments and the angles between each instrument and the (hidden) spin axis did not merely result in a non-classical distribution, but also in a qualitative characterization of any hidden variable theory with hopes to fulfill all three constraints: the theory must endow a model of a spin $`\mathrm{}`$/2 particle with a degeneracy that replaces one continuous degree of freedom with a two-valued one. We will now look at a theory that accomplishes this feat.
III.EXAMPLE: A HIDDEN VARIABLE MODEL BASED ON NON-FLAT SPACE-TIME STRUCTURE
A.Overview
In the foregoing, a weakness in Bell's counterfactual reasoning was exposed by peeling away unwarranted and mostly silent assumptions that are underlying his proof, until we were left with the nitty gritty observational data. Then we postulated, in spite of Bell's conclusion, that the direction of a particle's spin is really existing, although hidden from observation. We found that the QM-compatible configurations of three vectors, denoting the orientations of the measuring instruments and a candidate model's hidden variable, seem to live within one dimension less than classically expected.
The next step is the construction of a model with a characteristic that we rightly can call the spin direction of the model. We postulate that the stuff that the particle is made of is space time itself, or more specifically, the structure of space time. We will check that the full specification of a geodetic path herein - tentatively representing the world line of a measuring apparatus during a measurement - exhibits the same lack of one degree of freedom.
Such a postulated space-time structure is in part based on guesswork, in part on esthetic rules, such as simplicity and symmetry. The real structure of space time is perhaps unknowable, but we may hit upon a theory of the structure of space time that survives the observations that we perform to test it.
The presented model tries only to explain a very limited set of phenomena, namely the correlation between two spin component measurements. We have not tried very hard to incorporate and explain other phenomena. Thus, a simple thing like the spatial distance between the spinning particles is not expressed very well by the proposed model, nor their relative movements. In fact, the model explains two widely different phenomena without making a distinction, which shows that in any case the differences between these phenomena are not expressed in the model. The first phenomenon is the correlation between the measurements of spin components on two different particles that together form a system in the singlet state. The second phenomenon, that is explained equally well, is the passage of a single spinning particle first through one, then through a second Stern-Gerlach magnet at some distance from the first, that is inclined with respect to the first.
B.Metric and geodesic equations
Consider a metric $`g_{ik}`$
$$ds^2=\mathrm{cos}^2\vartheta dt^2dr^2r^2d\vartheta ^2r^2\mathrm{sin}^2\vartheta d\phi ^2$$
$$+\mathrm{\hspace{0.17em}2}r\mathrm{sin}^2\vartheta d\phi dt.$$
$`(12)`$
To get a feeling of how it would be like to be in a space time with this metric, imagine an infinite set of concentric spheres. Suppose that you momentarily were attached to one of these spheres. You would have the feeling that this sphere was rotating around the $`\vartheta =0`$ axis with angular velocity $`\frac{1}{r}=\frac{1}{radiusofsphere}`$. In other words, you would feel an acceleration away from the $`\vartheta =0`$ axis. If you were to move in the direction opposite to this intrinsic rotation with angular velocity $`1/r`$ the acceleration would disappear. Although everywhere the experienced acceleration could be explained as the effect of an angular velocity that gradually decreased as $`\frac{1}{r}`$, the spheres in this space time are fixed to each other forever. It is the curvature of space-time that gives the experience of being accelerated, just like the curvature of space-time caused by the Earth's mass lets us feel the force of gravitation: an acceleration without movement, as opposed to the acceleration experienced in a rocket.
The rotational effects in this space-time are due to the last term in the metric ground form, $`+2r\mathrm{sin}^2\vartheta d\phi d\vartheta `$. This term specifies the model's spin direction relative to the chosen frame of reference; by changing the sign of this term or by rotating $`180^{}`$ along an axis perpendicular to the $`\vartheta =0`$ axis we obtain a model with opposite spin direction.
In order to investigate how free test-particles move in this space-time we have to solve the equations of motion
$$\frac{d^2x_i}{ds^2}=\mathrm{\Gamma }_{jk}^i\frac{dx_j}{ds}\frac{dx_k}{ds}.$$
$`(13)`$
The non-zero coefficients of the affine connection $`\mathrm{\Gamma }_{jk}^i`$, which are symmetric in the lower indices, are:
| $`\mathrm{\Gamma }_{tr}^t=\frac{\mathrm{sin}^2\vartheta }{2r}`$ | | $`\mathrm{\Gamma }_{r\phi }^t=\frac{\mathrm{sin}^2\vartheta }{2}`$ | |
| --- | --- | --- | --- |
| $`\mathrm{\Gamma }_{t\phi }^r=\frac{\mathrm{sin}^2\vartheta }{2}`$ | | $`\mathrm{\Gamma }_{\vartheta \vartheta }^r=r`$ | |
| $`\mathrm{\Gamma }_{\phi \phi }^r=r\mathrm{sin}^2\vartheta `$ | | $`\mathrm{\Gamma }_{t\phi }^\vartheta =\frac{\mathrm{sin}\vartheta \mathrm{cos}\vartheta }{r}`$ | |
| $`\mathrm{\Gamma }_{r\vartheta }^\vartheta =\frac{1}{r}`$ | | $`\mathrm{\Gamma }_{\phi \phi }^\vartheta =\mathrm{sin}\vartheta \mathrm{cos}\vartheta `$ | |
| $`\mathrm{\Gamma }_{tt}^\vartheta =\frac{\mathrm{sin}\vartheta \mathrm{cos}\vartheta }{r^2}`$ | | $`\mathrm{\Gamma }_{tr}^\phi =\frac{\mathrm{cos}^2\vartheta }{2r^2}`$ | |
| $`\mathrm{\Gamma }_{r\phi }^\phi =\frac{1+\mathrm{cos}^2\vartheta }{2r}`$ | | $`\mathrm{\Gamma }_{\vartheta \phi }^\phi =\frac{\mathrm{cos}\vartheta }{\mathrm{sin}\vartheta }`$. | |
$`\left(14\right)`$
With the shorthand notation $`U^r=r^{}=\frac{dr}{ds}`$, $`U^r^{}=r^{\prime \prime }=\frac{d^2r}{ds^2}`$ etc. and the above connection $`\mathrm{\Gamma }_{jk}^i`$ the geodesic equations become
$$U_{}^{t}{}_{}{}^{}=\frac{U^r\mathrm{sin}\vartheta \left[\mathrm{sin}\vartheta (U^trU^\phi )\right]}{r}$$
$`(15a)`$
$$U_{}^{r}{}_{}{}^{}=\frac{(rU^\vartheta )^2(rU^\phi \mathrm{sin}\vartheta )[\mathrm{sin}\vartheta (U^trU^\phi )]}{r}$$
$`(15b)`$
$$U_{}^{\vartheta }{}_{}{}^{}=\frac{\mathrm{\hspace{0.17em}2}U^r(rU^\vartheta )+\frac{\mathrm{cos}\vartheta }{\mathrm{sin}\vartheta }[\mathrm{sin}\vartheta (U^trU^\phi )]^2}{r^2}$$
$`(15c)`$
$$U_{}^{\phi }{}_{}{}^{}=\frac{1}{r^2\mathrm{sin}\vartheta }\times \{U^r(rU^\phi \mathrm{sin}\vartheta )$$
$$+U^r\mathrm{cos}^2\vartheta \left[\mathrm{sin}\vartheta (U^trU^\phi )\right]$$
$$+\mathrm{\hspace{0.17em}2}\frac{\mathrm{cos}\vartheta }{\mathrm{sin}\vartheta }(rU^\vartheta )\left[\mathrm{sin}\vartheta (U^trU^\phi )\right]\}.$$
$`(15d)`$
The geodesic equations have the following solutions:
$$U^t=P+\frac{X}{r}$$
$`(16a)`$
$$U^\phi =\frac{P\frac{X\mathrm{cotg}^2\vartheta }{r}}{r}$$
$`(16b)`$
$$\left(U^\vartheta \right)^2=\frac{A\frac{X^2}{\mathrm{sin}^2\vartheta }}{r^4}$$
$`(16c)`$
$$\left(U^r\right)^2=\frac{AX^2}{r^2}+\frac{2PX}{r}+P^2W,$$
$`(16d)`$
where $`P,X,A`$ and $`W`$ are constants. We can also write
$$P=\mathrm{cos}^2\vartheta U^t+r\mathrm{sin}^2\vartheta U^\phi $$
$`(17a)`$
$$X=r\mathrm{sin}^2\vartheta \left(U^trU^\phi \right)$$
$`(17b)`$
$$A=\left(r^2U^\vartheta \right)^2+\left[r\mathrm{sin}\vartheta \left(U^trU^\phi \right)\right]^2$$
$`(17c)`$
$$W=\left(\mathrm{cos}\vartheta U^t\right)^2\left(U^r\right)^2\left(rU^\vartheta \right)^2\left(r\mathrm{sin}\vartheta U^\phi \right)^2$$
$$+\mathrm{\hspace{0.17em}2}r\mathrm{sin}^2\vartheta U^tU^\phi $$
$`(17d)`$
$`W`$ is simply $`g_{ij}U^iU^j`$, the length of the vector $`𝚄`$ squared. Time-like geodesics have $`W>0`$, space-like geodesics have $`W<0`$ and light-like geodesics have $`W=0`$. We will restrict $`W`$ to the values -1, 0 and 1. This restriction removes an arbitrary scaling factor.
C.Comparison with paths in central force fields
If we only look at time-like geodesics in the direction from past to future (the paths that test particles follow, $`W=1,U^t>0`$), then only three numbers $`P,X`$ and $`A`$ are needed to fully specify a geodesic. What are the consequences of this paucity of orbit-fixing numbers? Let us compare a geodesic in a central gravitational field and a geodesic in this model-space time. Four constants are needed to specify the orbit of a freely falling test particle A with respect to a massive body, such as the orbit of a planet around a star. The shape of the orbit or eccentricity \- whether the orbit is a circle, an ellipse, a parabola or a hyperbola - provides one number. The size of the orbit - e.g. the distance of closest approach to the central massive body \- provides another number. The orientation of the orbital plane, defined as the unit vector normal to the orbital plane, requires the specification of two angular coordinates. That adds two more degrees of freedom and brings the total number of constants to four.
For the time-like geodesics in our model-space-time the situation is different. If we assume that two numbers are needed to specify "shape" and "size", then only one number is left to specify the "orbital plane". The quotes indicate that we cannot be sure that it makes sense to talk about shape, size and orbital plane. We will later have to look at that.
If one number fixes an orbital plane then obviously the orbital orientation has only one degree of freedom. Given one orbital plane, then another orbital plane can be specified with reference to the given orbital plane by providing the difference of the "orbital plane numbers". In fact, there is an ambivalence in such a specification because of the sign of the difference. This sign can not be specified without breaking the symmetry between the two planes: we have to assume that either plane can play the role of reference plane and the expressions should not depend on this choice in any arbitrary way.
We have seen the same "directional degeneracy" before: the distribution of values $`Z_A,Z_B,Z_{AB}`$ that reproduces quantum mechanics is such that given one $`Z`$-value, the other $`Z`$-values can be specified with reference to this single $`Z`$-value. For example, if $`Z_A`$ is given then $`Z_B`$ is specified, up to a bivalent choice, by giving just one more number: $`Z_{AB}`$ (see 7a-d).
D.How the constants define shape and orientation of the geodesic
We will now look more closely at the constants of the motion and see whether it really is the case that there is only one number to specify the orbital plane.
$`A,P,X`$ and $`W`$ are real numbers that only to some degree can be chosen freely. The expression for $`\left(U^\vartheta \right)^2`$ indicates that only values $`\left|X\right|\sqrt{A}`$ can lead to geodesics. It also shows that such geodesics are restricted to points where $`\mathrm{sin}\vartheta \frac{\left|X\right|}{\sqrt{A}}`$. Geodesics with values of $`\left|X\right|`$ close to $`\sqrt{A}`$ are close to the equatorial plane $`\vartheta =\frac{\pi }{2}`$, while a value of $`\frac{\left|X\right|}{\sqrt{A}}`$ close to zero allows the geodesic to approach the poles very closely. Therefore we define
$$S=\frac{X}{\sqrt{A}}\left(1S1\right)$$
$${}_{}{}^{\prime \prime }tiltoforbitalplane^{\prime \prime }.$$
$`(18)`$
If $`S=1`$ or $`S=1`$ then $`\mathrm{sin}\vartheta =1`$; such geodesics are equatorial. All other geodesics are wavering north and south (and through) the equatorial plane and their orbital plane, if such a mathematical object can be defined, is tilted with respect to the equatorial plane. The maximum angular distance from the equatorial plane is reached when $`\left|U^\vartheta \right|=0`$, which is when $`\mathrm{sin}\vartheta =\left|S\right|`$. We call $`S`$ the "tilt of the orbital plane with respect to the equatorial plane", in analogy to the tilt of a planetary orbit with respect to the Sun's equatorial plane, which also is equal to the angle where the planet reaches its greatest angular distance from the equatorial plane. While we already have seen that the tilt of a planetary orbit is only one of two constants that define the orientation of the orbital plane, we still have to see whether there is such a second constant in the case of our model-space-time geodesics.
We can learn much about the orbital shape and size from investigating the radial velocity. If $`AX^2`$ then $`\left(U^r\right)^2`$ is a second order function of $`1/r`$ and otherwise it is a first order function of $`1/r`$. The coefficient of the $`1/r^2`$-term is zero or negative, because $`AX^2`$. From that follows that $`\left(U^r\right)^2`$ has a maximum if $`AX^2`$. If the equation $`\left(U^r\right)^2=0`$has no solutions, then this maximum is negative, which is forbidden, $`\left(U^r\right)^2`$ being the square of a real number and therefore necessarily non-negative. We can investigate the constraints on $`P,X,W`$ and $`A`$ that ensure that
$$\left(U^r\right)^2=\frac{AX^2}{r^2}+\frac{2PX}{r}+P^2W=0$$
$`(19)`$
has solutions.
The solutions of $`1/r`$ are:
$$\frac{1}{r}|_{U^r=0}=\frac{PX\pm \sqrt{AP^2AW+X^2W}}{AX^2}.$$
$`(20)`$
The condition that there be solutions is that
$$P^2W\left(1\frac{X^2}{A}\right).$$
$`(21)`$
The factor in parentheses is non-negative, because $`AX^2`$. If $`W0`$, then the relation is fulfilled for all values of $`P`$. For positive $`W`$ the above relation puts a lower bound on the absolute value of $`P`$.
Because the radial parameter $`r`$ is non-negative, we are only interested in non-negative solutions of $`1/r`$. If there is one positive solution, then the trajectory is unbound: the positive solution is the point of closest approach, but there is no point of greatest radial distance. If the other solution is zero, then the trajectory is just barely unbound and has the status of a parabolic trajectory in a central force field. If the other solution is negative, the trajectory is "hyperbolic". If there are two positive solutions, the trajectory is bound, like the elliptic trajectories in a central field with a $`1/r`$ potential. If the two positive solutions are equal, the trajectory has constant radius, i.e. it is comparable with circular orbits.
The close analogy between classical orbits in central force fields and geodesics in our space-time model is concisely expressed by
$$\frac{P^2W}{2}=\frac{\left(U^r\right)^2+\left(U^{}\right)^2}{2}\frac{XP}{r}\frac{X^2}{2r^2},$$
$`(22a)`$
where
$$\left(U^{}\right)^2=\left(rU^\vartheta \right)^2+\mathrm{sin}^2\vartheta (U^trU^\phi )^2$$
$${}_{}{}^{\prime \prime }squareoftangentialvelocity^{\prime \prime }.$$
$`(22b)`$
Eq. (22a) unmistakably has the signature of an energy, having a kinetic part depending on the squares of the radial and tangential velocities and two potential parts, one of which is due to a long range $`1/r`$ potential that can be attractive or repelling, like a Coulomb potential, while the other is due to a short range attractive $`1/r^2`$ potential. The $`1/r`$ potential gives rise to orbits having circular, elliptic, parabolic or hyperbolic shapes, while the $`1/r^2`$ potential adds a precession of pericentrum to the movement, like the shift of the perihelion of Mercury that also is caused by a $`1/r^2`$ term.
The energy expression \[Eq. (22a)\] contains two independent constants $`P`$ and $`X`$ that together define shape and size of an orbit in the same way that shape and size of planetary orbits are defined by the mass of the sun and the energy per kilogram of planetary mass, where $`X`$ plays the role of the solar mass and $`\left(P^2W\right)/2`$ the role of energy per unit of planetary mass.
Now we have "used" three constants to specify the tilt of the orbital plane \[Eq. (18)\], the size of the orbit \[Eq. (20)\] and the shape of the orbit \[Eq. (22a)\] and there is no constant to completely specify the orientation of the orbital plane. This situation is explained by the fact that the orbital plane has to rotate to ensure that a geodesic test particle does not leave the orbital plane. The pace $`\mathrm{\Omega }=\frac{d\psi }{dt}`$ with which the intersection line $`\left\{\vartheta =\pi /2,\phi =\psi \left(t\right)\right\}`$ of orbital plane and the equatorial plane rotates depends on the instantaneous radial distance $`r`$ of the test particle:
$$\mathrm{\Omega }=1/r.$$
$`(23)`$
The movement of the test particle for the case where the orbit has constant $`r`$ ("circular" orbit) is accurately modeled by the movement of a point on the rim of a coin that is set to spin on its side on a table. The coin may start almost upright, slowly falling due to frictional dissipation and decreasing its tilt with respect to the surface of the table until it lies down flat on the table. In the model space-time, of course, there is no friction and the initial tilt remains the same forever. While the coin is wobbling on the table, it rolls on its surface, which means that points on the rim not only take part in the rotation of the plane of the coin, but also in a rotation around the axis that is perpendicular to the coin, moving up and down and around in a complex dance.
The movement becomes even more complex if the radial distance is not constant, but it can still be understood easily if one imagines that the movement takes place in a tilted orbital plane that rotates. Figures 5-6 give depict a geodesic from these two perspectives.
E.Spin component measurements
In a Stern-Gerlach experiment the direction in which a spinning particle is deflected depends on the direction of the force
$$F_z=(\mu .B)=\mu _z\frac{B_z}{z},$$
$`(24)`$
where $`\mu `$ is the magnetic moment of the particle associated with the spin of the particle and $`B`$ is the magnetic field of the magnet, having a strong gradient $`B_z/z`$. As seen from the particle's point of view, depending on the force, the pole where the magnetic field is strongest is either pushed away or attracted by the particle as it passes through the gap between the poles. This effect is similar to the force exerted on a geodesic test particle in the model space-time. If we look at Eq. (22a) we can see that the $`1/r`$ potential gives rise to either an attractive force or a repelling force, depending on the sign of $`X`$. This close analogy is the reason why we can say that the sign of $`X`$ corresponds to the outcome of a spin component measurement.
IV.CONCLUSION
If we accept that we can not know for sure whether space-time is flat at the scale (in space and time) at which the Bohm-Aharonov experiment is performed, then we can not evade the conclusion that Bell's and similar proofs, which are all based on counterfactual statements, do not apply. Non-local angles and other geometric relations over some distance involving counterfactual set-ups are not measurable and have, even from a realist point of view, no definite values, because non-local geometric relations in curved space-time require operational definitions, which are of course not applicable to counterfactual set-ups.
A proper description of the EPR and Bohm-Aharonov experiments amalgamates the indefiniteness of the traditional quantum description with the realist point of view of the traditional classical description. Such a description is in the spirit of general relativity and makes plausible that, contrary to common opinion, the philosophical foundation of general relativity is also fundamental to the proper interpretation of quantum mechanics as a statistical theory in a non-flat and largely unknown playground.
An analysis of the results of the Bohm-Aharonov experiment (which are assumed to agree with the predictions of quantum mechanics) indicates that a hidden variable can be introduced to explain the results, but that the configuration of three directions (the orientations of the measuring instruments and the direction of the hidden variable) has one degree of freedom less than expected classically. Classically, given two of the three angles that identify a configuration, the third angle can be chosen freely from a continuous spectrum. On the other hand, configurations that are compatible with the predictions of QM restrict the third angle to a bivalent choice. This restriction is not severe and does not introduce non-locality by itself: even classically the third angle is restricted.
A model based on the simple, even naive assumption that spin has to do with space time structure with rotational symmetry in one direction, exhibits geodesic movements with fewer orbit defining constants than are needed for a Keplerian orbit of a test particle in a central force field. Of the three constants, only one constant defines the orientation of the orbital plane. Relative to an orbital plane, any other orbital plane can, up to a bivalent choice, be specified with a single number. Thus the model has exactly the required property to ensure that the predictions of quantum mechanics can be reproduced.
A. Einstein, B. Podolsky & N. Rosen, Phys. Rev. 47, 777 (1935)
N. Bohr, Phys. Rev, 48, 696 (1935)
D. Bohm and Y. Aharonov, Phys. Rev. 108, 1070 (1957)
J.S. Bell, Physics 1, 195 (1964)
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# Classification of extremal elliptic 𝐾3 surfaces and fundamental groups of open 𝐾3 surfaces
## 1. Introduction
A complex elliptic $`K3`$ surface $`f:X^1`$ with a section $`O`$ is said to be extremal if the Picard number $`\rho (X)`$ of $`X`$ is $`20`$ and the Mordell-Weil group $`MW_f`$ of $`f`$ is finite. The purpose of this paper is to present the complete list of all extremal elliptic $`K3`$ surfaces. As an application, we show that, if an $`ADE`$-configuration of smooth rational curves on a $`K3`$ surface satisfies a certain condition, then the topological fundamental group of the complement is trivial. (See Theorem 4.3 for the precise statement.)
Let $`f:X^1`$ be an elliptic $`K3`$ surface with a section $`O`$. We denote by $`R_f`$ the set of all points $`v^1`$ such that $`f^1(v)`$ is reducible. For a point $`vR_f`$, let $`f^1(v)^\mathrm{\#}`$ be the union of irreducible components of $`f^1(v)`$ that are disjoint from the zero section $`O`$. It is known that the cohomology classes of irreducible components of $`f^1(v)^\mathrm{\#}`$ form a negative definite root lattice $`S_{f,v}`$ of type $`A_l`$, $`D_m`$ or $`E_n`$ in $`H^2(X;)`$. Let $`\tau (S_{f,v})`$ be the type of this lattice. We define $`\mathrm{\Sigma }_f`$ to be the formal sum of these types;
$$\mathrm{\Sigma }_f:=\underset{vR_f}{}\tau (S_{f,v}).$$
The Néron-Severi lattice $`NS_X`$ of $`X`$ is defined to be $`H^{1,1}(X)H^2(X;)`$, and the transcendental lattice $`T_X`$ of $`X`$ is defined to be the orthogonal complement of $`NS_X`$ in $`H^2(X;)`$. We call the triple $`(\mathrm{\Sigma }_f,MW_f,T_X)`$ the data of the elliptic $`K3`$ surface $`f:X^1`$. When $`f:X^1`$ is extremal, the transcendental lattice $`T_X`$ is a positive definite even lattice of rank $`2`$.
###### Theorem 1.1.
There exists an extremal elliptic $`K3`$ surface $`f:X^1`$ with data $`(\mathrm{\Sigma }_f,MW_f,T_X)`$ if and only if $`(\mathrm{\Sigma }_f,MW_f,T_X)`$ appears in Table 2 given at the end of this paper.
In Table 2, the transcendental lattice $`T_X`$ is expressed by the coefficients of its Gram matrix
$$\left(\begin{array}{cc}a& b\\ b& c\end{array}\right).$$
See Subsection 2.1 on how to recover the $`K3`$ surface $`X`$ from $`T_X`$.
The classification of semi-stable extremal elliptic $`K3`$ surfaces has been done by Miranda and Persson and complemented by Artal-Bartolo, Tokunaga and Zhang . We can check that the semi-stable part of our list (No. 1- No. 112) coincides with theirs. Nishiyama classified all elliptic fibrations (not necessarily extremal) on certain $`K3`$ surfaces. On the other hand, Ye has independently classified all extremal elliptic $`K3`$ surfaces with no semi-stable singular fibers by different methods from ours.
Acknowledgment. The authors would like to thank Professors Shigeyuki Kondō, Ken-ichi Nishiyama and Keiji Oguiso for helpful discussions.
## 2. Preliminaries
### 2.1. Transcendental lattice of singular $`K3`$ surfaces
Let $`𝒬`$ be the set of symmetric matrices
$$Q=\left(\begin{array}{cc}a& b\\ b& c\end{array}\right)$$
of integer coefficients such that $`a`$ and $`c`$ are even and that the corresponding quadratic forms are positive definite. The group $`GL_2()`$ acts on $`𝒬`$ from right by
$$Q{}_{}{}^{t}gQg,$$
where $`gGL_2()`$. Let $`Q_1`$ and $`Q_2`$ be two matrices in $`𝒬`$, and let $`L_1`$ and $`L_2`$ be the positive definite even lattices of rank $`2`$ whose Gram matrices are $`Q_1`$ and $`Q_2`$, respectively. Then $`L_1`$ and $`L_2`$ are isomorphic as lattices if and only if $`Q_1`$ and $`Q_2`$ are in the same orbit under the action of $`GL_2()`$. On the other hand, each orbit in $`𝒬`$ under the action of $`SL_2()`$ contains a unique matrix with coefficients satisfying
$$a<2bac,\text{with}b0\text{if}a=c.$$
(See, for example, Conway and Sloane \[3, p. 358\].) Hence each orbit in $`𝒬`$ under the action of $`GL_2()`$ contains a unique matrix with coefficients satisfying
(2.1)
$$02bac.$$
In Table 2, the transcendental lattice is represented by the Gram matrix satisfying the condition (2.1).
Let $`X`$ be a $`K3`$ surface with $`\rho (X)=20`$; that is, $`X`$ is a singular $`K3`$ surface in the terminology of Shioda and Inose . The transcendental lattice $`T_X`$ can be naturally oriented by means of a holomorphic two form on $`X`$ (cf. \[16, p. 128\]). Let $`𝒮`$ denote the set of isomorphism classes of singular $`K3`$ surfaces. Using the natural orientation on the transcendental lattice, we can lift the map $`𝒮𝒬/GL_2()`$ given by $`XT_X`$ to the map $`𝒮𝒬/SL_2()`$.
###### Proposition 2.1 (Shioda and Inose ).
This map $`𝒮𝒬/SL_2()`$ is bijective. ∎
Moreover, Shioda and Inose gave us a method to construct explicitly the singular $`K3`$ surface corresponding to a given element of $`𝒬/SL_2()`$ by means of Kummer surfaces. The injectivity of the map $`𝒮𝒬/SL_2()`$ had been proved by Piateskii-Shapiro and Shafarevich .
Suppose that an orbit $`[Q]𝒬/GL_2()`$ is represented by a matrix $`Q`$ satisfying (2.1). Let $`\rho :𝒬/SL_2()𝒬/GL_2()`$ be the natural projection. Then we have
$$|\rho ^1([Q])|=\{\begin{array}{cc}2\hfill & \text{if }0<2b<a<c\hfill \\ 1\hfill & \text{otherwise.}\hfill \end{array}$$
Therefore, if a data in Table 2 satisfies $`a=c`$ or $`b=0`$ or $`2b=a`$ (resp. $`0<2b<a<c`$), then the number of the isomorphism classes of $`K3`$ surfaces that possess a structure of the extremal elliptic $`K3`$ surfaces with the given data is one (resp. two).
### 2.2. Roots of a negative definite even lattice
Let $`M`$ be a negative definite even lattice. A vector of $`M`$ is said to be a root of $`M`$ if its norm is $`2`$. We denote by $`\mathrm{root}(M)`$ the number of roots of $`M`$, and by $`M_{root}`$ the sublattice of $`M`$ generated by the roots of $`M`$. Suppose that a Gram matrix $`(a_{ij})`$ of $`M`$ is given. Then $`\mathrm{root}(M)`$ can be calculated by the following method. Let
$$g_r(x)=\underset{i,j=1}{\overset{r}{}}a_{ij}x_ix_j$$
be the positive definite quadratic form associated with the opposite lattice $`M^{}`$ of $`M`$, where $`r`$ is the rank of $`M`$. We consider the bounded closed subset
$$E(g_r,2):=\{x^r;g_r(x)2\}$$
of $`^r`$. Then we have
$$\mathrm{root}(M)+1=|E(g_r,2)^r|,$$
where $`+1`$ comes from the origin. For a positive integer $`k`$ less than $`r`$, we write by $`p_k:^r^k`$ the projection $`(x_1,\mathrm{},x_r)(x_1,\mathrm{},x_k)`$. Then there exist a positive definite quadratic form $`g_k`$ of variables $`(x_1,\mathrm{},x_k)`$ and a positive real number $`\sigma _k`$ such that
$$p_k(E(g_r,2))=E(g_k,\sigma _k):=\{y^k;g_k(y)\sigma _k\}.$$
The projection $`(x_1,\mathrm{},x_{k+1})(x_1,\mathrm{},x_k)`$ maps $`E(g_{k+1},\sigma _{k+1})`$ to $`E(g_k,\sigma _k)`$. Hence, if we have the list of the points of $`E(g_k,\sigma _k)^k`$, then it is easy to make the list of the points of $`E(g_{k+1},\sigma _{k+1})^{k+1}`$. Thus, starting from $`E(g_1,\sigma _1)`$, we can make the list of the points of $`E(g_r,2)^r`$ by induction on $`k`$.
### 2.3. Root lattices of type $`ADE`$
A root type is, by definition, a finite formal sum $`\mathrm{\Sigma }`$ of $`A_l`$, $`D_m`$ and $`E_n`$ with non-negative integer coefficients;
$$\mathrm{\Sigma }=\underset{l1}{}a_lA_l+\underset{m4}{}d_mD_m+\underset{n=6}{\overset{8}{}}e_nE_n.$$
We denote by $`L(\mathrm{\Sigma })`$ the negative definite root lattice corresponding to $`\mathrm{\Sigma }`$. The rank of $`L(\mathrm{\Sigma })`$ is given by
$$\mathrm{rank}(L(\mathrm{\Sigma }))=\underset{l1}{}a_ll+\underset{m4}{}d_mm+\underset{n=6}{\overset{8}{}}e_nn,$$
and the number of roots of $`L(\mathrm{\Sigma })`$ is given by
(2.2)
$$\mathrm{root}(L(\mathrm{\Sigma }))=\underset{l1}{}a_l(l^2+l)+\underset{m4}{}d_m(2m^22m)+72e_6+126e_7+240e_8.$$
(See, for example, Bourbaki .) Because of $`L(\mathrm{\Sigma })_{root}=L(\mathrm{\Sigma })`$, we have
(2.3)
$$L(\mathrm{\Sigma }_1)L(\mathrm{\Sigma }_2)\mathrm{\Sigma }_1=\mathrm{\Sigma }_2.$$
We also define $`\mathrm{𝑒𝑢}(\mathrm{\Sigma })`$ by
$$\mathrm{𝑒𝑢}(\mathrm{\Sigma }):=\underset{l1}{}a_l(l+1)+\underset{m4}{}d_m(m+2)+\underset{n=6}{\overset{8}{}}e_n(n+2).$$
###### Lemma 2.2.
Let $`f:X^1`$ be an elliptic $`K3`$ surface. Then $`\mathrm{𝑒𝑢}(\mathrm{\Sigma }_f)`$ is at most $`24`$. Moreover, if $`\mathrm{𝑒𝑢}(\mathrm{\Sigma }_f)<24`$, then there exists at least one singular fiber of type $`\mathrm{I}_1`$, $`\mathrm{II}`$, $`\mathrm{III}`$ or $`\mathrm{IV}`$.
###### Proof.
Let $`e(Y)`$ denote the topological euler number of a $`CW`$-complex $`Y`$. Then $`e(X)=24`$ is equal with the sum of topological euler numbers of singular fibers of $`f`$. Every singular fiber has a positive topological euler number. We have defined $`\mathrm{𝑒𝑢}(\mathrm{\Sigma })`$ in such a way that, if $`vR_f`$, then $`\mathrm{𝑒𝑢}(\tau (S_{f,v}))e(f^1(v))`$ holds, and if $`\mathrm{𝑒𝑢}(\tau (S_{f,v}))<e(f^1(v))`$, then the type of the fiber $`f^1(v)`$ is either $`\mathrm{III}`$ or $`\mathrm{IV}`$. Hence $`\mathrm{𝑒𝑢}(\mathrm{\Sigma }_f)`$ does not exceed the sum of the topological euler numbers of reducible singular fibers, and if $`\mathrm{𝑒𝑢}(\mathrm{\Sigma }_f)<24`$, then there is an irreducible singular fiber or a singular fiber of type $`\mathrm{III}`$ or $`\mathrm{IV}`$. ∎
### 2.4. Discriminant form and overlattices
Let $`L`$ be an even lattice, $`L^{}`$ the dual of $`L`$, $`D_L`$ the discriminant group $`L^{}/L`$ of $`L`$, and $`q_L`$ the discriminant form on $`D_L`$. (See Nikulin \[11, n. 4\] for the definitions.) An overlattice of $`L`$ is, by definition, an integral sublattice of the $``$-lattice $`L^{}`$ containing $`L`$.
###### Lemma 2.3 (Nikulin Proposition 1.4.2).
(1) Let $`A`$ be an isotopic subgroup of $`(D_L,q_L)`$. Then the pre-image $`M:=\varphi _L^1(A)`$ of $`A`$ by the natural projection $`\varphi _L:L^{}D_L`$ is an overlattice of $`L`$, and the discriminant form $`(D_M,q_M)`$ of $`M`$ is isomorphic to $`(A^{}/A,q_L|_{A^{}/A})`$, where $`A^{}`$ is the orthogonal complement of $`A`$ in $`D_L`$, and $`q_L|_{A^{}/A}`$ is the restriction of $`q_L`$ to $`A^{}/A`$. (2) The correspondence $`AM`$ gives a bijection from the set of isotopic subgroups of $`(D_L,q_L)`$ to the set of even overlattices of $`L`$. ∎
###### Lemma 2.4 (Nikulin Corollary 1.6.2).
Let $`S`$ and $`K`$ be two even lattices. Then the following two conditions are equivalent. (i) There is an isomorphism $`\gamma :D_S\stackrel{}{}D_K`$ of abelian groups such that $`\gamma ^{}q_K=q_S`$. (ii) There is an even unimodular overlattice of $`SK`$ into which $`S`$ and $`K`$ are primitively embedded. ∎
### 2.5. Néron-Severi groups of elliptic $`K3`$ surfaces
Let $`f:X^1`$ be an elliptic $`K3`$ surface with the zero section $`O`$. In the Néron-Severi lattice $`NS_X`$ of $`X`$, the cohomology classes of the zero section $`O`$ and a general fiber of $`f`$ generate a sublattice $`U_f`$ of rank $`2`$, which is isomorphic to the hyperbolic lattice
$$H:=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
Let $`W_f`$ be the orthogonal complement of $`U_f`$ in $`NS_X`$. Because $`U_f`$ is unimodular, we have $`NS_X=U_fW_f`$. Because $`U_f`$ is of signature $`(1,1)`$ and $`NS_X`$ is of signature $`(1,\rho (X)1)`$, $`W_f`$ is negative definite of rank $`\rho (X)2`$. Note that $`W_f`$ contains the sublattice
$$S_f:=\underset{vR_f}{}S_{f,v}$$
generated by the cohomology classes of irreducible components of reducible fibers of $`f`$ that are disjoint from the zero section. By definition, $`S_f`$ is isomorphic to $`L(\mathrm{\Sigma }_f)`$.
###### Lemma 2.5 (Nishiyama Lemma 6.1).
The sublattice $`S_f`$ of $`W_f`$ coincides with $`(W_f)_{root}`$, and the Mordell-Weil group $`MW_f`$ of $`f`$ is isomorphic to $`W_f/S_f`$. In particular, $`\mathrm{root}(L(\mathrm{\Sigma }_f))`$ is equal with $`\mathrm{root}(W_f)`$. ∎
Because $`W_fU_fT_X`$ has an even unimodular overlattice $`H^2(X;)`$ into which $`NS_X=W_fU_f`$ and $`T_X`$ are primitively embedded, and because the discriminant form of $`NS_X`$ is equal with the discriminant form of $`W_f`$ by $`D_{U_f}=(0)`$, Lemma 2.4 implies the following:
###### Corollary 2.6.
There is an isomorphism $`\gamma :D_{W_f}\stackrel{}{}D_{T_X}`$ of abelian groups such that $`\gamma ^{}q_{T_X}`$ coincides with $`q_{W_f}`$. ∎
### 2.6. Existence of elliptic $`K3`$ surfaces
Let $`\mathrm{\Lambda }`$ be the $`K3`$ lattice $`L(2E_8)H^3`$.
###### Lemma 2.7 (Kondō Lemma 2.1).
Let $`T`$ be a positive definite primitive sublattice of $`\mathrm{\Lambda }`$ with $`\mathrm{rank}(T)=2`$, and $`T^{}`$ the orthogonal complement of $`T`$ in $`\mathrm{\Lambda }`$. Suppose that $`T^{}`$ contains a sublattice $`H_T`$ isomorphic to the hyperbolic lattice. Let $`M_T`$ be the orthogonal complement of $`H_T`$ in $`T^{}`$. Then there exist an elliptic $`K3`$ surface $`f:X^1`$ such that $`T_XT`$ and $`W_fM_T`$.
###### Proof.
By the surjectivity of the period map of the moduli of $`K3`$ surfaces (cf. Todorov ), there exist a $`K3`$ surface $`X`$ and an isomorphism $`\alpha :H^2(X;)\mathrm{\Lambda }`$ of lattices such that $`\alpha ^1(T)=T_X`$. By Kondō \[5, Lemma 2.1\], the $`K3`$ surface $`X`$ has an elliptic fibration $`f:X^1`$ with a section such that $`[F]^{}/[F]M_T`$, where $`[F]U_f`$ is the cohomology class of a fiber of $`f`$, and $`[F]^{}`$ is the orthogonal complement of $`[F]`$ in the Néron-Severi lattice $`NS_X`$. Because $`NS_X`$ is equal with $`U_fW_f`$, and because $`[F]^{}U_f`$ coincides with $`[F]`$, we see that $`[F]^{}/[F]`$ is isomorphic to $`W_f`$. ∎
### 2.7. Datum of extremal elliptic $`K3`$ surfaces
###### Proposition 2.8.
A triple $`(\mathrm{\Sigma },MW,T)`$ consisting of a root type $`\mathrm{\Sigma }`$, a finite abelian group $`MW`$ and a positive definite even lattice $`T`$ of rank $`2`$ is a data of an extremal elliptic $`K3`$ surface if and only if the following hold:
1. $`\mathrm{length}(MW)2`$, $`\mathrm{rank}(L(\mathrm{\Sigma }))=18`$ and $`\mathrm{𝑒𝑢}(\mathrm{\Sigma })24`$.
2. There exists an overlattice $`M`$ of $`L(\mathrm{\Sigma })`$ satisfying the following:
1. $`M/L(\mathrm{\Sigma })MW`$,
2. there exists an isomorphism $`\gamma :D_M\stackrel{}{}D_T`$ of abelian groups such that $`\gamma ^{}q_T=q_M`$, and
3. $`\mathrm{root}(L(\mathrm{\Sigma }))=\mathrm{root}(M)`$.
###### Proof.
Suppose that there exists an extremal elliptic $`K3`$ surface $`f:X^1`$ with data equal with $`(\mathrm{\Sigma },MW,T)`$. It is obvious that $`\mathrm{\Sigma }`$ and $`MW`$ satisfies the condition $`(D1)`$. Via the isomorphism $`S_fL(\mathrm{\Sigma })`$, the overlattice $`W_f`$ of $`S_f`$ corresponds to an overlattice $`M`$ of $`L(\mathrm{\Sigma })`$, which satisfies the conditions $`(D2\text{-}a)`$-$`(D2\text{-}c)`$ by Lemma 2.5 and Corollary 2.6. Conversely, suppose that $`(\mathrm{\Sigma },MW,T)`$ satisfies the conditions $`(D1)`$ and $`(D2)`$. By Lemma 2.4, the condition $`(D2\text{-}b)`$ and $`D_H=0`$ imply that there exists an even unimodular overlattice of $`MHT`$ into which $`MH`$ and $`T`$ are primitively embedded. By the theorem of Milnor (see, for example, Serre ) on the classification of even unimodular lattices, any even unimodular lattice of signature $`(3,19)`$ is isomorphic to the $`K3`$ lattice $`\mathrm{\Lambda }`$. Then Lemma 2.7 implies that there exists an elliptic $`K3`$ surface $`f:X^1`$ satisfying $`W_fM`$ and $`T_XT`$. The condition $`(D2\text{-}c)`$ implies $`M_{root}=L(\mathrm{\Sigma })`$. Combining this with Lemma 2.5, we see that $`S_fL(\mathrm{\Sigma })`$. Then (2.2) implies that $`\mathrm{\Sigma }_f=\mathrm{\Sigma }`$. Using Lemma 2.5 and the condition $`(D2\text{-}a)`$, we see that $`MW_fMW`$. Thus the data of $`f:X^1`$ coincides with $`(\mathrm{\Sigma },MW,T)`$. ∎
###### Remark 2.9.
In the light of Lemma 2.3, the condition $`(D2)`$ is equivalent to the following:
1. There exists an isotopic subgroup $`A`$ of $`(D_{L(\mathrm{\Sigma })},q_{L(\mathrm{\Sigma })})`$ satisfying the following:
1. $`A`$ is isomorphic to $`MW`$,
2. there exists an isomorphism $`\gamma :A^{}/A\stackrel{}{}D_T`$ of abelian groups such that $`\gamma ^{}q_T=q_{L(\mathrm{\Sigma })}|_{A^{}/A}`$, and
3. $`\mathrm{root}(\varphi _{L(\mathrm{\Sigma })}^1(A))`$ is equal with $`\mathrm{root}(L(\mathrm{\Sigma }))`$, where $`\varphi _{L(\mathrm{\Sigma })}:L(\mathrm{\Sigma })^{}D_{L(\mathrm{\Sigma })}`$ is the natural projection.
###### Remark 2.10.
We did not use the conditions $`\mathrm{length}(MW)2`$ and $`\mathrm{𝑒𝑢}(\mathrm{\Sigma })24`$ in the proof of the “ if ” part of Proposition 2.8. It follows that, if $`(\mathrm{\Sigma },MW,T)`$ satisfies $`\mathrm{rank}(L(\mathrm{\Sigma }))=18`$ and the condition $`(D2)`$, then $`\mathrm{length}(MW)2`$ and $`\mathrm{𝑒𝑢}(\mathrm{\Sigma })24`$ follow automatically. This fact can be used when we check the computer program described in the next section.
## 3. Making the list
First we list up all root types $`\mathrm{\Sigma }`$ satisfying $`\mathrm{rank}(L(\mathrm{\Sigma }))=18`$ and $`\mathrm{𝑒𝑢}(\mathrm{\Sigma })24`$. This list $``$ consists of $`712`$ elements.
Next we run a program that takes an element $`\mathrm{\Sigma }`$ of the list $``$ as an input and proceeds as follows.
###### Step 1.
The program calculates the intersection matrix of $`L(\mathrm{\Sigma })^{}`$. Using this matrix, it calculates the discriminant form of $`L(\mathrm{\Sigma })`$, and decomposes it into $`p`$-parts;
$$(D_{L(\mathrm{\Sigma })},q_{L(\mathrm{\Sigma })})=\underset{p}{}(D_{L(\mathrm{\Sigma })},q_{L(\mathrm{\Sigma })})_p,$$
where $`p`$ runs through the set $`\{p_1,\mathrm{},p_k\}`$ of prime divisors of the discriminant $`|D_{L(\mathrm{\Sigma })}|`$ of $`L(\mathrm{\Sigma })`$. We write the $`p_i`$-part of $`(D_{L(\mathrm{\Sigma })},q_{L(\mathrm{\Sigma })})`$ by $`(D_{L(\mathrm{\Sigma }),i},q_{L(\mathrm{\Sigma }),i})`$.
###### Step 2.
For each $`p_i`$, it calculates the set $`I(p_i)`$ of all pairs $`(A,A^{})`$ of an isotopic subgroup $`A`$ of $`(D_{L(\mathrm{\Sigma }),i},q_{L(\mathrm{\Sigma }),i})`$ and its orthogonal complement $`A^{}`$ such that $`\mathrm{length}(A)2`$.
###### Step 3.
For each element
$$𝒜:=((A_1,A_1^{}),\mathrm{},(A_k,A_k^{}))I(p_1)\times \mathrm{}\times I(p_k),$$
it calculates the $`/2`$-valued quadratic form
$$q_𝒜:=q_{L(\mathrm{\Sigma }),1}|_{A_1^{}/A_1}\times \mathrm{}\times q_{L(\mathrm{\Sigma }),k}|_{A_k^{}/A_k}$$
on the finite abelian group
$$D_𝒜:=A_1^{}/A_1\times \mathrm{}\times A_k^{}/A_k.$$
Let $`d(𝒜)`$ be the order of $`D_𝒜`$.
###### Step 4.
It generates the list $`𝒯(d(𝒜))`$ of positive definite even lattices of rank $`2`$ with discriminant equal with $`d(𝒜)`$. For each $`T𝒯(d(𝒜))`$, it calculates the discriminant form of $`T`$ and decomposes it into $`p`$-parts. If $`D_T`$ is isomorphic to $`D_𝒜`$ and $`q_T`$ is isomorphic to $`q_𝒜`$, then it proceeds to the next step. Note that the automorphism group of a finite abelian $`p`$-group of length $`2`$ is easily calculated, and hence it is an easy task to check whether two given quadratic forms on the finite abelian $`p`$-group of length $`2`$ are isomorphic or not.
###### Step 5.
It calculates the Gram matrix of the sublattice $`\stackrel{~}{L}(𝒜)`$ of $`L(\mathrm{\Sigma })^{}`$ generated by $`L(\mathrm{\Sigma })L(\mathrm{\Sigma })^{}`$ and the pull-backs of generators of the subgroups $`A_iD_{L(\mathrm{\Sigma }),i}`$ by the projection $`L(\mathrm{\Sigma })^{}D_{L(\mathrm{\Sigma })}D_{L(\mathrm{\Sigma }),i}`$. Then it calculates $`\mathrm{root}(\stackrel{~}{L}(𝒜))`$ by the method described in the subsection 2.2. If $`\mathrm{root}(\stackrel{~}{L}(𝒜))`$ is equal with $`\mathrm{root}(L(\mathrm{\Sigma }))`$ calculated by (2.2), then it puts out the pair of the finite abelian group
$$MW:=A_1\times \mathrm{}\times A_k$$
and the lattice $`T`$.
Then $`(\mathrm{\Sigma },MW,T)`$ satisfies the conditions $`(D1)`$ and $`(D3)`$, and all triples $`(\mathrm{\Sigma },MW,T)`$ satisfying $`(D1)`$ and $`(D3)`$ are obtained by this program.
## 4. Fundamental groups of open $`K3`$ surfaces
A simple normal crossing divisor $`\mathrm{\Delta }`$ on a $`K3`$ surface $`X`$ is said to be an $`ADE`$-configuration of smooth rational curves if each irreducible component of $`\mathrm{\Delta }`$ is a smooth rational curve and the intersection matrix of the irreducible components of $`\mathrm{\Delta }`$ is a direct sum of the Cartan matrices of type $`A_l`$, $`D_m`$ or $`E_n`$ multiplied by $`1`$. It is known that $`\mathrm{\Delta }`$ is an $`ADE`$-configuration of smooth rational curves if and only if each connected component of $`\mathrm{\Delta }`$ can be contracted to a rational double point. We consider the following quite plausible hypothesis. Let $`\mathrm{\Delta }`$ be an $`ADE`$-configuration of smooth rational curves on a $`K3`$ surface $`X`$.
Hypothesis. If $`\pi _1^{alg}(X\mathrm{\Delta })`$ is trivial, then so is $`\pi _1(X\mathrm{\Delta })`$.
Here $`\pi _1^{alg}(X\mathrm{\Delta })`$ is the algebraic fundamental group of $`X\mathrm{\Delta }`$, which is the pro-finite completion of the topological fundamental group $`\pi _1(X\mathrm{\Delta })`$.
###### Proposition 4.1.
Suppose that Hypothesis is true for any $`ADE`$-configuration of smooth rational curves on an arbitrary $`K3`$ surface. Let $`\mathrm{\Delta }`$ be an $`ADE`$-configuration of smooth rational curves on a $`K3`$ surface $`X`$. Then $`\pi _1(X\mathrm{\Delta })`$ satisfies one of the following:
1. $`\pi _1(X\mathrm{\Delta })`$ is trivial.
2. There exist a complex torus $`T`$ of dimension $`2`$ and a finite automorphism group $`G`$ of $`T`$ such that $`T/G`$ is birational to $`X`$ and that $`\pi _1(X\mathrm{\Delta })`$ fits in the exact sequence
$$1\pi _1(T)\pi _1(X\mathrm{\Delta })G1.$$
3. $`\pi _1(X\mathrm{\Delta })`$ is isomorphic to a symplectic automorphism group of a $`K3`$ surface.
###### Remark 4.2.
Fujiki classified the automorphism groups of complex tori of dimension $`2`$. In particular, the $`G`$ in (ii) is either one of $`/(n)`$ ($`n=2,3,4,6`$), $`Q_8`$ (Quaternion of order 8), $`D_{12}`$ (Dihedral of order 12) and $`T_{24}`$ (Tetrahedral of order 24), whence the $`\pi _1(X\mathrm{\Delta })`$ in (ii) is a soluble group. Mukai presented the complete list of symplectic automorphism groups of $`K3`$ surfaces. (See also Kondō and Xiao .) Under Hypothesis, therefore, we know what groups can appear as $`\pi _1(X\mathrm{\Delta })`$.
Proof of Proposition 4.1. Suppose that $`\pi _1(X\mathrm{\Delta })`$ is non-trivial. By Hypothesis, $`\pi _1^{alg}(X\mathrm{\Delta })`$ is also non-trivial. For a surjective homomorphism $`\varphi :\pi _1(X\mathrm{\Delta })G`$ from $`\pi _1(X\mathrm{\Delta })`$ to a finite group $`G`$, we denote by
$$\psi _\varphi :\stackrel{~}{Y}_\varphi X$$
the finite Galois cover of $`X`$ corresponding to $`\varphi `$, which is étale over $`X\mathrm{\Delta }`$ and whose Galois group is canonically isomorphic to $`G`$. Let $`\rho :\stackrel{~}{Y}_\varphi ^{}\stackrel{~}{Y}_\varphi `$ be the resolution of singularities, and $`\gamma :\stackrel{~}{Y}_\varphi ^{}Y_\varphi `$ the contraction of $`(1)`$-curves. We denote by $`\mathrm{\Delta }_\varphi `$ the union of one-dimensional irreducible components of $`\gamma (\rho ^1(\psi _\varphi ^1(\mathrm{\Delta })))`$. Then it is easy to see that $`Y_\varphi `$ is either a $`K3`$ surface or a complex torus of dimension $`2`$, and that the Galois group $`G`$ of $`\psi _\varphi `$ acts on $`Y_\varphi `$ symplectically. Moreover, $`\mathrm{\Delta }_\varphi `$ is an empty set or an $`ADE`$-configuration of smooth rational curves. We have an exact sequence
$$1\pi _1(Y_\varphi \mathrm{\Delta }_\varphi )\pi _1(X\mathrm{\Delta })G\mathrm{\hspace{0.33em}1},$$
because $`\pi _1(\stackrel{~}{Y}_\varphi \psi _\varphi ^1(\mathrm{\Delta }))`$ is isomorphic to $`\pi _1(Y_\varphi \mathrm{\Delta }_\varphi )`$. Suppose that there exists $`\varphi :\pi _1(X\mathrm{\Delta })G`$ such that $`Y_\varphi `$ is a complex torus of dimension $`2`$. Then $`\mathrm{\Delta }_\varphi `$ is empty, and hence (ii) occurs. Suppose that no complex tori of dimension $`2`$ appear as a finite Galois cover of $`X`$ branched in $`\mathrm{\Delta }`$. Then any finite quotient group of $`\pi _1(X\mathrm{\Delta })`$ must appear in Mukai’s list of symplectic automorphism groups of $`K3`$ surfaces. Because this list consists of finite number of isomorphism classes of finite groups, there exists a maximal finite quotient $`\varphi _{max}:\pi _1(X\mathrm{\Delta })G_{max}`$ of $`\pi _1(X\mathrm{\Delta })`$. Then $`\pi _1(Y_{\varphi _{max}}\mathrm{\Delta }_{\varphi _{max}})`$ has no non-trivial finite quotient group, and hence it is trivial by Hypothesis. Thus (iii) occurs. ∎
For an $`ADE`$-configuration $`\mathrm{\Delta }`$ of smooth rational curves on a $`K3`$ surface $`X`$, we denote by $`[\mathrm{\Delta }]`$ the sublattice of $`H^2(X;)`$ generated by the cohomology classes of the irreducible components of $`\mathrm{\Delta }`$, which is isomorphic to a negative definite root lattice of type $`ADE`$. We denote by $`\mathrm{\Sigma }_\mathrm{\Delta }`$ the root type such that $`[\mathrm{\Delta }]`$ is isomorphic to $`L(\mathrm{\Sigma }_\mathrm{\Delta })`$. Using the list of extremal elliptic $`K3`$ surfaces, we prove the following theorem. We first consider the following conditions on a root type $`\mathrm{\Sigma }`$ (see (2.4) for the definition of $`D_{L(\mathrm{\Sigma })}`$).
1. $`\mathrm{rank}(L(\mathrm{\Sigma }))18`$, and
2. $`\mathrm{length}(D_{L(\mathrm{\Sigma })})20\mathrm{rank}(L(\mathrm{\Sigma }))`$.
###### Theorem 4.3.
Let $`X`$ be a $`K3`$ surface and $`\mathrm{\Delta }`$ an $`ADE`$-configuration of smooth rational curves on $`X`$. Suppose that the root type $`\mathrm{\Sigma }_\mathrm{\Delta }`$ satisfies conditions $`(N1)`$ and $`(N2)`$. If $`[\mathrm{\Delta }]`$ is primitive in $`H^2(X;)`$ then $`\pi _1(X\mathrm{\Delta })`$ is trivial.
In virtue of Lemma 4.6 below, we can easily derive the following:
###### Corollary 4.4.
Let $`X`$ be a $`K3`$ surface and $`\mathrm{\Delta }`$ an $`ADE`$-configuration of smooth rational curves on $`X`$. Suppose that $`\mathrm{\Sigma }_\mathrm{\Delta }`$ satisfies the conditions $`(N1)`$ and $`(N2)`$. Then Hypothesis is true. ∎
###### Remark 4.5.
The conditions $`(N1)`$ and $`(N2)`$ come from Nikulin \[11, Theorem 1.14.1\] (see also Morrison \[8, Theorem 2.8\]), which gives a sufficient condition for the uniqueness of the primitive embedding of $`L(\mathrm{\Sigma })`$ into the $`K3`$ lattice $`\mathrm{\Lambda }`$.
First we prepare some lemmas. Let $`\overline{[\mathrm{\Delta }]}`$ be the primitive closure of $`[\mathrm{\Delta }]`$ in $`H^2(X;)`$.
###### Lemma 4.6 (Xiao Lemma 2).
The dual of the abelianisation of $`\pi _1(X\mathrm{\Delta })`$ is canonically isomorphic to $`\overline{[\mathrm{\Delta }]}/[\mathrm{\Delta }]`$. In particular, if $`\pi _1^{alg}(X\mathrm{\Delta })`$ is trivial, then $`[\mathrm{\Delta }]`$ is primitive in $`H^2(X;)`$. ∎
Let $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ be graphs with the set of vertices denoted by $`\mathrm{Vert}(\mathrm{\Gamma }_1)`$ and $`\mathrm{Vert}(\mathrm{\Gamma }_2)`$, respectively. An embedding of $`\mathrm{\Gamma }_1`$ into $`\mathrm{\Gamma }_2`$ is, by definition, an injection $`f:\mathrm{Vert}(\mathrm{\Gamma }_1)\mathrm{Vert}(\mathrm{\Gamma }_2)`$ such that, for any $`u,v\mathrm{Vert}(\mathrm{\Gamma }_1)`$, $`f(u)`$ and $`f(v)`$ are connected by an edge of $`\mathrm{\Gamma }_2`$ if and only if $`u`$ and $`v`$ are connected by an edge of $`\mathrm{\Gamma }_1`$.
Let $`\mathrm{\Gamma }(\mathrm{\Sigma })`$ denote the Dynkin graph of $`\mathrm{\Sigma }`$.
###### Lemma 4.7.
Suppose that $`\mathrm{\Sigma }`$ satisfies the conditions $`(N1)`$ and $`(N2)`$. Then there exists $`\mathrm{\Sigma }^{}`$ satisfying $`\mathrm{rank}(L(\mathrm{\Sigma }^{}))=18`$ and the condition $`(N2)`$ such that $`\mathrm{\Gamma }(\mathrm{\Sigma })`$ can be embedded in $`\mathrm{\Gamma }(\mathrm{\Sigma }^{})`$.
###### Proof.
This is checked by listing up all $`\mathrm{\Sigma }`$ satisfying the conditions $`(N1)`$ and $`(N2)`$ using computer. ∎
###### Lemma 4.8.
Let $`f:X^1`$ be an elliptic surface with the zero section $`O`$. Suppose that a fiber $`f^1(v)`$ over $`v^1`$ is a singular fiber of type $`\mathrm{III}`$ or $`\mathrm{IV}`$. Let $`\mathrm{\Xi }`$ be a union of some irreducible components of $`f^1(v)`$ that does not coincide with the whole fiber $`f^1(v)`$. If $`U`$ is a small open disk on $`^1`$ with the center $`v`$, then $`f^1(U)(\mathrm{\Xi }(f^1(U)O))`$ has an abelian fundamental group.
###### Proof.
This can be proved easily by the van-Kampen theorem. ∎
###### Lemma 4.9.
Let $`\mathrm{\Sigma }`$ be satisfying the conditions $`(N1)`$ and $`(N2)`$. Suppose that $`(X,\mathrm{\Delta })`$ and $`(X^{},\mathrm{\Delta }^{})`$ satisfy the following:
1. $`\mathrm{\Sigma }_\mathrm{\Delta }=\mathrm{\Sigma }_\mathrm{\Delta }^{}=\mathrm{\Sigma }`$,
2. $`\overline{[\mathrm{\Delta }]}=[\mathrm{\Delta }]`$ and $`\overline{[\mathrm{\Delta }^{}]}=[\mathrm{\Delta }^{}]`$.
Then there exists a connected continuous family $`(X_t,\mathrm{\Delta }_t)`$ parameterized by $`t[0,1]`$ such that $`(X_0,\mathrm{\Delta }_0)=(X,\mathrm{\Delta })`$, $`(X_1,\mathrm{\Delta }_1)=(X^{},\mathrm{\Delta }^{})`$ and that $`(X_t,\mathrm{\Delta }_t)`$ are diffeomorphic to one another. In particular, $`\pi _1(X\mathrm{\Delta })`$ is isomorphic to $`\pi _1(X^{}\mathrm{\Delta }^{})`$.
###### Proof.
By Nikulin \[11, Theorem 1.14.1\], the primitive embedding of $`L(\mathrm{\Sigma })`$ into the $`K3`$ lattice $`\mathrm{\Lambda }`$ is unique up to $`\mathrm{Aut}(\mathrm{\Lambda })`$. Hence the assertion follows from Nikulin’s connectedness theorem \[10, Theorem 2.10\]. ∎
Proof of Theorem 4.3. Let us consider the following:
###### Claim 1.
Suppose that $`\mathrm{\Sigma }`$ satisfies $`\mathrm{rank}(L(\mathrm{\Sigma }))=18`$ and the condition $`(N2)`$. Then there exists an $`ADE`$-configuration of smooth rational curves $`\mathrm{\Delta }_\mathrm{\Sigma }`$ on a $`K3`$ surface $`X_\mathrm{\Sigma }`$ such that $`\mathrm{\Sigma }_{\mathrm{\Delta }_\mathrm{\Sigma }}=\mathrm{\Sigma }`$ and $`\pi _1(X_\mathrm{\Sigma }\mathrm{\Delta }_\mathrm{\Sigma })=\{1\}`$.
We deduce Theorem 4.3 from Claim 1. Suppose that $`\mathrm{\Delta }`$ is an $`ADE`$-configuration of smooth rational curves on a $`K3`$ surface $`X`$ such that $`\mathrm{\Sigma }_\mathrm{\Delta }`$ satisfies the conditions $`(N1)`$ and $`(N2)`$, and that $`[\mathrm{\Delta }]`$ is primitive in $`H^2(X;)`$. By Lemma 4.7, there exists $`\mathrm{\Sigma }_1`$ satisfying $`\mathrm{rank}(L(\mathrm{\Sigma }_1))=18`$ and the condition $`(N2)`$ such that $`\mathrm{\Gamma }(\mathrm{\Sigma }_\mathrm{\Delta })`$ is embedded into $`\mathrm{\Gamma }(\mathrm{\Sigma }_1)`$. By Claim 1, we have $`(X_1,\mathrm{\Delta }_1)`$ such that $`\mathrm{\Sigma }_{\mathrm{\Delta }_1}=\mathrm{\Sigma }_1`$ and $`\pi _1(X_1\mathrm{\Delta }_1)=\{1\}`$. Let $`\mathrm{\Delta }^{}\mathrm{\Delta }_1`$ be the sub-configuration of smooth rational curves on $`X_1`$ corresponding to the subgraph $`\mathrm{\Gamma }(\mathrm{\Sigma }_\mathrm{\Delta })\mathrm{\Gamma }(\mathrm{\Sigma }_1)=\mathrm{\Gamma }(\mathrm{\Sigma }_{\mathrm{\Delta }_1})`$. There is a surjection from $`\pi _1(X_1\mathrm{\Delta }_1)`$ to $`\pi _1(X_1\mathrm{\Delta }^{})`$, and hence $`\pi _1(X_1\mathrm{\Delta }^{})`$ is trivial. In particular, $`[\mathrm{\Delta }^{}]`$ is primitive in $`H^2(X_1;)`$. Because of $`\mathrm{\Sigma }_\mathrm{\Delta }^{}=\mathrm{\Sigma }_\mathrm{\Delta }`$, Lemma 4.9 implies that $`\pi _1(X\mathrm{\Delta })`$ is isomorphic to $`\pi _1(X_1\mathrm{\Delta }^{})`$. Thus $`\pi _1(X\mathrm{\Delta })`$ is trivial.
Let $`f:X^1`$ be an extremal elliptic $`K3`$ surface. For a point $`vR_f`$, we denote the total fiber of $`f`$ over $`v`$ by
$$\underset{i=1}{\overset{r_v}{}}m_{v,i}C_{v,i},$$
where $`m_{v,i}`$ is the multiplicity of the irreducible component $`C_{v,i}`$ of $`f^1(v)`$. We denote by $`\mathrm{\Gamma }_f`$ the union of the zero section and all irreducible fibers $`f^1(v)`$ $`(vR_f)`$.
###### Claim 2.
Suppose that $`MW_f=(0)`$. Suppose that a sub-configuration $`\mathrm{\Delta }`$ of $`\mathrm{\Gamma }_f`$ satisfies the following two conditions.
1. The number of $`vR_f`$ such that $`m_{v,i}=1C_{v,i}\mathrm{\Delta }`$ holds is at most one.
2. Either one of the following holds:
1. The configuration $`\mathrm{\Delta }`$ does not contain the zero section,
2. there is a point $`v_1R_f`$ such that the type $`\tau (S_{f,v_1})`$ is $`A_1`$ and that $`F_1:=f^1(v_1)`$ and $`\mathrm{\Delta }`$ have no common irreducible components, or
3. $`\mathrm{𝑒𝑢}(\mathrm{\Sigma }_f)23`$.
Then $`\pi _1(X\mathrm{\Delta })`$ is trivial.
Proof of Claim 2. By Lemma 2.5, the assumption $`MW_f=(0)`$ implies that the cohomology classes $`[O]`$ and $`[C_{v,i}]`$ $`(vR_f,i=1,\mathrm{},r_v)`$ of the irreducible components of $`\mathrm{\Gamma }_f`$ span $`NS_X`$. The relations among these generators are generated by
$$\underset{i=1}{\overset{r_v}{}}m_{v,i}C_{v,i}=\underset{i=1}{\overset{r_v^{}}{}}m_{v^{},i}C_{v^{},i}(v,v^{}R_f).$$
Therefore the condition $`(Z1)`$ implies that the cohomology classes of the irreducible components of $`\mathrm{\Delta }`$ constitute a subset of a $``$-basis of $`NS_X`$. Hence $`[\mathrm{\Delta }]`$ is primitive in $`H^2(X;)`$. In particular, $`\pi _1(X\mathrm{\Delta })`$ is a perfect group by Lemma 4.6. On the other hand, the condition $`(Z1)`$ implies that there exists a point $`v_0^1`$ such that every fiber of the restriction
$$f|_{X(\mathrm{\Delta }f^1(v_0))}:X(\mathrm{\Delta }f^1(v_0))^1\{v_0\}$$
of $`f`$ has a reduced irreducible component. Then, by Nori’s lemma \[13, Lemma 1.5 (C)\], if $`U`$ is a non-empty connected classically open subset of $`^1\{v_0\}`$, then the inclusion of $`f^1(U)(f^1(U)\mathrm{\Delta }))`$ into $`X(\mathrm{\Delta }f^1(v_0))`$ induces a surjection on the fundamental groups. The inclusion of $`X(\mathrm{\Delta }f^1(v_0))`$ into $`X\mathrm{\Delta }`$ also induces a surjection on the fundamental groups. We shall show that there exists a small open disk $`U`$ on $`^1\{v_0\}`$ such that
$$G_U:=\pi _1(f^1(U)(f^1(U)\mathrm{\Delta }))$$
is abelian. When $`(Z2\text{-}a)`$ occurs, we take a small open disk disjoint from $`R_f`$ as $`U`$. Then $`G_U`$ is abelian, because of $`f^1(U)\mathrm{\Delta }=\mathrm{}`$. Suppose that $`(Z2\text{-}b)`$ occurs. We can take $`v_0`$ from $`^1\{v_1\}`$, because $`F_1`$ has no irreducible components of multiplicity $`2`$. We choose as $`U`$ a small open disk with the center $`v_1`$. There is a contraction from $`f^1(U)(f^1(U)\mathrm{\Delta })`$ to $`F_1(F_1\mathrm{\Delta })`$. Because $`\pi _1(F_1(F_1\mathrm{\Delta }))`$ is abelian, so is $`G_U`$. Suppose that $`(Z2\text{-}c)`$ occurs. By Lemma 2.2, there exists a singular fiber $`F_2:=f^1(v_2)`$ of type $`\mathrm{I}_1`$, $`\mathrm{II}`$, $`\mathrm{III}`$ or $`\mathrm{IV}`$. Because $`F_2`$ has no irreducible components of multiplicity $`2`$, we can choose $`v_0`$ from $`^1\{v_2\}`$. If $`F_2`$ is of type $`\mathrm{I}_1`$ or $`\mathrm{II}`$, then $`F_2\mathrm{\Delta }`$ consists of a nonsingular point of $`F_2`$, and $`\pi _1(F_2(F_2\mathrm{\Delta }))`$ is abelian. Hence $`G_U`$ is also abelian. If $`F_2`$ is of type $`\mathrm{III}`$ or $`\mathrm{IV}`$, then $`F_2\mathrm{\Delta }`$ cannot coincide with the whole fiber $`F_2`$. Hence Lemma 4.8 implies that $`G_U`$ is abelian. Therefore we see that $`\pi _1(X\mathrm{\Delta })`$ is abelian. Being both perfect and abelian, $`\pi _1(X\mathrm{\Delta })`$ is trivial. ∎
Now we proceed to the proof of Claim 1. We list up all $`\mathrm{\Sigma }`$ satisfying the condition $`(N2)`$ and $`\mathrm{rank}(L(\mathrm{\Sigma }))=18`$. It consists of $`297`$ elements. Among them, $`199`$ elements can be the type $`\mathrm{\Sigma }_f`$ of singular fibers of some extremal elliptic $`K3`$ surface $`f:X^1`$ with $`MW_f=0`$. For these configurations, $`\pi _1(X\mathrm{\Delta })`$ is trivial by Claim 2. The remaining $`98`$ configurations are listed in the second column of Table 1 below. Each of them is a sub-configuration of $`\mathrm{\Gamma }_f`$ satisfying the conditions $`(Z1)`$ and $`(Z2)`$, where $`f:X^1`$ is the extremal elliptic $`K3`$ surface with $`MW_f=0`$ whose number in Table 2 is given in the third column of Table 1. The fourth and fifth columns of Table 1 indicate $`\mathrm{\Sigma }_f`$ and $`\mathrm{𝑒𝑢}(\mathrm{\Sigma }_f)`$, respectively. In the case nos. 20, 28, 39, 41 and 85 in Table 1, we can choose the embedding of $`\mathrm{\Delta }`$ into $`\mathrm{\Gamma }_f`$ in such a way that $`(Z2\text{-}b)`$ holds. In the case nos. 30, 37, 57 and 63 in Table 1, we can choose the embedding of $`\mathrm{\Delta }`$ into $`\mathrm{\Gamma }_f`$ in such a way that $`(Z2\text{-}a)`$ holds. By Claim 2 again, $`\pi _1(X\mathrm{\Delta })`$ is trivial for these $`98`$ configurations $`\mathrm{\Delta }`$. ∎
###### Remark 4.10.
The graph $`\mathrm{\Gamma }(A_{19})`$ (resp. $`\mathrm{\Gamma }(D_{19})`$) can be embedded into $`\mathrm{\Gamma }_f`$ in such a way that $`(Z1)`$ and $`(Z2)`$ are satisfied, where $`f:X^1`$ is the extremal elliptic $`K3`$ surfaces whose number in Table 2 is 312 (resp. 320). Therefore, if $`\mathrm{\Gamma }(\mathrm{\Delta })`$ is embedded in $`\mathrm{\Gamma }(A_{19})`$ or $`\mathrm{\Gamma }(D_{19})`$, then $`\mathrm{\Gamma }(\mathrm{\Delta })`$ can be embedded in $`\mathrm{\Gamma }_f`$ in such a way that $`(Z1)`$ and $`(Z2)`$ are satisfied.
Table 1. List of embedding of $`\mathrm{\Delta }`$ in $`\mathrm{\Gamma }_f`$
| no | $`\mathrm{\Delta }`$ | No | $`\mathrm{\Sigma }_f`$ | $`\mathrm{𝑒𝑢}(\mathrm{\Sigma }_f)`$ |
| --- | --- | --- | --- | --- |
| 1 | $`A_2+A_3+2A_4+A_5`$ | 19 | $`A_2+2A_3+A_4+A_6`$ | $`23`$ |
| 2 | $`A_1+A_2+A_3+2A_6`$ | 23 | $`A_1+A_2+A_4+A_5+A_6`$ | $`23`$ |
| 3 | $`2A_1+A_4+2A_6`$ | 23 | $`A_1+A_2+A_4+A_5+A_6`$ | $`23`$ |
| 4 | $`2A_2+2A_4+A_6`$ | 23 | $`A_1+A_2+A_4+A_5+A_6`$ | $`23`$ |
| 5 | $`A_1+A_5+2A_6`$ | 40 | $`A_1+A_4+A_6+A_7`$ | $`22`$ |
| 6 | $`A_4+2A_7`$ | 52 | $`A_4+A_6+A_8`$ | $`21`$ |
| 7 | $`A_1+A_2+2A_4+A_7`$ | 23 | $`A_1+A_2+A_4+A_5+A_6`$ | $`23`$ |
| 8 | $`A_3+2A_4+A_7`$ | 24 | $`A_3+A_4+A_5+A_6`$ | $`22`$ |
| 9 | $`A_2+2A_4+A_8`$ | 36 | $`A_2+A_4+A_5+A_7`$ | $`22`$ |
| 10 | $`2A_3+A_4+A_8`$ | 46 | $`A_1+A_2+A_3+A_4+A_8`$ | $`23`$ |
| 11 | $`A_3+A_7+A_8`$ | 53 | $`A_1+A_2+A_7+A_8`$ | $`22`$ |
| 12 | $`A_1+2A_2+A_4+A_9`$ | 46 | $`A_1+A_2+A_3+A_4+A_8`$ | $`23`$ |
| 13 | $`A_2+A_3+A_4+A_9`$ | 71 | $`2A_2+A_4+A_{10}`$ | $`22`$ |
| 14 | $`A_3+A_4+A_{11}`$ | 93 | $`A_2+A_4+A_{12}`$ | $`21`$ |
| 15 | $`A_7+A_{11}`$ | 312 | $`A_{10}+E_8`$ | $`21`$ |
| 16 | $`2A_3+A_{12}`$ | 93 | $`A_2+A_4+A_{12}`$ | $`21`$ |
| 17 | $`A_3+A_{15}`$ | 312 | $`A_{10}+E_8`$ | $`21`$ |
| 18 | $`A_2+2A_6+D_4`$ | 99 | $`A_2+A_3+A_{13}`$ | $`21`$ |
| 19 | $`2A_4+A_6+D_4`$ | 18 | $`A_1+A_3+2A_4+A_6`$ | $`23`$ |
| 20 | $`2A_2+A_4+A_6+D_4`$ | 20 | $`A_1+2A_2+A_3+A_4+A_6`$ | $`24`$ |
| 21 | $`A_2+A_4+A_8+D_4`$ | 44 | $`2A_1+2A_4+A_8`$ | $`23`$ |
| 22 | $`A_6+A_8+D_4`$ | 50 | $`2A_1+A_2+A_6+A_8`$ | $`23`$ |
| 23 | $`2A_2+A_{10}+D_4`$ | 72 | $`2A_1+A_2+A_4+A_{10}`$ | $`23`$ |
| 24 | $`A_4+A_{10}+D_4`$ | 72 | $`2A_1+A_2+A_4+A_{10}`$ | $`23`$ |
| 25 | $`A_2+A_{12}+D_4`$ | 90 | $`2A_1+2A_2+A_{12}`$ | $`23`$ |
| 26 | $`A_{14}+D_4`$ | 320 | $`D_{10}+E_8`$ | $`22`$ |
| 27 | $`2A_2+A_4+2D_5`$ | 210 | $`2A_2+D_{14}`$ | $`22`$ |
| 28 | $`A_1+2A_2+2A_4+D_5`$ | 157 | $`A_1+A_2+2A_4+D_7`$ | $`24`$ |
| 29 | $`A_2+A_3+2A_4+D_5`$ | 46 | $`A_1+A_2+A_3+A_4+A_8`$ | $`23`$ |
| 30 | $`A_2+A_6+2D_5`$ | 193 | $`A_2+A_6+D_{10}`$ | $`22`$ |
| 31 | $`A_3+A_4+A_6+D_5`$ | 18 | $`A_1+A_3+2A_4+A_6`$ | $`23`$ |
| 32 | $`A_2+A_4+A_7+D_5`$ | 72 | $`2A_1+A_2+A_4+A_{10}`$ | $`23`$ |
| 33 | $`A_6+A_7+D_5`$ | 50 | $`2A_1+A_2+A_6+A_8`$ | $`23`$ |
| 34 | $`A_2+A_3+A_8+D_5`$ | 50 | $`2A_1+A_2+A_6+A_8`$ | $`23`$ |
| 35 | $`A_3+A_{10}+D_5`$ | 69 | $`A_1+2A_2+A_3+A_{10}`$ | $`23`$ |
| 36 | $`A_2+A_{11}+D_5`$ | 90 | $`2A_1+2A_2+A_{12}`$ | $`23`$ |
| 37 | $`A_4+2D_7`$ | 213 | $`A_4+D_{14}`$ | $`21`$ |
| 38 | $`A_3+2A_4+D_7`$ | 44 | $`2A_1+2A_4+A_8`$ | $`23`$ |
| 39 | $`2A_2+A_3+A_4+D_7`$ | 20 | $`A_1+2A_2+A_3+A_4+A_6`$ | $`24`$ |
| 40 | $`A_2+A_4+A_5+D_7`$ | 23 | $`A_1+A_2+A_4+A_5+A_6`$ | $`23`$ |
| 41 | $`A_1+2A_2+A_6+D_7`$ | 14 | $`2A_1+2A_2+2A_6`$ | $`24`$ |
Table 1. List of embedding of $`\mathrm{\Delta }`$ in $`\mathrm{\Gamma }_f`$
| no | $`\mathrm{\Delta }`$ | No | $`\mathrm{\Sigma }_f`$ | $`\mathrm{𝑒𝑢}(\mathrm{\Sigma }_f)`$ |
| --- | --- | --- | --- | --- |
| 42 | $`2A_2+A_7+D_7`$ | 90 | $`2A_1+2A_2+A_{12}`$ | $`23`$ |
| 43 | $`A_4+A_7+D_7`$ | 44 | $`2A_1+2A_4+A_8`$ | $`23`$ |
| 44 | $`A_1+A_2+A_8+D_7`$ | 50 | $`2A_1+A_2+A_6+A_8`$ | $`23`$ |
| 45 | $`A_3+A_8+D_7`$ | 44 | $`2A_1+2A_4+A_8`$ | $`23`$ |
| 46 | $`A_{11}+D_7`$ | 320 | $`D_{10}+E_8`$ | $`22`$ |
| 47 | $`A_2+A_4+D_5+D_7`$ | 200 | $`A_2+A_5+D_{11}`$ | $`22`$ |
| 48 | $`A_6+D_5+D_7`$ | 186 | $`A_9+D_9`$ | $`21`$ |
| 49 | $`A_2+2A_4+D_8`$ | 66 | $`A_2+A_7+A_9`$ | $`21`$ |
| 50 | $`A_4+A_6+D_8`$ | 23 | $`A_1+A_2+A_4+A_5+A_6`$ | $`23`$ |
| 51 | $`A_2+A_8+D_8`$ | 50 | $`2A_1+A_2+A_6+A_8`$ | $`23`$ |
| 52 | $`A_{10}+D_8`$ | 320 | $`D_{10}+E_8`$ | $`22`$ |
| 53 | $`A_1+2A_4+D_9`$ | 44 | $`2A_1+2A_4+A_8`$ | $`23`$ |
| 54 | $`A_2+A_3+A_4+D_9`$ | 46 | $`A_1+A_2+A_3+A_4+A_8`$ | $`23`$ |
| 55 | $`A_3+A_6+D_9`$ | 76 | $`2A_1+A_6+A_{10}`$ | $`22`$ |
| 56 | $`A_2+A_7+D_9`$ | 50 | $`2A_1+A_2+A_6+A_8`$ | $`23`$ |
| 57 | $`2A_2+D_5+D_9`$ | 210 | $`2A_2+D_{14}`$ | $`22`$ |
| 58 | $`A_2+D_7+D_9`$ | 186 | $`A_9+D_9`$ | $`21`$ |
| 59 | $`2A_2+A_4+D_{10}`$ | 72 | $`2A_1+A_2+A_4+A_{10}`$ | $`23`$ |
| 60 | $`A_3+A_4+D_{11}`$ | 44 | $`2A_1+2A_4+A_8`$ | $`23`$ |
| 61 | $`A_7+D_{11}`$ | 320 | $`D_{10}+E_8`$ | $`22`$ |
| 62 | $`A_2+D_5+D_{11}`$ | 186 | $`A_9+D_9`$ | $`21`$ |
| 63 | $`D_7+D_{11}`$ | 218 | $`D_{18}`$ | $`20`$ |
| 64 | $`A_2+A_4+D_{12}`$ | 72 | $`2A_1+A_2+A_4+A_{10}`$ | $`23`$ |
| 65 | $`A_6+D_{12}`$ | 320 | $`D_{10}+E_8`$ | $`22`$ |
| 66 | $`A_1+2A_2+D_{13}`$ | 90 | $`2A_1+2A_2+A_{12}`$ | $`23`$ |
| 67 | $`A_2+A_3+D_{13}`$ | 72 | $`2A_1+A_2+A_4+A_{10}`$ | $`23`$ |
| 68 | $`A_3+D_{15}`$ | 320 | $`D_{10}+E_8`$ | $`22`$ |
| 69 | $`A_2+D_{16}`$ | 320 | $`D_{10}+E_8`$ | $`22`$ |
| 70 | $`2A_1+A_4+2E_6`$ | 303 | $`A_1+A_4+A_5+E_8`$ | $`23`$ |
| 71 | $`2A_1+A_2+2A_4+E_6`$ | 23 | $`A_1+A_2+A_4+A_5+A_6`$ | $`23`$ |
| 72 | $`A_2+2A_3+A_4+E_6`$ | 46 | $`A_1+A_2+A_3+A_4+A_8`$ | $`23`$ |
| 73 | $`2A_6+E_6`$ | 37 | $`A_1+2A_2+A_6+A_7`$ | $`23`$ |
| 74 | $`2A_3+A_6+E_6`$ | 41 | $`A_5+A_6+A_7`$ | $`21`$ |
| 75 | $`A_2+A_3+A_7+E_6`$ | 37 | $`A_1+2A_2+A_6+A_7`$ | $`23`$ |
| 76 | $`2A_4+D_4+E_6`$ | 182 | $`A_4+A_5+D_9`$ | $`22`$ |
| 77 | $`A_2+A_6+D_4+E_6`$ | 183 | $`A_1+A_2+A_6+D_9`$ | $`23`$ |
| 78 | $`A_8+D_4+E_6`$ | 186 | $`A_9+D_9`$ | $`21`$ |
| 79 | $`A_1+D_5+2E_6`$ | 320 | $`D_{10}+E_8`$ | $`22`$ |
| 80 | $`A_2+2D_5+E_6`$ | 320 | $`D_{10}+E_8`$ | $`22`$ |
| 81 | $`A_1+A_2+A_4+D_5+E_6`$ | 193 | $`A_2+A_6+D_{10}`$ | $`22`$ |
| 82 | $`A_2+A_3+D_7+E_6`$ | 200 | $`A_2+A_5+D_{11}`$ | $`22`$ |
Table 1. List of embedding of $`\mathrm{\Delta }`$ in $`\mathrm{\Gamma }_f`$
| no | $`\mathrm{\Delta }`$ | No | $`\mathrm{\Sigma }_f`$ | $`\mathrm{𝑒𝑢}(\mathrm{\Sigma }_f)`$ |
| --- | --- | --- | --- | --- |
| 83 | $`A_5+D_7+E_6`$ | 320 | $`D_{10}+E_8`$ | $`22`$ |
| 84 | $`A_2+D_{10}+E_6`$ | 193 | $`A_2+A_6+D_{10}`$ | $`22`$ |
| 85 | $`A_1+A_2+2A_4+E_7`$ | 17 | $`2A_1+A_2+2A_4+A_6`$ | $`24`$ |
| 86 | $`A_3+2A_4+E_7`$ | 18 | $`A_1+A_3+2A_4+A_6`$ | $`23`$ |
| 87 | $`2A_2+D_7+E_7`$ | 210 | $`2A_2+D_{14}`$ | $`22`$ |
| 88 | $`A_2+2A_4+E_8`$ | 36 | $`A_2+A_4+A_5+A_7`$ | $`22`$ |
| 89 | $`2A_1+2A_2+A_4+E_8`$ | 30 | $`2A_2+A_3+A_4+A_7`$ | $`23`$ |
| 90 | $`2A_3+A_4+E_8`$ | 24 | $`A_3+A_4+A_5+A_6`$ | $`22`$ |
| 91 | $`A_3+A_7+E_8`$ | 46 | $`A_1+A_2+A_3+A_4+A_8`$ | $`23`$ |
| 92 | $`A_2+A_4+D_4+E_8`$ | 182 | $`A_4+A_5+D_9`$ | $`22`$ |
| 93 | $`A_6+D_4+E_8`$ | 186 | $`A_9+D_9`$ | $`21`$ |
| 94 | $`A_1+2A_2+D_5+E_8`$ | 210 | $`2A_2+D_{14}`$ | $`22`$ |
| 95 | $`A_2+A_3+D_5+E_8`$ | 198 | $`2A_2+A_3+D_{11}`$ | $`23`$ |
| 96 | $`A_3+D_7+E_8`$ | 213 | $`A_4+D_{14}`$ | $`21`$ |
| 97 | $`A_2+D_8+E_8`$ | 210 | $`2A_2+D_{14}`$ | $`22`$ |
| 98 | $`2A_1+A_2+E_6+E_8`$ | 320 | $`D_{10}+E_8`$ | $`22`$ |
Table 2. List of extremal elliptic $`K3`$ surfaces
| No | $`\mathrm{\Sigma }`$ | $`MW`$ | $`a`$ | $`b`$ | $`c`$ |
| --- | --- | --- | --- | --- | --- |
| 1 | $`6A_3`$ | $`/(4)\times /(4)`$ | $`4`$ | $`0`$ | $`4`$ |
| 2 | $`2A_1+4A_4`$ | $`/(5)`$ | $`10`$ | $`0`$ | $`10`$ |
| 3 | $`2A_2+2A_3+2A_4`$ | $`(0)`$ | $`60`$ | $`0`$ | $`60`$ |
| 4 | $`3A_1+3A_5`$ | $`/(2)\times /(6)`$ | $`2`$ | $`0`$ | $`6`$ |
| 5 | $`4A_2+2A_5`$ | $`/(3)\times /(3)`$ | $`6`$ | $`0`$ | $`6`$ |
| 6 | $`A_3+3A_5`$ | $`/(6)`$ | $`4`$ | $`0`$ | $`6`$ |
| 7 | $`2A_1+2A_3+2A_5`$ | $`/(2)\times /(2)`$ | $`12`$ | $`0`$ | $`12`$ |
| 8 | $`A_1+2A_2+A_3+2A_5`$ | $`/(6)`$ | $`6`$ | $`0`$ | $`12`$ |
| 9 | $`2A_4+2A_5`$ | $`(0)`$ | $`30`$ | $`0`$ | $`30`$ |
| 10 | $`2A_2+A_4+2A_5`$ | $`/(3)`$ | $`6`$ | $`0`$ | $`30`$ |
| 11 | $`A_1+A_3+A_4+2A_5`$ | $`/(2)`$ | $`12`$ | $`0`$ | $`30`$ |
| 12 | $`A_1+A_2+2A_3+A_4+A_5`$ | $`/(2)`$ | $`24`$ | $`12`$ | $`36`$ |
| 13 | $`3A_6`$ | $`/(7)`$ | $`2`$ | $`1`$ | $`4`$ |
| 14 | $`2A_1+2A_2+2A_6`$ | $`(0)`$ | $`42`$ | $`0`$ | $`42`$ |
| 15 | $`2A_3+2A_6`$ | $`(0)`$ | $`28`$ | $`0`$ | $`28`$ |
| 16 | $`A_2+A_4+2A_6`$ | $`(0)`$ | $`28`$ | $`7`$ | $`28`$ |
| 17 | $`2A_1+A_2+2A_4+A_6`$ | $`(0)`$ | $`50`$ | $`20`$ | $`50`$ |
| 18 | $`A_1+A_3+2A_4+A_6`$ | $`(0)`$ | $`10`$ | $`0`$ | $`140`$ |
| | | | $`20`$ | $`0`$ | $`70`$ |
| 19 | $`A_2+2A_3+A_4+A_6`$ | $`(0)`$ | $`24`$ | $`12`$ | $`76`$ |
| 20 | $`A_1+2A_2+A_3+A_4+A_6`$ | $`(0)`$ | $`30`$ | $`0`$ | $`84`$ |
| 21 | $`2A_1+2A_5+A_6`$ | $`/(2)`$ | $`12`$ | $`6`$ | $`24`$ |
| 22 | $`A_1+2A_3+A_5+A_6`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`84`$ |
| 23 | $`A_1+A_2+A_4+A_5+A_6`$ | $`(0)`$ | $`30`$ | $`0`$ | $`42`$ |
| | | | $`18`$ | $`6`$ | $`72`$ |
| 24 | $`A_3+A_4+A_5+A_6`$ | $`(0)`$ | $`12`$ | $`0`$ | $`70`$ |
| 25 | $`4A_1+2A_7`$ | $`/(2)\times /(4)`$ | $`4`$ | $`0`$ | $`4`$ |
| 26 | $`2A_2+2A_7`$ | $`(0)`$ | $`24`$ | $`0`$ | $`24`$ |
| | | $`/(2)`$ | $`12`$ | $`0`$ | $`12`$ |
| 27 | $`A_1+A_3+2A_7`$ | $`/(8)`$ | $`2`$ | $`0`$ | $`4`$ |
| 28 | $`2A_1+3A_3+A_7`$ | $`/(2)\times /(4)`$ | $`4`$ | $`0`$ | $`8`$ |
| 29 | $`A_2+3A_3+A_7`$ | $`/(4)`$ | $`4`$ | $`0`$ | $`24`$ |
| 30 | $`2A_2+A_3+A_4+A_7`$ | $`(0)`$ | $`12`$ | $`0`$ | $`120`$ |
| 31 | $`2A_1+A_2+A_3+A_4+A_7`$ | $`/(2)`$ | $`20`$ | $`0`$ | $`24`$ |
| 32 | $`A_1+2A_5+A_7`$ | $`/(2)`$ | $`6`$ | $`0`$ | $`24`$ |
| 33 | $`3A_1+A_3+A_5+A_7`$ | $`/(2)\times /(2)`$ | $`8`$ | $`0`$ | $`12`$ |
Table 2. List of extremal elliptic $`K3`$ surfaces
| No | $`\mathrm{\Sigma }`$ | $`MW`$ | $`a`$ | $`b`$ | $`c`$ |
| --- | --- | --- | --- | --- | --- |
| 34 | $`A_1+A_2+A_3+A_5+A_7`$ | $`/(2)`$ | $`12`$ | $`0`$ | $`24`$ |
| 35 | $`2A_1+A_4+A_5+A_7`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`120`$ |
| 36 | $`A_2+A_4+A_5+A_7`$ | $`(0)`$ | $`6`$ | $`0`$ | $`120`$ |
| | | | $`24`$ | $`0`$ | $`30`$ |
| 37 | $`A_1+2A_2+A_6+A_7`$ | $`(0)`$ | $`24`$ | $`0`$ | $`42`$ |
| 38 | $`2A_1+A_3+A_6+A_7`$ | $`/(2)`$ | $`12`$ | $`4`$ | $`20`$ |
| 39 | $`A_2+A_3+A_6+A_7`$ | $`(0)`$ | $`4`$ | $`0`$ | $`168`$ |
| 40 | $`A_1+A_4+A_6+A_7`$ | $`(0)`$ | $`2`$ | $`0`$ | $`280`$ |
| | | | $`18`$ | $`4`$ | $`32`$ |
| 41 | $`A_5+A_6+A_7`$ | $`(0)`$ | $`16`$ | $`4`$ | $`22`$ |
| 42 | $`2A_1+2A_8`$ | $`(0)`$ | $`18`$ | $`0`$ | $`18`$ |
| | | $`/(3)`$ | $`4`$ | $`2`$ | $`10`$ |
| 43 | $`A_1+3A_2+A_3+A_8`$ | $`/(3)`$ | $`12`$ | $`0`$ | $`18`$ |
| 44 | $`2A_1+2A_4+A_8`$ | $`(0)`$ | $`20`$ | $`10`$ | $`50`$ |
| 45 | $`3A_2+A_4+A_8`$ | $`/(3)`$ | $`12`$ | $`3`$ | $`12`$ |
| 46 | $`A_1+A_2+A_3+A_4+A_8`$ | $`(0)`$ | $`6`$ | $`0`$ | $`180`$ |
| 47 | $`A_1+2A_2+A_5+A_8`$ | $`/(3)`$ | $`6`$ | $`0`$ | $`18`$ |
| 48 | $`A_2+A_3+A_5+A_8`$ | $`/(3)`$ | $`4`$ | $`0`$ | $`18`$ |
| 49 | $`A_1+A_4+A_5+A_8`$ | $`(0)`$ | $`18`$ | $`0`$ | $`30`$ |
| 50 | $`2A_1+A_2+A_6+A_8`$ | $`(0)`$ | $`18`$ | $`0`$ | $`42`$ |
| 51 | $`A_1+A_3+A_6+A_8`$ | $`(0)`$ | $`10`$ | $`4`$ | $`52`$ |
| 52 | $`A_4+A_6+A_8`$ | $`(0)`$ | $`18`$ | $`9`$ | $`22`$ |
| 53 | $`A_1+A_2+A_7+A_8`$ | $`(0)`$ | $`18`$ | $`0`$ | $`24`$ |
| 54 | $`2A_9`$ | $`(0)`$ | $`10`$ | $`0`$ | $`10`$ |
| | | $`/(5)`$ | $`2`$ | $`0`$ | $`2`$ |
| 55 | $`A_1+A_2+2A_3+A_9`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`60`$ |
| 56 | $`2A_1+2A_2+A_3+A_9`$ | $`/(2)`$ | $`6`$ | $`0`$ | $`60`$ |
| 57 | $`A_1+2A_4+A_9`$ | $`/(5)`$ | $`2`$ | $`0`$ | $`10`$ |
| 58 | $`3A_1+A_2+A_4+A_9`$ | $`/(2)`$ | $`20`$ | $`10`$ | $`20`$ |
| 59 | $`2A_1+A_3+A_4+A_9`$ | $`/(2)`$ | $`10`$ | $`0`$ | $`20`$ |
| 60 | $`2A_1+A_2+A_5+A_9`$ | $`/(2)`$ | $`12`$ | $`6`$ | $`18`$ |
| 61 | $`A_1+A_3+A_5+A_9`$ | $`/(2)`$ | $`10`$ | $`0`$ | $`12`$ |
| 62 | $`A_4+A_5+A_9`$ | $`(0)`$ | $`10`$ | $`0`$ | $`30`$ |
| | | $`/(2)`$ | $`10`$ | $`5`$ | $`10`$ |
| 63 | $`3A_1+A_6+A_9`$ | $`/(2)`$ | $`4`$ | $`2`$ | $`36`$ |
| 64 | $`A_1+A_2+A_6+A_9`$ | $`(0)`$ | $`10`$ | $`0`$ | $`42`$ |
Table 2. List of extremal elliptic $`K3`$ surfaces
| No | $`\mathrm{\Sigma }`$ | $`MW`$ | $`a`$ | $`b`$ | $`c`$ |
| --- | --- | --- | --- | --- | --- |
| 65 | $`A_3+A_6+A_9`$ | $`(0)`$ | $`2`$ | $`0`$ | $`140`$ |
| 66 | $`A_2+A_7+A_9`$ | $`(0)`$ | $`10`$ | $`0`$ | $`24`$ |
| 67 | $`A_1+A_8+A_9`$ | $`(0)`$ | $`10`$ | $`0`$ | $`18`$ |
| 68 | $`A_2+2A_3+A_{10}`$ | $`(0)`$ | $`24`$ | $`12`$ | $`28`$ |
| 69 | $`A_1+2A_2+A_3+A_{10}`$ | $`(0)`$ | $`12`$ | $`0`$ | $`66`$ |
| 70 | $`2A_4+A_{10}`$ | $`(0)`$ | $`10`$ | $`5`$ | $`30`$ |
| 71 | $`2A_2+A_4+A_{10}`$ | $`(0)`$ | $`6`$ | $`3`$ | $`84`$ |
| | | | $`24`$ | $`9`$ | $`24`$ |
| 72 | $`2A_1+A_2+A_4+A_{10}`$ | $`(0)`$ | $`2`$ | $`0`$ | $`330`$ |
| 73 | $`A_1+A_3+A_4+A_{10}`$ | $`(0)`$ | $`20`$ | $`0`$ | $`22`$ |
| | | | $`12`$ | $`4`$ | $`38`$ |
| 74 | $`A_1+A_2+A_5+A_{10}`$ | $`(0)`$ | $`6`$ | $`0`$ | $`66`$ |
| | | | $`18`$ | $`6`$ | $`24`$ |
| 75 | $`A_3+A_5+A_{10}`$ | $`(0)`$ | $`4`$ | $`0`$ | $`66`$ |
| | | | $`12`$ | $`0`$ | $`22`$ |
| 76 | $`2A_1+A_6+A_{10}`$ | $`(0)`$ | $`12`$ | $`2`$ | $`26`$ |
| 77 | $`A_2+A_6+A_{10}`$ | $`(0)`$ | $`4`$ | $`1`$ | $`58`$ |
| | | | $`16`$ | $`5`$ | $`16`$ |
| 78 | $`A_1+A_7+A_{10}`$ | $`(0)`$ | $`2`$ | $`0`$ | $`88`$ |
| | | | $`10`$ | $`2`$ | $`18`$ |
| 79 | $`A_8+A_{10}`$ | $`(0)`$ | $`10`$ | $`1`$ | $`10`$ |
| 80 | $`A_1+3A_2+A_{11}`$ | $`/(3)`$ | $`6`$ | $`0`$ | $`12`$ |
| 81 | $`3A_1+2A_2+A_{11}`$ | $`/(6)`$ | $`2`$ | $`0`$ | $`12`$ |
| 82 | $`A_1+2A_3+A_{11}`$ | $`/(4)`$ | $`4`$ | $`0`$ | $`6`$ |
| 83 | $`2A_2+A_3+A_{11}`$ | $`/(3)`$ | $`4`$ | $`0`$ | $`12`$ |
| | | $`/(6)`$ | $`4`$ | $`2`$ | $`4`$ |
| 84 | $`2A_1+A_2+A_3+A_{11}`$ | $`/(4)`$ | $`6`$ | $`0`$ | $`6`$ |
| | | $`/(2)`$ | $`12`$ | $`0`$ | $`12`$ |
| 85 | $`3A_1+A_4+A_{11}`$ | $`/(2)`$ | $`6`$ | $`0`$ | $`20`$ |
| 86 | $`A_1+A_2+A_4+A_{11}`$ | $`(0)`$ | $`12`$ | $`0`$ | $`30`$ |
| 87 | $`2A_1+A_5+A_{11}`$ | $`/(2)`$ | $`6`$ | $`0`$ | $`12`$ |
| | | $`/(6)`$ | $`2`$ | $`0`$ | $`4`$ |
| 88 | $`A_2+A_5+A_{11}`$ | $`/(3)`$ | $`4`$ | $`0`$ | $`6`$ |
| 89 | $`A_1+A_6+A_{11}`$ | $`(0)`$ | $`4`$ | $`0`$ | $`42`$ |
| 90 | $`2A_1+2A_2+A_{12}`$ | $`(0)`$ | $`12`$ | $`6`$ | $`42`$ |
| 91 | $`A_1+A_2+A_3+A_{12}`$ | $`(0)`$ | $`6`$ | $`0`$ | $`52`$ |
Table 2. List of extremal elliptic $`K3`$ surfaces
| No | $`\mathrm{\Sigma }`$ | $`MW`$ | $`a`$ | $`b`$ | $`c`$ |
| --- | --- | --- | --- | --- | --- |
| 92 | $`2A_1+A_4+A_{12}`$ | $`(0)`$ | $`2`$ | $`0`$ | $`130`$ |
| | | | $`18`$ | $`8`$ | $`18`$ |
| 93 | $`A_2+A_4+A_{12}`$ | $`(0)`$ | $`6`$ | $`3`$ | $`34`$ |
| 94 | $`A_1+A_5+A_{12}`$ | $`(0)`$ | $`10`$ | $`2`$ | $`16`$ |
| 95 | $`A_6+A_{12}`$ | $`(0)`$ | $`2`$ | $`1`$ | $`46`$ |
| 96 | $`A_1+2A_2+A_{13}`$ | $`(0)`$ | $`6`$ | $`0`$ | $`42`$ |
| | | $`/(2)`$ | $`6`$ | $`3`$ | $`12`$ |
| 97 | $`3A_1+A_2+A_{13}`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`42`$ |
| 98 | $`2A_1+A_3+A_{13}`$ | $`/(2)`$ | $`6`$ | $`2`$ | $`10`$ |
| 99 | $`A_2+A_3+A_{13}`$ | $`(0)`$ | $`4`$ | $`0`$ | $`42`$ |
| 100 | $`A_1+A_4+A_{13}`$ | $`(0)`$ | $`2`$ | $`0`$ | $`70`$ |
| | | | $`8`$ | $`2`$ | $`18`$ |
| | | $`/(2)`$ | $`2`$ | $`1`$ | $`18`$ |
| 101 | $`A_5+A_{13}`$ | $`(0)`$ | $`4`$ | $`2`$ | $`22`$ |
| 102 | $`2A_2+A_{14}`$ | $`/(3)`$ | $`4`$ | $`1`$ | $`4`$ |
| 103 | $`2A_1+A_2+A_{14}`$ | $`(0)`$ | $`12`$ | $`6`$ | $`18`$ |
| | | $`/(3)`$ | $`2`$ | $`0`$ | $`10`$ |
| 104 | $`A_1+A_3+A_{14}`$ | $`(0)`$ | $`10`$ | $`0`$ | $`12`$ |
| 105 | $`A_4+A_{14}`$ | $`(0)`$ | $`10`$ | $`5`$ | $`10`$ |
| 106 | $`3A_1+A_{15}`$ | $`/(4)`$ | $`2`$ | $`0`$ | $`4`$ |
| 107 | $`A_1+A_2+A_{15}`$ | $`(0)`$ | $`10`$ | $`2`$ | $`10`$ |
| | | $`/(2)`$ | $`4`$ | $`0`$ | $`6`$ |
| 108 | $`A_3+A_{15}`$ | $`/(4)`$ | $`2`$ | $`0`$ | $`2`$ |
| 109 | $`2A_1+A_{16}`$ | $`(0)`$ | $`2`$ | $`0`$ | $`34`$ |
| | | | $`4`$ | $`2`$ | $`18`$ |
| 110 | $`A_2+A_{16}`$ | $`(0)`$ | $`6`$ | $`3`$ | $`10`$ |
| 111 | $`A_1+A_{17}`$ | $`(0)`$ | $`4`$ | $`2`$ | $`10`$ |
| | | $`/(3)`$ | $`2`$ | $`0`$ | $`2`$ |
| 112 | $`A_{18}`$ | $`(0)`$ | $`2`$ | $`1`$ | $`10`$ |
| 113 | $`2A_4+2D_5`$ | $`(0)`$ | $`20`$ | $`0`$ | $`20`$ |
| 114 | $`A_3+2A_5+D_5`$ | $`/(2)`$ | $`12`$ | $`0`$ | $`12`$ |
| 115 | $`2A_4+A_5+D_5`$ | $`(0)`$ | $`20`$ | $`0`$ | $`30`$ |
| 116 | $`A_1+A_3+A_4+A_5+D_5`$ | $`/(2)`$ | $`12`$ | $`0`$ | $`20`$ |
| 117 | $`A_1+2A_6+D_5`$ | $`(0)`$ | $`14`$ | $`0`$ | $`28`$ |
| 118 | $`2A_2+A_3+A_6+D_5`$ | $`(0)`$ | $`12`$ | $`0`$ | $`84`$ |
| 119 | $`A_1+A_2+A_4+A_6+D_5`$ | $`(0)`$ | $`20`$ | $`0`$ | $`42`$ |
Table 2. List of extremal elliptic $`K3`$ surfaces
| No | $`\mathrm{\Sigma }`$ | $`MW`$ | $`a`$ | $`b`$ | $`c`$ |
| --- | --- | --- | --- | --- | --- |
| 120 | $`A_2+A_5+A_6+D_5`$ | $`(0)`$ | $`6`$ | $`0`$ | $`84`$ |
| | | | $`12`$ | $`0`$ | $`42`$ |
| 121 | $`A_1+A_7+2D_5`$ | $`/(4)`$ | $`2`$ | $`0`$ | $`8`$ |
| 122 | $`A_1+A_2+A_3+A_7+D_5`$ | $`/(4)`$ | $`6`$ | $`0`$ | $`8`$ |
| 123 | $`2A_1+A_4+A_7+D_5`$ | $`/(2)`$ | $`8`$ | $`0`$ | $`20`$ |
| 124 | $`A_8+2D_5`$ | $`(0)`$ | $`8`$ | $`4`$ | $`20`$ |
| 125 | $`A_1+A_4+A_8+D_5`$ | $`(0)`$ | $`2`$ | $`0`$ | $`180`$ |
| | | | $`18`$ | $`0`$ | $`20`$ |
| 126 | $`A_5+A_8+D_5`$ | $`(0)`$ | $`12`$ | $`0`$ | $`18`$ |
| 127 | $`2A_2+A_9+D_5`$ | $`(0)`$ | $`6`$ | $`0`$ | $`60`$ |
| 128 | $`2A_1+A_2+A_9+D_5`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`60`$ |
| 129 | $`A_1+A_3+A_9+D_5`$ | $`/(2)`$ | $`8`$ | $`4`$ | $`12`$ |
| 130 | $`A_4+A_9+D_5`$ | $`(0)`$ | $`10`$ | $`0`$ | $`20`$ |
| 131 | $`A_1+A_2+A_{10}+D_5`$ | $`(0)`$ | $`14`$ | $`4`$ | $`20`$ |
| 132 | $`2A_1+A_{11}+D_5`$ | $`/(4)`$ | $`2`$ | $`0`$ | $`6`$ |
| 133 | $`A_2+A_{11}+D_5`$ | $`/(2)`$ | $`6`$ | $`0`$ | $`6`$ |
| 134 | $`A_1+A_{12}+D_5`$ | $`(0)`$ | $`2`$ | $`0`$ | $`52`$ |
| | | | $`6`$ | $`2`$ | $`18`$ |
| 135 | $`A_{13}+D_5`$ | $`(0)`$ | $`6`$ | $`2`$ | $`10`$ |
| 136 | $`3D_6`$ | $`/(2)\times /(2)`$ | $`2`$ | $`0`$ | $`2`$ |
| 137 | $`2A_3+2D_6`$ | $`/(2)\times /(2)`$ | $`4`$ | $`0`$ | $`4`$ |
| 138 | $`2A_2+2A_4+D_6`$ | $`(0)`$ | $`30`$ | $`0`$ | $`30`$ |
| 139 | $`2A_1+2A_5+D_6`$ | $`/(2)\times /(2)`$ | $`6`$ | $`0`$ | $`6`$ |
| 140 | $`A_1+2A_3+A_5+D_6`$ | $`/(2)\times /(2)`$ | $`4`$ | $`0`$ | $`12`$ |
| 141 | $`A_3+A_4+A_5+D_6`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`30`$ |
| 142 | $`2A_6+D_6`$ | $`(0)`$ | $`14`$ | $`0`$ | $`14`$ |
| 143 | $`A_2+A_4+A_6+D_6`$ | $`(0)`$ | $`6`$ | $`0`$ | $`70`$ |
| 144 | $`A_1+2A_2+A_7+D_6`$ | $`/(2)`$ | $`6`$ | $`0`$ | $`24`$ |
| 145 | $`A_2+A_3+A_7+D_6`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`24`$ |
| 146 | $`A_1+A_4+A_7+D_6`$ | $`/(2)`$ | $`6`$ | $`2`$ | $`14`$ |
| 147 | $`A_4+A_8+D_6`$ | $`(0)`$ | $`4`$ | $`2`$ | $`46`$ |
| 148 | $`A_1+A_2+A_9+D_6`$ | $`/(2)`$ | $`6`$ | $`0`$ | $`10`$ |
| | | $`/(2)`$ | $`4`$ | $`2`$ | $`16`$ |
| 149 | $`A_3+A_9+D_6`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`10`$ |
| 150 | $`A_2+A_{10}+D_6`$ | $`(0)`$ | $`6`$ | $`0`$ | $`22`$ |
| 151 | $`A_1+A_{11}+D_6`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`6`$ |
Table 2. List of extremal elliptic $`K3`$ surfaces
| No | $`\mathrm{\Sigma }`$ | $`MW`$ | $`a`$ | $`b`$ | $`c`$ |
| --- | --- | --- | --- | --- | --- |
| 152 | $`A_{12}+D_6`$ | $`(0)`$ | $`4`$ | $`2`$ | $`14`$ |
| 153 | $`A_2+A_5+D_5+D_6`$ | $`/(2)`$ | $`6`$ | $`0`$ | $`12`$ |
| 154 | $`A_7+D_5+D_6`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`8`$ |
| 155 | $`2A_2+2D_7`$ | $`(0)`$ | $`12`$ | $`0`$ | $`12`$ |
| 156 | $`A_2+3A_3+D_7`$ | $`/(4)`$ | $`8`$ | $`4`$ | $`8`$ |
| 157 | $`A_1+A_2+2A_4+D_7`$ | $`(0)`$ | $`10`$ | $`0`$ | $`60`$ |
| 158 | $`A_2+A_3+A_6+D_7`$ | $`(0)`$ | $`8`$ | $`4`$ | $`44`$ |
| 159 | $`A_1+A_4+A_6+D_7`$ | $`(0)`$ | $`4`$ | $`0`$ | $`70`$ |
| 160 | $`A_5+A_6+D_7`$ | $`(0)`$ | $`2`$ | $`0`$ | $`84`$ |
| 161 | $`2A_1+A_2+A_7+D_7`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`24`$ |
| 162 | $`A_1+A_3+A_7+D_7`$ | $`/(4)`$ | $`2`$ | $`0`$ | $`8`$ |
| 163 | $`2A_1+A_9+D_7`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`10`$ |
| 164 | $`A_2+A_9+D_7`$ | $`(0)`$ | $`2`$ | $`0`$ | $`60`$ |
| 165 | $`A_1+A_{10}+D_7`$ | $`(0)`$ | $`4`$ | $`0`$ | $`22`$ |
| 166 | $`A_{11}+D_7`$ | $`/(4)`$ | $`2`$ | $`1`$ | $`2`$ |
| 167 | $`A_1+A_5+D_5+D_7`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`12`$ |
| 168 | $`A_5+D_6+D_7`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`12`$ |
| 169 | $`2A_1+2D_8`$ | $`/(2)\times /(2)`$ | $`2`$ | $`0`$ | $`2`$ |
| 170 | $`2A_2+2A_3+D_8`$ | $`/(2)`$ | $`12`$ | $`0`$ | $`12`$ |
| 171 | $`2A_5+D_8`$ | $`/(2)`$ | $`6`$ | $`0`$ | $`6`$ |
| 172 | $`2A_1+A_3+A_5+D_8`$ | $`/(2)\times /(2)`$ | $`2`$ | $`0`$ | $`12`$ |
| 173 | $`A_1+A_4+A_5+D_8`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`30`$ |
| 174 | $`2A_2+A_6+D_8`$ | $`(0)`$ | $`12`$ | $`6`$ | $`24`$ |
| 175 | $`A_1+A_2+A_7+D_8`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`24`$ |
| 176 | $`A_1+A_9+D_8`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`10`$ |
| 177 | $`2D_5+D_8`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`4`$ |
| 178 | $`A_1+A_3+D_6+D_8`$ | $`/(2)\times /(2)`$ | $`2`$ | $`0`$ | $`4`$ |
| 179 | $`2D_9`$ | $`(0)`$ | $`4`$ | $`0`$ | $`4`$ |
| 180 | $`A_1+2A_2+A_4+D_9`$ | $`(0)`$ | $`12`$ | $`0`$ | $`30`$ |
| 181 | $`A_1+A_3+A_5+D_9`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`12`$ |
| 182 | $`A_4+A_5+D_9`$ | $`(0)`$ | $`4`$ | $`0`$ | $`30`$ |
| 183 | $`A_1+A_2+A_6+D_9`$ | $`(0)`$ | $`4`$ | $`0`$ | $`42`$ |
| 184 | $`2A_1+A_7+D_9`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`8`$ |
| 185 | $`A_1+A_8+D_9`$ | $`(0)`$ | $`4`$ | $`0`$ | $`18`$ |
| 186 | $`A_9+D_9`$ | $`(0)`$ | $`4`$ | $`0`$ | $`10`$ |
| 187 | $`A_4+D_5+D_9`$ | $`(0)`$ | $`4`$ | $`0`$ | $`20`$ |
Table 2. List of extremal elliptic $`K3`$ surfaces
| No | $`\mathrm{\Sigma }`$ | $`MW`$ | $`a`$ | $`b`$ | $`c`$ |
| --- | --- | --- | --- | --- | --- |
| 188 | $`2A_1+2A_3+D_{10}`$ | $`/(2)\times /(2)`$ | $`4`$ | $`0`$ | $`4`$ |
| 189 | $`2A_4+D_{10}`$ | $`(0)`$ | $`10`$ | $`0`$ | $`10`$ |
| 190 | $`A_1+A_3+A_4+D_{10}`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`20`$ |
| 191 | $`3A_1+A_5+D_{10}`$ | $`/(2)\times /(2)`$ | $`4`$ | $`2`$ | $`4`$ |
| 192 | $`A_3+A_5+D_{10}`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`12`$ |
| 193 | $`A_2+A_6+D_{10}`$ | $`(0)`$ | $`2`$ | $`0`$ | $`42`$ |
| 194 | $`A_8+D_{10}`$ | $`(0)`$ | $`2`$ | $`0`$ | $`18`$ |
| 195 | $`A_1+A_2+D_5+D_{10}`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`6`$ |
| 196 | $`A_2+D_6+D_{10}`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`6`$ |
| 197 | $`A_1+D_7+D_{10}`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`4`$ |
| 198 | $`2A_2+A_3+D_{11}`$ | $`(0)`$ | $`12`$ | $`0`$ | $`12`$ |
| 199 | $`A_1+A_2+A_4+D_{11}`$ | $`(0)`$ | $`6`$ | $`0`$ | $`20`$ |
| 200 | $`A_2+A_5+D_{11}`$ | $`(0)`$ | $`6`$ | $`0`$ | $`12`$ |
| 201 | $`A_1+A_6+D_{11}`$ | $`(0)`$ | $`6`$ | $`2`$ | $`10`$ |
| 202 | $`2A_1+2A_2+D_{12}`$ | $`/(2)`$ | $`6`$ | $`0`$ | $`6`$ |
| 203 | $`A_1+A_2+A_3+D_{12}`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`6`$ |
| 204 | $`2A_1+A_4+D_{12}`$ | $`/(2)`$ | $`4`$ | $`2`$ | $`6`$ |
| 205 | $`A_1+D_5+D_{12}`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`4`$ |
| 206 | $`D_6+D_{12}`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`2`$ |
| 207 | $`A_1+A_4+D_{13}`$ | $`(0)`$ | $`2`$ | $`0`$ | $`20`$ |
| 208 | $`A_5+D_{13}`$ | $`(0)`$ | $`2`$ | $`0`$ | $`12`$ |
| 209 | $`D_5+D_{13}`$ | $`(0)`$ | $`4`$ | $`0`$ | $`4`$ |
| 210 | $`2A_2+D_{14}`$ | $`(0)`$ | $`6`$ | $`0`$ | $`6`$ |
| 211 | $`2A_1+A_2+D_{14}`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`6`$ |
| 212 | $`A_1+A_3+D_{14}`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`4`$ |
| 213 | $`A_4+D_{14}`$ | $`(0)`$ | $`4`$ | $`2`$ | $`6`$ |
| 214 | $`A_1+A_2+D_{15}`$ | $`(0)`$ | $`4`$ | $`0`$ | $`6`$ |
| 215 | $`2A_1+D_{16}`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`2`$ |
| 216 | $`A_2+D_{16}`$ | $`/(2)`$ | $`2`$ | $`1`$ | $`2`$ |
| 217 | $`A_1+D_{17}`$ | $`(0)`$ | $`2`$ | $`0`$ | $`4`$ |
| 218 | $`D_{18}`$ | $`(0)`$ | $`2`$ | $`0`$ | $`2`$ |
| 219 | $`3E_6`$ | $`/(3)`$ | $`2`$ | $`1`$ | $`2`$ |
| 220 | $`2A_3+2E_6`$ | $`(0)`$ | $`12`$ | $`0`$ | $`12`$ |
| 221 | $`A_1+A_3+2A_4+E_6`$ | $`(0)`$ | $`20`$ | $`0`$ | $`30`$ |
| 222 | $`A_1+A_5+2E_6`$ | $`/(3)`$ | $`2`$ | $`0`$ | $`6`$ |
| 223 | $`A_2+2A_5+E_6`$ | $`/(3)`$ | $`6`$ | $`0`$ | $`6`$ |
Table 2. List of extremal elliptic $`K3`$ surfaces
| No | $`\mathrm{\Sigma }`$ | $`MW`$ | $`a`$ | $`b`$ | $`c`$ |
| --- | --- | --- | --- | --- | --- |
| 224 | $`2A_2+A_3+A_5+E_6`$ | $`/(3)`$ | $`6`$ | $`0`$ | $`12`$ |
| 225 | $`A_3+A_4+A_5+E_6`$ | $`(0)`$ | $`12`$ | $`0`$ | $`30`$ |
| 226 | $`A_6+2E_6`$ | $`(0)`$ | $`6`$ | $`3`$ | $`12`$ |
| 227 | $`A_1+A_2+A_3+A_6+E_6`$ | $`(0)`$ | $`6`$ | $`0`$ | $`84`$ |
| | | | $`12`$ | $`0`$ | $`42`$ |
| 228 | $`2A_1+A_4+A_6+E_6`$ | $`(0)`$ | $`20`$ | $`10`$ | $`26`$ |
| 229 | $`A_2+A_4+A_6+E_6`$ | $`(0)`$ | $`18`$ | $`3`$ | $`18`$ |
| 230 | $`A_1+A_5+A_6+E_6`$ | $`(0)`$ | $`6`$ | $`0`$ | $`42`$ |
| 231 | $`A_1+A_4+A_7+E_6`$ | $`(0)`$ | $`2`$ | $`0`$ | $`120`$ |
| 232 | $`A_5+A_7+E_6`$ | $`(0)`$ | $`6`$ | $`0`$ | $`24`$ |
| 233 | $`2A_2+A_8+E_6`$ | $`/(3)`$ | $`6`$ | $`3`$ | $`6`$ |
| 234 | $`2A_1+A_2+A_8+E_6`$ | $`/(3)`$ | $`2`$ | $`0`$ | $`18`$ |
| 235 | $`A_1+A_3+A_8+E_6`$ | $`(0)`$ | $`12`$ | $`0`$ | $`18`$ |
| 236 | $`A_4+A_8+E_6`$ | $`(0)`$ | $`12`$ | $`3`$ | $`12`$ |
| 237 | $`A_1+A_2+A_9+E_6`$ | $`(0)`$ | $`12`$ | $`6`$ | $`18`$ |
| 238 | $`A_3+A_9+E_6`$ | $`(0)`$ | $`10`$ | $`0`$ | $`12`$ |
| 239 | $`2A_1+A_{10}+E_6`$ | $`(0)`$ | $`2`$ | $`0`$ | $`66`$ |
| 240 | $`A_2+A_{10}+E_6`$ | $`(0)`$ | $`6`$ | $`3`$ | $`18`$ |
| 241 | $`A_1+A_{11}+E_6`$ | $`(0)`$ | $`6`$ | $`0`$ | $`12`$ |
| | | $`/(3)`$ | $`2`$ | $`0`$ | $`4`$ |
| 242 | $`A_{12}+E_6`$ | $`(0)`$ | $`4`$ | $`1`$ | $`10`$ |
| 243 | $`A_3+A_4+D_5+E_6`$ | $`(0)`$ | $`12`$ | $`0`$ | $`20`$ |
| 244 | $`A_1+A_6+D_5+E_6`$ | $`(0)`$ | $`2`$ | $`0`$ | $`84`$ |
| 245 | $`A_7+D_5+E_6`$ | $`(0)`$ | $`8`$ | $`0`$ | $`12`$ |
| 246 | $`D_6+2E_6`$ | $`(0)`$ | $`6`$ | $`0`$ | $`6`$ |
| 247 | $`A_2+A_4+D_6+E_6`$ | $`(0)`$ | $`6`$ | $`0`$ | $`30`$ |
| 248 | $`A_6+D_6+E_6`$ | $`(0)`$ | $`4`$ | $`2`$ | $`22`$ |
| 249 | $`A_1+A_4+D_7+E_6`$ | $`(0)`$ | $`4`$ | $`0`$ | $`30`$ |
| 250 | $`D_5+D_7+E_6`$ | $`(0)`$ | $`4`$ | $`0`$ | $`12`$ |
| 251 | $`A_4+D_8+E_6`$ | $`(0)`$ | $`8`$ | $`2`$ | $`8`$ |
| 252 | $`A_1+A_2+D_9+E_6`$ | $`(0)`$ | $`6`$ | $`0`$ | $`12`$ |
| 253 | $`A_3+D_9+E_6`$ | $`(0)`$ | $`4`$ | $`0`$ | $`12`$ |
| 254 | $`A_1+D_{11}+E_6`$ | $`(0)`$ | $`2`$ | $`0`$ | $`12`$ |
| 255 | $`D_{12}+E_6`$ | $`(0)`$ | $`4`$ | $`2`$ | $`4`$ |
| 256 | $`2A_2+2E_7`$ | $`(0)`$ | $`6`$ | $`0`$ | $`6`$ |
| 257 | $`A_1+A_3+2E_7`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`4`$ |
Table 2. List of extremal elliptic $`K3`$ surfaces
| No | $`\mathrm{\Sigma }`$ | $`MW`$ | $`a`$ | $`b`$ | $`c`$ |
| --- | --- | --- | --- | --- | --- |
| 258 | $`A_4+2E_7`$ | $`(0)`$ | $`4`$ | $`2`$ | $`6`$ |
| 259 | $`A_1+2A_3+A_4+E_7`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`20`$ |
| 260 | $`2A_2+A_3+A_4+E_7`$ | $`(0)`$ | $`12`$ | $`0`$ | $`30`$ |
| 261 | $`2A_3+A_5+E_7`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`12`$ |
| 262 | $`A_1+A_2+A_3+A_5+E_7`$ | $`/(2)`$ | $`6`$ | $`0`$ | $`12`$ |
| 263 | $`2A_1+A_4+A_5+E_7`$ | $`/(2)`$ | $`8`$ | $`2`$ | $`8`$ |
| 264 | $`A_2+A_4+A_5+E_7`$ | $`(0)`$ | $`6`$ | $`0`$ | $`30`$ |
| 265 | $`A_1+2A_2+A_6+E_7`$ | $`(0)`$ | $`6`$ | $`0`$ | $`42`$ |
| 266 | $`A_2+A_3+A_6+E_7`$ | $`(0)`$ | $`4`$ | $`0`$ | $`42`$ |
| 267 | $`A_1+A_4+A_6+E_7`$ | $`(0)`$ | $`2`$ | $`0`$ | $`70`$ |
| | | | $`8`$ | $`2`$ | $`18`$ |
| 268 | $`A_5+A_6+E_7`$ | $`(0)`$ | $`4`$ | $`2`$ | $`22`$ |
| 269 | $`2A_2+A_7+E_7`$ | $`(0)`$ | $`6`$ | $`0`$ | $`24`$ |
| 270 | $`2A_1+A_2+A_7+E_7`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`24`$ |
| 271 | $`A_1+A_3+A_7+E_7`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`8`$ |
| 272 | $`A_4+A_7+E_7`$ | $`(0)`$ | $`6`$ | $`2`$ | $`14`$ |
| 273 | $`A_1+A_2+A_8+E_7`$ | $`(0)`$ | $`6`$ | $`0`$ | $`18`$ |
| 274 | $`A_3+A_8+E_7`$ | $`(0)`$ | $`4`$ | $`0`$ | $`18`$ |
| 275 | $`2A_1+A_9+E_7`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`10`$ |
| 276 | $`A_2+A_9+E_7`$ | $`(0)`$ | $`6`$ | $`0`$ | $`10`$ |
| | | $`/(2)`$ | $`4`$ | $`1`$ | $`4`$ |
| 277 | $`A_1+A_{10}+E_7`$ | $`(0)`$ | $`2`$ | $`0`$ | $`22`$ |
| | | | $`6`$ | $`2`$ | $`8`$ |
| 278 | $`A_{11}+E_7`$ | $`(0)`$ | $`4`$ | $`0`$ | $`6`$ |
| 279 | $`D_4+2E_7`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`2`$ |
| 280 | $`A_2+A_4+D_5+E_7`$ | $`(0)`$ | $`6`$ | $`0`$ | $`20`$ |
| 281 | $`A_1+A_5+D_5+E_7`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`12`$ |
| 282 | $`A_6+D_5+E_7`$ | $`(0)`$ | $`6`$ | $`2`$ | $`10`$ |
| 283 | $`A_2+A_3+D_6+E_7`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`6`$ |
| 284 | $`A_5+D_6+E_7`$ | $`/(2)`$ | $`4`$ | $`2`$ | $`4`$ |
| 285 | $`D_5+D_6+E_7`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`4`$ |
| 286 | $`A_1+A_3+D_7+E_7`$ | $`/(2)`$ | $`4`$ | $`0`$ | $`4`$ |
| 287 | $`A_4+D_7+E_7`$ | $`(0)`$ | $`2`$ | $`0`$ | $`20`$ |
| 288 | $`A_1+A_2+D_8+E_7`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`6`$ |
| 289 | $`A_2+D_9+E_7`$ | $`(0)`$ | $`4`$ | $`0`$ | $`6`$ |
| 290 | $`A_1+D_{10}+E_7`$ | $`/(2)`$ | $`2`$ | $`0`$ | $`2`$ |
Table 2. List of extremal elliptic $`K3`$ surfaces
| No | $`\mathrm{\Sigma }`$ | $`MW`$ | $`a`$ | $`b`$ | $`c`$ |
| --- | --- | --- | --- | --- | --- |
| 291 | $`D_{11}+E_7`$ | $`(0)`$ | $`2`$ | $`0`$ | $`4`$ |
| 292 | $`A_2+A_3+E_6+E_7`$ | $`(0)`$ | $`6`$ | $`0`$ | $`12`$ |
| 293 | $`A_1+A_4+E_6+E_7`$ | $`(0)`$ | $`2`$ | $`0`$ | $`30`$ |
| 294 | $`A_5+E_6+E_7`$ | $`(0)`$ | $`6`$ | $`0`$ | $`6`$ |
| 295 | $`D_5+E_6+E_7`$ | $`(0)`$ | $`2`$ | $`0`$ | $`12`$ |
| 296 | $`2A_1+2E_8`$ | $`(0)`$ | $`2`$ | $`0`$ | $`2`$ |
| 297 | $`A_2+2E_8`$ | $`(0)`$ | $`2`$ | $`1`$ | $`2`$ |
| 298 | $`2A_2+2A_3+E_8`$ | $`(0)`$ | $`12`$ | $`0`$ | $`12`$ |
| 299 | $`2A_1+2A_4+E_8`$ | $`(0)`$ | $`10`$ | $`0`$ | $`10`$ |
| 300 | $`A_1+A_2+A_3+A_4+E_8`$ | $`(0)`$ | $`6`$ | $`0`$ | $`20`$ |
| 301 | $`2A_5+E_8`$ | $`(0)`$ | $`6`$ | $`0`$ | $`6`$ |
| 302 | $`A_2+A_3+A_5+E_8`$ | $`(0)`$ | $`6`$ | $`0`$ | $`12`$ |
| 303 | $`A_1+A_4+A_5+E_8`$ | $`(0)`$ | $`2`$ | $`0`$ | $`30`$ |
| 304 | $`2A_2+A_6+E_8`$ | $`(0)`$ | $`6`$ | $`3`$ | $`12`$ |
| 305 | $`2A_1+A_2+A_6+E_8`$ | $`(0)`$ | $`2`$ | $`0`$ | $`42`$ |
| 306 | $`A_1+A_3+A_6+E_8`$ | $`(0)`$ | $`6`$ | $`2`$ | $`10`$ |
| 307 | $`A_4+A_6+E_8`$ | $`(0)`$ | $`2`$ | $`1`$ | $`18`$ |
| 308 | $`A_1+A_2+A_7+E_8`$ | $`(0)`$ | $`2`$ | $`0`$ | $`24`$ |
| 309 | $`2A_1+A_8+E_8`$ | $`(0)`$ | $`2`$ | $`0`$ | $`18`$ |
| 310 | $`A_2+A_8+E_8`$ | $`(0)`$ | $`6`$ | $`3`$ | $`6`$ |
| 311 | $`A_1+A_9+E_8`$ | $`(0)`$ | $`2`$ | $`0`$ | $`10`$ |
| 312 | $`A_{10}+E_8`$ | $`(0)`$ | $`2`$ | $`1`$ | $`6`$ |
| 313 | $`2D_5+E_8`$ | $`(0)`$ | $`4`$ | $`0`$ | $`4`$ |
| 314 | $`A_1+A_4+D_5+E_8`$ | $`(0)`$ | $`2`$ | $`0`$ | $`20`$ |
| 315 | $`A_5+D_5+E_8`$ | $`(0)`$ | $`2`$ | $`0`$ | $`12`$ |
| 316 | $`2A_2+D_6+E_8`$ | $`(0)`$ | $`6`$ | $`0`$ | $`6`$ |
| 317 | $`A_4+D_6+E_8`$ | $`(0)`$ | $`4`$ | $`2`$ | $`6`$ |
| 318 | $`A_1+A_2+D_7+E_8`$ | $`(0)`$ | $`4`$ | $`0`$ | $`6`$ |
| 319 | $`A_1+D_9+E_8`$ | $`(0)`$ | $`2`$ | $`0`$ | $`4`$ |
| 320 | $`D_{10}+E_8`$ | $`(0)`$ | $`2`$ | $`0`$ | $`2`$ |
| 321 | $`A_1+A_3+E_6+E_8`$ | $`(0)`$ | $`2`$ | $`0`$ | $`12`$ |
| 322 | $`A_4+E_6+E_8`$ | $`(0)`$ | $`2`$ | $`1`$ | $`8`$ |
| 323 | $`D_4+E_6+E_8`$ | $`(0)`$ | $`4`$ | $`2`$ | $`4`$ |
| 324 | $`A_1+A_2+E_7+E_8`$ | $`(0)`$ | $`2`$ | $`0`$ | $`6`$ |
| 325 | $`A_3+E_7+E_8`$ | $`(0)`$ | $`2`$ | $`0`$ | $`4`$ |
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# UFR-HEP/ 2000-07 On Hyperkahler Singularities
## 1 Introduction
Over the few past years, there has been an increasing interest in studying the moduli space of vacua of the Coulomb and Higgs branches of supersymmeric gauge theories with eight supercharges in various dimensions. This interest is mostly due to the fact that supersymmetry severely restricts the quantum corrections to the moduli space metric and allows to make many exact computations . A large class of these gauge theories can be realized by D-brane configurations in type II strings on Calabi Yau manifolds .
Recently, a special interest has been given to the analysis of the hypermultiplet gauge invariant moduli space near the Higgs branch singularity. This analysis has been shown to be relevant for the study of many aspects in supersymmetric gauge theories with eight supercharges amongst which we quote: (1) the understanding of the assymptotic regions of the infrared IR low energy limits of the $`N=(4,4)`$ supersymmetric gauge theories in two dimensions especially within the so called throats of the Coulomb and Higgs branches where the theories are typically described by $`2d`$ $`N=4`$ conformal Liouville theories . (2) the study of vector and hypermultiplet moduli spaces in the context of strings compactifications to four dimensions where the low energy supergravity has scalar fields in both vector and matter multiplets . (3) the description of stringy instanton moduli space as hyperkahler deformations of the classical instanton moduli space with non zero NS-NS B field . (4) the rederivation of $`2d`$ CFT’s from compactifications of superstrings on Calabi Yau fourfolds with a $`sp(1)`$ hyperKahler singularity .
The aim of this paper is to study new aspects of hyperkahler singularities of the Higgs branch of supersymmeric gauge theories with eight supercharges by using the linear sigma model method. Actually this study extends the analysis of the ADE singularities of Kahler manifolds to hyperKahler ones. It also generalizes known results, on the Coulomb branch of supersymetric gauge theories with four supercharges, to the Higgs branch of supersymetric gauge theories with eight supercharges.
The presentation of this paper is as follows: In section 2, we review hyperKahler singularities and describe brefly the standard way used in handling the D-flateness eq of supersymmetric theories with eight supercharges. In section 3, we use the harmonic superspace method of supersymmeric gauge theories with eight supercharges to study the problem of hyperKahler singularities. This way of doing has many advantagesas first it preserves manifestly the eight supercharges secon it keeps all the three FI parameters non zero without need to make any appropriate choice. Our method , which is based on realizing the $`SU(2)_R`$ invariance a by helf harmonic analysis on $`S^2`$, permits to exhibit explicitly the role of the three Kahler structures of the gauge invariant hyperkahler moduli space. In Section 4, we study the solving of the D-flateness eqs by introducing a method based on the factorization of the charges of the gauge and the $`SU(2)_R`$ symmetries. We give two classes of solutions of these eqs preserving manifestly the eight supercharges and recverring known resultts on ADE singularities by particular breaking half of the supersymetry. In section 5, we give our conclusion.
## 2 Generalities on Hyperkahler singularities
We begin by noting that Calabi Yau fourfolds can develop singularities of many types; this includes the $`\frac{C^4}{Z_4}`$ orbifold, the ADE hypersurfaces considered recently in and the so called hyperkahler singularity we are intersted in here. We review hereafter the example of the hyperkahler cotangent bundle $`T^{}(CP^2)`$ of complex projective space $`(CP^2)`$ and its generalisation $`T^{}(CP^n)`$ for $`n>2`$. To that purpose, consider $`2d`$ $`N=4`$ supersymmetric $`U(1)`$ gauge theory with one FI isotriplet coupling $`\stackrel{}{\xi }=(\xi ^1,\xi ^2,\xi ^3)`$ and three hypermultiplets of charges $`q_a^i=q^i=1;i=1,2,3`$. The moduli space of this gauge model are obtained by solving
$$\underset{i=1}{\overset{3}{}}[\overline{\phi }_{i\alpha }\phi _i^\beta +\phi _{i\alpha }\overline{\phi }_i^\beta ]=\stackrel{}{\xi }\stackrel{}{\sigma }_\alpha ^\beta ,$$
(2.1)
which by the way is just a special situation of the following eqs:
$$\underset{i=1}{\overset{n}{}}q_a^i[\phi _i^\alpha \overline{\phi }_{i\beta }+\phi _{i\beta }\overline{\phi }_i^\alpha ]=\stackrel{}{\xi _a}\stackrel{}{\sigma }_\beta ^\alpha ;a=1,\mathrm{},r.$$
(2.2)
For later use it is interesting to note that eqs(2.2) have a formal analogy with the sigma model vaccum eqs of $`2d`$ $`N=2`$ supersymmetric $`U(1)^r`$ gauge theory involved in the analysis of the Coulomb branch of IIA superstrings on Calabi Yau threefolds with ADE singularities.
$$\underset{i}{}q_a^i|X_i|^2=R_a;a=1,\mathrm{},r.$$
(2.3)
In eqs (2.3), the $`X_i`$’s are complex scalar fields, the $`R_a`$’s are FI couplings and the $`q_a^i`$’s are the $`U(1)^r`$ gauge charges of the $`X_i`$’s which, for reference, read in the case of a $`SU(n)`$ singularity as :
$$q_a^i=2\delta _a^i+\delta _a^{i1}+\delta _a^{i+1},$$
(2.4)
with the remarkable equality
$$\underset{i}{}q_a^i=0.$$
(2.5)
Eqs (2.1) is a system of three eqs which, up to replacing the Pauli matrices by their expressions and using the $`SU(2)_R`$ transformations $`\phi ^\alpha =\epsilon ^{\alpha \beta }\phi _\beta `$ with $`\epsilon _{12}=\epsilon ^{21}=1`$ and $`\overline{(\phi ^\alpha )}=\overline{\phi }_\alpha `$, split as follows:
$$\begin{array}{ccc}\underset{i=1}{\overset{3}{}}(|\phi _i^1|^2|\phi _i^2|^2)\hfill & =\xi ^3& \hfill (a)\\ \underset{i=1}{\overset{3}{}}\phi _i^1\overline{\phi }_{i2}\hfill & =\xi ^1+i\xi ^2& \hfill (b)\\ \underset{i=1}{\overset{3}{}}\phi _i^2\overline{\phi }_{i1}\hfill & =\xi ^1i\xi ^2& \hfill (c).\end{array}$$
(2.6)
The moduli space of zero energy states of the classical gauge theory is the space of the solutions of eqs (2.1) and (2.6) divided by the action of the $`U(1)`$ gauge group. The solutions of eqs (2.1) depend on the values of the FI couplings. For the case where $`\xi ^1=\xi ^2=\xi ^3=0`$, the moduli space has an $`SU(3)\times SU(2)_R`$ symmetry; it is a cone over a seven manifold described by the eqs:
$$\underset{i=1}{\overset{3}{}}(\phi _{\alpha i}\overline{\phi }_i^\beta \phi _i^\beta \overline{\phi }_{\alpha i})=\delta {}_{\alpha }{}^{\beta }.$$
(2.7)
For the case $`\stackrel{}{\xi }\stackrel{}{0}`$, the abovementioned $`SU(3)\times SU(2)_R`$ symmetry is explicitly broken down to $`SU(3)\times U(1)_R`$. In the remarkable case where $`\xi ^1=\xi ^2=0`$ and $`\xi ^3>0`$, it is not difficult to see that eqs (2.1) describe the cotangent bundle of $`CP^2`$. Indeed making the change
$$\psi _i=\frac{\phi _i^1}{[\underset{i=1}{\overset{3}{}}|\overline{\phi }_{i2}|^2+\xi ^3]^{\frac{1}{2}}},$$
(2.8)
and putting back into eq (2.6.a), one discovers that $`\psi _i`$’s satisfy $`\underset{i}{}|\psi _i|^2=1`$. The $`\psi _i`$’s parametrize the $`CP^2`$ space. On the other hand with $`\xi ^1=\xi ^2=0`$ conditions, eqs (2.6-b,c) may be interpreted to mean that $`\overline{\phi }_{2i}`$ lies in the cotangent space to $`CP^2`$ at the point determined by $`\psi _i`$. Note that though we are usually allowed to make the choice $`\xi ^1=\xi ^2=0`$ by using an appropriate $`SU(2)_R`$ transformation, we shall consider in section 4, the generic cases where $`\xi ^1,\xi ^2`$ and $`\xi ^3`$ are all of them non zero. For the time being let us note that the previous analysis may be extended to the cases of $`2d`$ $`N=4`$ supersymmetric $`U(1)`$ gauge linear sigma model involving $`n+1`$ hypermutiplets with charges $`q^i=1;i=1,\mathrm{},n+1`$ and transforming in the fundamental representation of $`SU(n+1)`$. The vaccum energy equations of this $`U(1)`$ gauge model read as :
$$\begin{array}{ccc}\underset{j=1}{\overset{n+1}{}}(|\phi _j^1|^2|\overline{\phi }_{2j}|^2)\hfill & =\xi ^3& \\ \underset{j=1}{\overset{n+1}{}}(\phi _j^1\overline{\phi }_{2j})\hfill & =\xi ^1+i\xi ^2.& \end{array}$$
(2.9)
For $`\xi ^1=\xi ^2=0`$ and $`\xi ^3>0`$, the classical moduli space of the gauge theory is given by the cotangent bundle of complex $`n`$ projective space: $`T^{}CP^n`$ . For $`\xi ^1=\xi ^2=\xi ^3=0`$, one has just the conifold singularity of $`n`$ dimensional complex manifolds.
## 3 Hyperkahler singularities in Harmonic superspace
Eqs(2.2) is a system of $`3r`$ equations which shares some general features with the usual $`2d`$ $`N=2`$ supersymmetic D-flatness conditions eqs(2.3). To solve them, we shall use a different method inspired from harmonic superspace. The latter was shown to be the appropriate space to deal with supersymmetric quantum field theories with eight supercharges. We shall not develop here the basis of this superspace; details may be found in . Our approach will be done in two steps: First we use $`SU(2)_R`$ harmonic analysis allowing us to realise $`SU(2)_R`$ representations as special functions on $`S_R^2`$. The index R carried by $`S_R^2`$ and $`U(1)_R`$ refers to the $`SU(2)_R`$. Second, we introduce a convenient change of variables based on factorizing the $`U(1)_G^r`$ gauge charges and the $`SU(2)_R`$ ones. We begin by describing the first step of this program. Introducing the following $`2\times 2`$ matrix
$$U=\left(\begin{array}{cc}u_1^+& u_2^+\\ u_1^{}& u_2^{}\end{array}\right)$$
(3.1)
and solving the isospin $`\frac{1}{2}`$ $`SU(2)_R`$ representation constraints namely the unimodularity $`detU=1`$ and the unitarity $`U^+U=U^+U=I`$ conditions, one discovers the defining eq of the unit $`S_R^3`$ sphere:
$$\begin{array}{ccc}u^{\pm \alpha }=ϵ^{\alpha \beta }u_\beta ^\pm ;\overline{u}^{+\alpha }=u_\alpha ^{};ϵ_{\alpha \beta }=ϵ_{\beta \alpha }\hfill & & \\ u^{+\alpha }u_\alpha ^{}=1,u^{+\alpha }u_\alpha ^+=u^\alpha u_\alpha ^{}=0.\hfill & & \end{array}$$
(3.2)
Moreover, using the $`u_\alpha ^\pm `$ variables, the $`SU(2)_R`$ algebra is realized as differential operators on the space of harmonic functions on $`S_R^3`$:
$$\begin{array}{ccc}D^{++}=u^{+\alpha }\frac{}{u^\alpha };D^{}=u^\alpha \frac{}{u^{+\alpha }}\hfill & & \\ 2D^{++}=[D^0,D^{++}];2D^{}=[D^0,D^{}]\hfill & & \\ D^0=[D^{++},D^{}]=u^{+\alpha }\frac{}{u^{+\alpha }}u^\alpha \frac{}{u^\alpha }\hfill & & \end{array}$$
(3.3)
To study the $`SU(2)_R`$ representations by using the harmonic variables, it is more convenient to consider harmonic functions $`F^q(u_\alpha ^\pm )`$ with definte $`U(1)_R`$ charge q ; that is functions $`F^q(u_\alpha ^\pm )`$ satisfying the eigenfunction eq
$$[D^0,F^q]=qF^q.$$
(3.4)
These functions $`F^q`$ have a global harmonic expansion of total charge $`q`$ and carry $`SU(2)_R`$ representations. For example, taking $`q=2`$ and choosing $`F^{++}`$ as:
$$F^{++}(u_\alpha ^\pm )=u_{(\alpha }^+u_{\beta )}^+F^{(\alpha \beta )},$$
(3.5)
one sees that $`F^{++}`$ is the highest state of the isovector representation of $`SU(2)_R`$;i.e:
$$[D^0,F^q]=qF^q,[D^{++},F^q]=0.$$
(3.6)
Note that one can usually realize $`F^{++}`$ as bilinears of isospinors $`f^+`$ and $`\overline{f}^+`$ as follows:
$$\begin{array}{ccc}F^{++}=if^+\overline{f}^+=iu_{(\alpha }^+u_{\beta )}^+f^{(\alpha }\overline{f}^{\beta )}.\hfill & & \end{array}$$
(3.7)
The complex number $`i`$ in front of the the factor of the right hand side of the above eqs ensures the reality condition of the isotriplet representation. After this digression on the $`SU(2)_R`$ harmonic analysis, we turn now to eqs (2.2) which up on multiplying its both sides by $`u_{(\alpha }^+u_{\beta )}^+`$, we get:
$$\underset{j}{}q_a^j\phi _j^+\overline{\phi }_j^+=i\xi _a^{++}.$$
(3.8)
Eqs (3.8) are also the D-flatness eqs one gets if one is using the $`2d`$ $`N=(4,4)`$ harmonic superspace formulation of $`2d`$ $`N=4`$ gauge theories .
In the end of this section, we would like to make two comments: First one can use the isospinor bilinear realization of isotriplets eqs (3.7) to represent the Kahler parameters $`\xi _a^{++}`$ as follows:
$$\begin{array}{ccc}\xi _a^{++}=i\zeta _a^+\overline{\zeta }_a^+=iu_{(\alpha }^+u_{\beta )}^+\zeta _a^{(\alpha }\overline{\zeta }_a^{\beta )}\hfill & & \\ \zeta _a^\pm =u_\alpha ^\pm \zeta _a^\alpha ;\overline{\zeta }_a^\pm =u_\alpha ^\pm \overline{\zeta }_a^\alpha ,\hfill & & \end{array}$$
(3.9)
where $`\zeta _a^\alpha `$ and $`\overline{\zeta }_a^\alpha `$ may, roughly speaking, be viewed as the square roots of the FI couplings $`\xi _a^{(\alpha \beta )}`$. Similar relations involving $`\zeta _a^\pm `$ and $`\overline{\zeta }_a^\pm `$ for $`\xi _a^0`$ and $`\xi _a^{}`$ may be also written down. Putting back these relations in eqs (3.8), one gets
$$\underset{j}{}q_a^j\phi _j^+\overline{\phi }_j^+=\zeta _a^+\overline{\zeta }_a^+=u_{(\alpha }^+u_{\beta )}^+\zeta _a^{(\alpha }\overline{\zeta }_a^{\beta )}.$$
(3.10)
Second, simplify further eqs (3.10) by making an extra change of variables which turns out to convenient when discussing the moduli space of gauge invariant vacua of eqs (2.2). This extra change consists to use the mapping $`R^3=R^+\times S^2`$ to write the FI isovectors $`\xi _a^{++}`$ as
$$\xi _a^{++}=R_a\eta _a^+\overline{\eta }_a^+=r_a^2\eta _a^+\overline{\eta }_a^+,$$
(3.11)
or equivalently by using the isospinors $`\zeta _a^\alpha `$ and $`\overline{\zeta }_a^\alpha `$ introduced previously:
$$\begin{array}{ccc}\zeta _a^\pm \hfill & =u_\alpha ^\pm \zeta _a^\alpha ;\overline{\zeta }_a^\pm =u_\alpha ^\pm \overline{\zeta }_a^\alpha & \hfill (a)\\ \zeta _a^\alpha \hfill & =r_a\eta _a^\alpha ;\overline{\zeta }_a^\alpha =r_a\overline{\eta }_a^\alpha & \hfill (b)\\ \zeta _a^\alpha \overline{\zeta }_{a\alpha }\hfill & =r_a^20& \hfill (c).\end{array}$$
(3.12)
Eqs (3.11) and (3.12) tell us that the $`R_a`$’s $`(R_a=r_a^20)`$ are the radial variables and the $`\eta _a^\alpha `$’s and $`\overline{\eta }_{a\alpha }`$’s, which satisfy
$$\eta _a^\alpha \overline{\eta }_{a\alpha }=1;\eta _a^\alpha \eta _{a\alpha }=\overline{\eta }_a^\alpha \overline{\eta }_{a\alpha }=0,$$
(3.13)
parametrize the two spheres $`S_a^2`$. The $`r_a^2`$ and $`\eta _a^\alpha `$ and $`\overline{\eta }_{a\alpha }`$ are in one to one correspondance with the $`r`$ FI isovectors. In other words eqs (3.13) describe a collection of $`r`$ unit two spheres which altogether with the $`r_a`$ conical variables of $`(R^3)^r`$ give the $`3r`$ parameters of the $`r`$ FI isovector couplings.
## 4 Solving the D-flateness eqs
In this section, we solve eqs(2.2) by developing a factorization method of the gauge and $`SU(2)_R`$ charges carried by the hyperKahler moduli. This factorization should satisfy the following natural constraints:
(a) The splitting of the gauge and $`SU(2)_R`$ charges preserves the eight supercharges of the gauge theory.
(b)It recovers the results of , extends the ADE models of Kahler backgrounds and, up on breaking half of the eight supercharges, gives the standard ADE results.
(c) It has a geometrical interpretation.
The factorization of hyperKahler moduli solving the above constraints is given by:
$$\begin{array}{ccc}\phi _j^+=X_j\eta _j^++\gamma Y_j\overline{\eta }_j^+\hfill & & \\ \overline{\phi }_j^+=\overline{X}_j\overline{\eta }_j^+\gamma \overline{Y}_j\overline{\eta }_j^+,\hfill & & \end{array}$$
(4.1)
where $`\eta _j^+`$ and $`\overline{\eta }_j^+`$ are as in eqs (3.12); that is:
$$\begin{array}{ccc}\eta _j^+=u_\alpha ^+\eta _j^\alpha ;\overline{\eta }_j^+=u_\alpha ^+\overline{\eta }_j^\alpha \hfill & & \\ \eta _j^\alpha \overline{\eta }_{\alpha j}=1;\eta _j^\alpha \eta _{\alpha j}=\overline{\eta }_j^\alpha \overline{\eta }_{\alpha j}=0.\hfill & & \end{array}$$
(4.2)
In eqs(4.1), $`X_j`$ and $`Y_j`$, $`j=1,\mathrm{},n`$ are complex fields carrying no $`SU(2)_R`$ charge. The parameter $`\gamma `$ takes the values $`\gamma =0`$ or $`\gamma =1`$ and allow to distinguish two classes of solutions we will give hereafter. Moreover the quantities $`X_j`$ , $`Y_j`$ and $`\eta _j^+`$ and $`\overline{\eta }_j^+`$ of eqs(4.1) behave under $`U(1)^r`$ gauge and $`SU(2)_R`$ transformations as follows:
$$\begin{array}{ccc}U(1)^r:X_jX_j^{}=\lambda ^{q_a^j}X_j\hfill & & \\ Y_jY_j^{}=\lambda ^{q_a^j}y_j\hfill & & \\ \eta _j^+\eta _j^+=\eta _j^+\hfill & & \\ \overline{\eta }_j^+\overline{\eta }_j^+=\overline{\eta }_j^+;\hfill & & \end{array}$$
(4.3)
and
$$\begin{array}{ccc}U(1)_R:X_jX_j^{}=X_j\hfill & & \\ Y_jY_j^{}=Y_j\hfill & & \\ \eta _j^+\eta _j^+=e^{i\theta }\eta _j^+\hfill & & \\ \overline{\eta }_j^+\overline{\eta }_j^+=e^{i\theta }\overline{\eta }_j^+.\hfill & & \end{array}$$
(4.4)
Actually eqs (4.3-4) define the factorization of the gauge charges and $`SU(2)_R`$ ones of the vacum moduli. In what follows we shall use the splitting eqs(4.1) to solve eqs (2.2) which, by help of the analysis of section 3, may be put in the form:
$$\underset{j}{}q_a^j\phi _j^+\overline{\phi }_j^+=r_a^2\eta _a^+\overline{\eta }_a^+.$$
(4.5)
We give herebelow two classes of solutions of these eqs. The first class describes a generalisation of the usual ALE surfaces with ADE singularities. The second class leads to new models which flow in the infrared to $`2d`$ $`N=(4,4)`$ conformal field theories.
### 4.1 Generalized ADE hypersurfaces
Eqs(2.2) looks formally like eqs (2.3); their solutions are then expected to describe generalisations of the standard eqs of ADE singularities associated with $`2d`$ $`N=2`$ linear sigma models. Imposing the ADE Calabi Yau condition
$$\begin{array}{ccc}\underset{j}{}q_a^j=0,a=1,\mathrm{},n1,\hfill & & \end{array}$$
(4.6)
one can imitate the analysis of $`N=2`$ linear $`\sigma `$ models and build the gauge invariant moduli in terms of the $`\phi _j^+`$ fields. In the $`SU(n)`$ case for instance where $`q_a^j`$ is given by eq(2.4), there are three gauge invariant moduli; $`U^{+\frac{n(n+1)}{2}}`$, $`V^{+\frac{n(n+1)}{2}}`$ and $`Z^{+(n+1)}`$ carrying $`\frac{n(n+1)}{2}`$ , $`\frac{n(n+1)}{2}`$ and $`(n+1)`$ $`U(1)_R`$ Cartan charges respectively. These invariant read as:
$$\begin{array}{ccc}U^{+\frac{n(n+1)}{2}}=\underset{j=0}{\overset{n}{}}(\phi _j^+)^{nj}\hfill & & \\ V^{+\frac{n(n+1)}{2}}=\underset{j=0}{\overset{n}{}}(\phi _j^+)^j\hfill & & \\ Z^{+(n+1)}=\underset{j=0}{\overset{n}{}}(\phi _j^+).\hfill & & \end{array}$$
(4.7)
They satisfy the following remarkable equation,
$$U^{+\frac{n(n+1)}{2}}V^{+\frac{n(n+1)}{2}}=[Z^{+(n+1)}]^n.$$
(4.8)
Eq (4.8) generalizes the usual equation of the ALE surface with $`SU(n)`$ singularity:
$$uv=z^n.$$
(4.9)
To better see the structure of eq (4.8), we use the $`\phi _j^+`$’s moduli factorization described earlier. Taking $`\gamma =0`$, the general splitting eqs (4.1) reduces to:
$$\phi _j^+=x_j\eta _j^+;\overline{\phi }_j^+=\overline{x}_j\overline{\eta }_j^+,$$
(4.10)
where $`X_j`$ and $`\eta _j^+`$ behave under gauge and $`SU(2)_R`$ transformations as in eqs (4.3-4). Note that like $`\phi _j^\pm `$, the realization $`X_j\eta _j^\pm `$ carries for each value of j, four real degrees of freedom; two coming from $`X_j`$ and the two others from the parameters the sphere $`S^2`$ described by $`\eta _j^\pm `$. Under $`2d`$ $`N=4`$ supersymmetric transformations which may be conveniently expressed as $`4d`$ $`N=2`$ supersymmetric transformations of fermionic parameters $`ϵ^\pm `$ and $`\overline{ϵ}^\pm `$, we have:
$$\delta \phi _j^+=ϵ^+\psi _j+\overline{ϵ}^+\overline{\chi }_j;$$
(4.11)
where $`\psi _j`$ and $`\overline{\chi }_j`$ are the Fermi partners of the $`\phi _j^\pm `$ scalars. $`(\phi _j^\pm ,\psi _j,\overline{\chi }_j)`$ constitute altogether the $`4d`$ $`N=2`$ free hypermultiplets. Using the splitting principle by factorizing $`ϵ^+`$ as $`ϵ\eta ^+`$ and $`\overline{ϵ}^+=\overline{ϵ}\overline{\eta }^+`$, and using eqs (4.10-11), we get
$$\eta _j^+\delta X_j+X_j\delta \eta _j^+=\overline{\eta }^+ϵ\psi _j+\overline{\eta }^+\overline{ϵ}\overline{\chi }_j,$$
(4.12)
or equivalenty
$$\eta _j^\alpha \delta X_j+X_j\delta \eta _j^\alpha =\overline{\eta }^\alpha ϵ\psi _j+\overline{\eta }^\alpha \overline{ϵ}\overline{\chi }_j.$$
(4.13)
Putting eqs (4.10) back into eqs (4.5), we obtain
$$\underset{j}{}q_a^j|X_j|^2\eta _j^+\overline{\eta }_j^+=r_a^2\eta _a^+\overline{\eta }_a^+$$
(4.14)
and
$$\begin{array}{ccc}U^{+\frac{n(n+1)}{2}}=uM^{+\frac{n(n+1)}{2}}\hfill & & \\ V^{+\frac{n(n+1)}{2}}=vN^{+\frac{n(n+1)}{2}}\hfill & & \\ Z^{+(n+1)}=zS^{+(n+1)};\hfill & & \end{array}$$
(4.15)
where $`u,v,z`$ and $`M^{+\frac{n(n+1)}{2}}`$, $`N^{+\frac{n(n+1)}{2}}`$ and $`S^{+(n+1)}`$ are gauge invariants given by:
$$\begin{array}{ccc}u=\underset{j=0}{\overset{n}{}}X_j^j\hfill & ;N^{+\frac{n(n+1)}{2}}& \hfill =\underset{j=0}{\overset{n}{}}(\eta _j^+)^{nj}\\ v=\underset{j=0}{\overset{n}{}}X_j^{nj}\hfill & ;M^{+\frac{n(n+1)}{2}}& \hfill =\underset{j=0}{\overset{n}{}}(\eta _j^+)^j\\ z=\underset{j=0}{\overset{n}{}}X_j\hfill & ;S^{+(n+1)}& \hfill =\underset{j=0}{\overset{n}{}}\eta _j^+.\end{array}$$
(4.16)
Note that $`u,v`$ and $`z`$ verify the relation (4.9) and $`M^{+\frac{n(n+1)}{2}}`$ , $`N^{+\frac{n(n+1)}{2}}`$ and $`S^{+(n+1)}`$ satisfy eq (4.8). Note moreover that eqs (4.14-15) may be brought to more familiar forms if one requires that all 2-spheres are identified; i.e:
(i) The $`(n+1)`$ 2- spheres parametrized by the $`\eta _j`$’s .
(ii) The $`(n1)`$ $`\eta _a`$ 2-spheres used in the parametrization of the FI couplings eqs (3.13).
(iii) The $`\eta ^+`$ 2-sphere involved in the factorization of the supersymmetric parameter $`ϵ^+`$ eq(4.12).
Put differently, we require the following constraint eq to hold:
$$\eta _j^+=\eta _a^+=\eta ^+.$$
(4.17)
With this identification, eqs (4.14) reduce to the well known D- flatness conditions of the $`U(1)^r`$ supersymmetric gauge theory with four supercharges; namely:
$$\underset{j}{}q_a^j|X_j|^2=r_a^2.$$
(4.18)
Moreover eq (4.8) reduces to the usual ALE surface with $`SU(n)`$ singularity (4.9) since the $`M^{+\frac{n(n+1)}{2}}`$, $`N^{+\frac{n(n+1)}{2}}`$ and $`S^{+(n+1)}`$ gauge invariant become trivial; they are given by powers of $`\eta ^+`$ as shown here below.
$$\begin{array}{ccc}M^{+\frac{n(n+1)}{2}}=(\eta ^+)^{\frac{n(n+1)}{2}}=N^{+\frac{n(n+1)}{2}}\hfill & & \\ S^{+(n+1)}=(\eta ^+)^{n+1}.\hfill & & \end{array}$$
(4.19)
In the general case where the gauge charges $`q_a^j`$ of the $`X_j`$’s satisfy the constraints(2.5), eqs(4.18) is the vaccum energy of $`2d`$ $`N=2`$ supersymmetric linear $`\sigma `$ models. Thus, the classical moduli space $``$ of the gauge invariant vacua of eqs(4.17-18) is given by the product of the 2-sphere parametrized by $`\eta ^+`$; eq(4.17)and the moduli space of the gauge invariant solutions of $`2d`$ $`N=2`$ vaccum energy states. In other words:
$$=\frac{C^{n+1}}{C_{}^{}{}_{}{}^{n1}}\times S^2.$$
Note that the identification constraint eq(4.17) has a nice interpretation; it breaks explicitly half of the eight supersymmetries leaving then four supercharges manifest. These four supercharges are behind the reduction of eqs(4.14) down to eqs(4.18)leading to the standard ADE models. This feature is immediately derived by combining eqs(4.13) and (4.17) as follows :
$$\eta ^\alpha \delta X_j+X_j\delta \eta ^\alpha =\overline{\eta }^\alpha ϵ\psi _j+\overline{\eta }^\alpha \overline{ϵ}\overline{\chi }_j.$$
(4.20)
Then multiplying both sides of this identity by $`\overline{\eta }_\alpha `$; one gets, after using eqs(4.2):
$$\delta X_j=ϵ\psi _j.$$
(4.21)
This eq gives the usual supersymmetric transformations of the complex scalars of the $`2d`$ $`N=2`$ chiral multiplets. This completes our check of consistency of the generalised $`SU(n)`$ hypersurface singularity (4.8). Before going ahead let us summarize in few words what we have been doing. Starting from eqs(2.2), we have shown that it is possible to put them into the equivalent form (4.5). The corresponding moduli space of gauge invaraint vacua is given by eq (4.8) which reduces to the standard ALE space with $`A_{n1}`$ singularity up on imposing the factorization eqs (4.10) and the conditions (4.17). The latters break four of the original eight supercharges. To restore the eight supersymmetries by still using the constraint (4.17), we should take $`\gamma `$ non zero; say $`\gamma =1`$. Non zero $`\gamma `$ restores four extra supercharges which add to the old four existing ones carried by eq (4.10). Note in passing that a naive analysis of eqs(4.5) suggests that the gauge invariant moduli space of vacua $``$ of eq (4.10) is given, for the generalized $`SU(n)`$ singularity eq (4.8), by the usual ALE space with $`SU(n)`$ singularity times the two-sphere power $`2n`$ :
$$=\frac{C^{n+1}}{C_{}^{}{}_{}{}^{n1}}\times (S^2)^{2n},$$
where $`(n+1)`$ two spheres come from the $`\phi _j^+`$’s as shown in eqs (4.10) and $`(n1)`$ two spheres come from the FI couplings.
### 4.2 Second Solution
Putting $`\gamma =1`$ in eqs(4.1) and (4.5), we get the system of three eqs given by:
$$\begin{array}{ccc}\underset{j}{}q_a^j(|X_j|^2|Y_j|^2)\eta _j^+\overline{\eta }_j^+\hfill & =r_a^2\eta _a^+\overline{\eta }_a^+& \hfill (a)\\ \underset{j}{}q_a^j(X_j\overline{Y}_j)\eta _j^+\eta _j^+\hfill & =0& \hfill (b)\\ \underset{j}{}q_a^j(\overline{X}_jY_j)\overline{\eta }_j^+\overline{\eta }_j^+\hfill & =0.& \hfill (c)\end{array}$$
(4.22)
At this level no constraint on the FI couplings has been imposed yet. If moreover we require that all the two sphere $`\eta _j^+`$ and $`\eta _a^+`$ are identified as in eq(4.17); the above system reduces to:
$$\begin{array}{ccc}\underset{j}{}q_a^j(|X_j|^2|Y_j|^2)\hfill & =r_a^2& \hfill (a)\\ \underset{j}{}q_a^j(X_j\overline{Y}_j)\hfill & =0& \hfill (b)\\ \underset{j}{}q_a^j(\overline{X}_jY_j)\hfill & =0.& \hfill (c)\end{array}$$
(4.23)
Eqs(4.22) have remarkable features which have nice interpretations. Though the $`q_a^j`$ gauge charges of the hypermultiplet moduli are not required to add to zero as in eq (2.5), eq (4.22.a) behave exactly as the D-flatness condition of $`2d`$ $`N=4`$ supersymmetric $`U(1)^r`$ gauge theory. The point is that eqs(4.22.a) involve twice the number of fields of eqs(2.3), but with opposite charges $`q_a^j`$. Put differently; eq(4.22) involve two sets of fields $`X_j`$ and $`Y_j`$ of charge $`q_a^j`$ and ($`q_a^j`$) respectively. The sum of gauge charges of the $`X_j`$’s and $`Y_j`$’s add automatically to zero eventhough eq(2.5)is not fulfilled.
$$\underset{j}{}q_a^j+\underset{j}{}(q_a^j)=0.$$
(4.24)
Eqs(4.24) go beyond the constraint eqs(2.5)and so models with $`\gamma =1`$ flow in the IR to $`2d`$ $`N=(4,4)`$ superconformal models extending the usual $`2d`$ $`N=(2,2)`$ ADE ones. Moreover eqs(4.24) may be fulfilled in different ways; either by taking all charges $`q_a^j`$ of the $`U(1)^r`$ gauge theory to be positive; say $`q_a^j=1;a=1,\mathrm{},r;j=1,\mathrm{},n`$, or part of the $`q_a^j`$’s are positive and the remaining ones are negative. In the case of a supersymmetric $`U(1)`$ gauge theory with $`(n+1)`$ hypermutiplets with gauge charges equal to one, eqs(4.23.a) describe a $`CP^n`$ manifold while eqs(4.23-b,c) which read as
$$\underset{j}{}X_j\overline{Y_j}=0,$$
(4.25)
together with their complex conjugate, show that the $`\overline{Y_j},`$’s are in the cotangent space of $`CP^n`$ at the point $`x_j=X_j/[\underset{i}{}|Y_i|^2+r_a^2]^{\frac{1}{2}}`$. Observe that in case where some of the positive charges $`q_a^j`$ of $`U(1)^r`$ gauge theory are not equal to one, the corresponding moduli space is just the cotangent bundle of some weighted complex projective space, $`T^{}(WP^n)`$. Obesrve moreover that in the infrared limit this gauge theory flows to a $`2d`$ $`N=(4,4)`$ conformal field theory with central charge $`C=6n`$.
## 5 Conclusion
In this paper we have studied the solutions of the D-flatness eqs of supersymmetric $`U(1)^r`$ gauge theories with eight supercharges by using the harmonic superspace linear sigma model approach . We have studied the blown up of hyperkahler singularities and shown that they are given by cotangent bunbles of compact weighted projective spaces. The latters depend on the number of hypermultiplets and gauge supermultiplets involved in the gauge model one is considering. This examination extends the standard linear sigma model analysis performed for the Kahler Coulomb branch of supersymmetric gauge theories with four supercharges. Our method, which go beyond standard analysis where only half of the eight supersymmetries are apparent, preserves manifesty all the eight supersymmetries and is done in two steps. First we have used a geometric realization of the $`SU(2)_R`$ symmetry to transform the D-flatness eq in a more convenient form quite easy to handle. Second we have introduced a method of factorization of gauge and $`SU(2)_R`$ charges of the hypermultiplet moduli. This factorisation involves an index $`\gamma `$ taking the values $`\gamma =0`$ or $`\gamma =1`$ which allow to distinguish two classes of solutions of eqs(2.2)preserving the eight supercharges. For $`\gamma =0`$, we have obtained a generalisation of the ADE complex surfaces which read in the non affine case as:
$$\begin{array}{ccc}A_{n1}:U^{+\frac{n(n+1)}{2}}V^{+\frac{n(n+1)}{2}}=[Z^{+(n+1)}]^n\hfill & & \\ D_n:(x^{++})^n+x^{++}(y^{+(n1)})^2+(z^{+n})^2=0\hfill & & \\ E_6:(x^{+6})^2+(y^{+4})^3+(z^{+3})^4=0\hfill & & \\ E_7:(x^{+9})^2+(y^{+6})^3+y^{+6}(z^{+4})^3=0\hfill & & \\ E_8:(x^{+15})^2+(y^{+10})^3+(z^{+6})^5=0.\hfill & & \end{array}$$
(5.1)
The above eqs reproduce the standard ADE singularities by partial breaking of $`2d`$ $`N=4`$ supersymmetry down to $`2d`$ $`N=2`$. For $`\gamma =1`$, we have found new models which flow in the infrared to $`2d`$ $`N=(4,4)`$ scale invariant models with central charge c= 6k.
This research work has been supported by the program PARS number phys 27/ 372/98 CNR.
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# Superconductivity in the SU(N) Anderson Lattice at 𝑈=∞
## I Introduction
The superconducting behavior of heavy-fermion materials has attracted much attention due to its non-conventional properties. Despite the large amount of work trying to understand heavy-fermion superconductivity, the normal state properties, the symmetry of the order parameter, the origin of superconductivity and the interplay between superconductivity and magnetism are still interesting and open questions.
Some of these materials, such as $`UAgCu_4`$, $`UCu_7`$, $`U_2Zn_{17}`$, order antiferromagnetically at low temperatures while others (such as $`UBe_{13}`$, $`CeCu_2Si_2`$, $`UPt_3`$) order in a superconducting state and others show no ordering (such as $`CeAl_3`$, $`UAuPt_4`$, $`CeCu_6`$, $`UAl_2`$). There are materials which order both antiferromagnetically and become superconducting as the temperature drops (e.g. $`URu_2Si_2`$, $`U_{0.97}Th_{0.03}Be_{13}`$) and it has recently been found that $`UPd_2Al_3`$ shows coexistence of superconductivity and local moment antiferromagnetism. All these materials have very large specific heat coefficients $`\gamma `$, indicating very large effective masses, hence the designation heavy fermions.
The superconducting properties of a system depend on the type of ground state that the system exhibits in the normal phase. The large specific heat $`\gamma `$coefficient can have two very different origins: a Kondo-impurity behavior, in which case $`\gamma `$ behaves as the inverse of the Kondo temperature $`T_K`$; or a Kondo-lattice behavior, in which case $`\gamma `$ is controlled by a large density of states at the Fermi energy. The large density of states arises from a hybridization mechanism between the conduction band and localized electronic states ($`f`$states, say).
Even though the large effective masses indicate strong correlations between the electrons, these behave in many cases as essentially “free”, with renormalized parameters, as explained by the Fermi liquid theory. However, there has recently been growing evidence that other materials have properties that do not fit the Fermi liquid picture. The reason could be either disorder, vicinity to a quantum phase transition, or unusual impurity-like behavior such as the one described by generalized models, as the $`n`$-channel Kondo model. The n-channel Kondo lattice shows interesting behavior and it has been shown to be an incoherent metal at low temperatures with a residual entropy that is usually lifted via ordering at very low temperatures.
A consistent description of the overall properties of the heavy-fermion behavior has been achieved assuming that a generalization of the impurity Anderson model to the lattice case is valid. In the Anderson lattice the energy of a single electron in an $`f`$orbital (e. g. $`4f^1`$) is $`ϵ_0`$, and the energy of two electrons in the same $`f`$orbital ($`4f^2`$) is $`2ϵ_0+U`$, where $`U`$ is the on-site Coulomb repulsion. The energy of the $`4f^2`$ state is much larger than the energy of the $`4f^1`$ state. Moreover, these systems are often characterized by large angular momentum, due to the spin-orbit coupling. In general, both the large values of $`U`$ and the large total angular momentum must be included in any model used to describe the properties of heavy-fermion materials.
The SU(N) Anderson lattice Hamiltonian is believed to give a good description of the normal state of Kondo-lattice systems. The limit $`U=\mathrm{}`$ is considered in many calculations since the experimental $`U`$ values are large. The Anderson lattice model predicts Fermi-liquid like properties in the normal non-magnetic state. The theoretical results give a good description of many materials and explain the main features at low temperatures such as universality, large effective masses, the Kondo resonance at the Fermi level. At the single-impurity level the picture is clear. In the Kondo limit the $`f`$-level has an occupation close to one leading to a localized spin that is shielded by a conduction electron spin cloud. This compensation of the spin explains why some of these compounds do not order magnetically. The main point to be explained in the lattice case is the competition between the Kondo compensation of the localized spins and the magnetic interactions between them. In these materials this interaction is mediated by the conduction electrons ($`RKKY`$-type). Actually, since the Kondo temperature is very small it is difficult to explain why the RKKY does not always prevail. Related to this competition is the effectiveness of the compensating cloud around each $`f`$-level. The size of this cloud has been subject of controversy. Arguments show that it should be a large scale of the order of $`v_F/T_K`$ but other arguments claim to be $`a`$ ($`a`$ is the lattice constant). This is a relevant issue in the lattice case related to Nozières exhaustion problem which states that there are not enough conduction electrons to screen the $`f`$-levels.
To increase the complexity the system may also order into a superconducting state. Many questions have been raised starting from the result that the discontinuity of the specific heat at $`T_c`$ is large, of the order of the specific heat itself in the normal phase (which originates in the heavy fermions). This indicates that pairing occurs between the heavy $`f`$-level electrons, which will then form the condensate. Within the Anderson lattice model the strong correlations and the hybridization are responsible for the high effective masses and it has been proposed that the mechanism for superconductivity lies in the strong Coulomb interaction between the $`f`$-electrons, not in a phonon mediated attraction.
Using Coleman’s slave boson formalism together with a large-$`N`$ approach, various attempts have been made to search for the existence of superconducting instabilities in the infinite-$`U`$ Anderson-lattice model. It was proposed that slave bosons fluctuations can provide an effective attraction between the electrons to leading order in $`1/N`$. Later, a calculation of the electron-electron scattering amplitude to order $`1/N^2`$ revealed an effective attractive interaction in the $`p`$ and $`d`$ channels, which was interpreted as a manifestation of the $`RKKY`$ interaction, showing that spin fluctuations are an important mechanism.
Assuming that the normal state is a Fermi liquid, several other studies of superconductivity have been carried out on the Anderson lattice model and generalizations of it. By adding an attractive nearest-neighbor interaction between the $`f`$-electrons, so as to explicitly provide an attractive channel leading to superconductivity, a mean-field study has been carried out as a function of the local repulsion $`U`$. Romano, Noce, and Micnas, have found a superconducting ground state for finite values of $`U`$, but no superconductivity was found for large values of onsite Coulomb repulsion, in the Anderson lattice. This is so because the authors consider the Kondo regime (this is, $`ϵ_0\mu `$ where $`\mu `$ is the chemical potential), where the occupation number of an $`f`$orbital, $`n_f`$, is close to two for small $`U`$. Therefore, upon increasing the interaction $`U`$, this number is reduced to one, blocking charge transport in the $`f`$-band.
In this paper we carry out a mean field study of superconductivity in the $`U=\mathrm{}`$ Anderson lattice where an attractive interaction between neighboring $`f`$-orbitals is explicitly introduced in order to simulate an effective interaction (which might have various causes) leading to superconductivity. Since $`U=\mathrm{}`$, we are restricted to $`f`$-level occupancies in the range $`0<n_f<1`$. In the mixed valent regime, where $`n_f`$ is between zero and one, charge movement is allowed among the $`f`$ orbitals, even when $`U=\mathrm{}`$. We study the dependence of the critical temperature and $`f`$-level ocupancy on the various model parameters for different Cooper pairing symmetries. The paper is organized as follows: in section II we present the model Hamiltonian we use in our study and derive the mean field equations. Particular attention is paid on the form of the superconducting pairing term. In section III we present our calculations of the critical temperature as function of the several parameters of the model and we summarize our findings in section IV.
## II The Model Hamiltonian
We consider an extended version of the Anderson lattice model, which includes a density-density attraction between the electrons occupying neighboring $`f`$-orbitals. This form of interaction enables us to consider three possible symmetries for electron pairing: $`s`$, $`d`$ and $`p`$-wave. The Hamiltonian is given by
$$H=H_c^0+H_f^0+H_{cf}+H_U+H_J,$$
(1)
where
$`H_f^0`$ $`=`$ $`{\displaystyle \underset{i,m}{}}(ϵ_0\mu )f_{i,m}^{}f_{i,m},`$ (2)
$`H_c^0`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{k},m}{}}(ϵ_\stackrel{}{k}\mu )c_{\stackrel{}{k},m}^{}c_{\stackrel{}{k},m},`$ (3)
$`H_{cf}`$ $`=`$ $`V{\displaystyle \underset{i,m}{}}\left(c_{i,m}^{}f_{i,m}+f_{i,m}^{}c_{i,m}\right),`$ (4)
$`H_U`$ $`=`$ $`U{\displaystyle \underset{i,mm^{}}{}}n_{i,m}n_{i,m^{}},`$ (5)
and
$$H_J=\frac{1}{2}J\underset{<i,j>,m,m^{}}{}n_{i,m}n_{j,m^{}},$$
(6)
where $`i`$ and $`j`$ are nearest neighbor sites and $`n_{i,m}=f_{i,m}^{}f_{i,m}`$. The $`c`$ and $`f`$ operators are fermionic and obey the usual anti-commutation relations. The hybridization potential $`V`$ is assumed to be momentum independent. The term $`H_U`$ represents the strong onsite repulsion between the $`f`$-orbitals and in the rest of this work we shall consider $`U=\mathrm{}`$. The term $`H_J`$ explicitly describes an effective attraction between neighboring $`f`$-sites ($`J<0`$) which is responsible for superconductivity. The total angular momentum projection $`m`$ takes on $`N`$ values. We shall assume that the local angular momentum of the $`f`$-sites is half-integer and, therefore, that $`N`$ is even.
The term $`H_J`$ may be re-written in momentum space as
$$H_J=\underset{\stackrel{}{Q},\stackrel{}{k},\stackrel{}{k}^{}}{}\underset{m,m^{}}{}\frac{J_{\stackrel{}{k},\stackrel{}{k}^{}}}{2}f_{}^{}{}_{\frac{\stackrel{}{Q}}{2}+\stackrel{}{k}^{},m}{}^{}f_{}^{}{}_{\frac{\stackrel{}{Q}}{2}\stackrel{}{k}^{},m^{}}{}^{}f_{\frac{\stackrel{}{Q}}{2}\stackrel{}{k},m^{}}f_{\frac{\stackrel{}{Q}}{2}+\stackrel{}{k},m},$$
(7)
where the interaction $`J_{\stackrel{}{k},\stackrel{}{k}^{}}=J_\stackrel{}{\delta }\mathrm{exp}i(\stackrel{}{k}\stackrel{}{k}^{})\stackrel{}{\delta }`$ and the summation over $`\stackrel{}{\delta }`$ runs over the nearest neighbors. Considering the case of a cubic lattice, the interaction $`J_{\stackrel{}{k},\stackrel{}{k}^{}}`$ may be separated into terms with $`s`$, $`p`$ and $`d`$wave symmetries as:
$`J_{\stackrel{}{k},\stackrel{}{k}^{}}`$ $`=`$ $`J\left(\eta _\stackrel{}{k}^{(s)}\eta _\stackrel{}{k}^{}^{(s)}+{\displaystyle \underset{i=x,y,z}{}}\eta _\stackrel{}{k}^{(p,i)}\eta _\stackrel{}{k}^{}^{(p,i)}\right)`$
$`+`$ $`J\left(\eta _\stackrel{}{k}^{(d_{x^2y^2})}\eta _\stackrel{}{k}^{}^{(d_{x^2y^2})}+\eta _\stackrel{}{k}^{(d_{r^23z^2})}\eta _\stackrel{}{k}^{}^{(d_{r^23z^2})}\right),`$
where
$`\eta _\stackrel{}{k}^{(s)}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}\left[\mathrm{cos}(k_x)+\mathrm{cos}(k_y)+\mathrm{cos}(k_z)\right],`$ (8)
$`\eta _\stackrel{}{k}^{(p,i)}`$ $`=`$ $`\sqrt{2}\mathrm{sin}(k_i),`$ (9)
$`\eta _\stackrel{}{k}^{(d_{x^2y^2})}`$ $`=`$ $`\mathrm{cos}(k_x)\mathrm{cos}(k_y),`$ (10)
$`\eta _\stackrel{}{k}^{(d_{r^23z^2})}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\left[\mathrm{cos}(k_x)+\mathrm{cos}(k_y)2\mathrm{cos}(k_z)\right].`$ (11)
Electron pairing in the superconducting phase will occur in the state with total pair momentum $`\stackrel{}{Q}=0`$.
We implement the condition $`U=\mathrm{}`$ within the slave-boson formulation due to Coleman, in which the empty $`f`$-site is represented by a slave boson $`b_i`$ and the physical operator $`f_i`$ in equation (4) is replaced with $`b_{}^{}{}_{i}{}^{}f_i`$. Condensation of the slave-bosons can be described by the replacement $`b_i<b_i>=<b_{}^{}{}_{i}{}^{}>=\sqrt{z}`$. The mean-field treatment of the interaction term $`H_J`$ involves the usual decoupling of destruction and annihilation operators but, in keeping with the spirit of Coleman’s slave boson formalism, we associate a boson operator with every $`f`$ operator in (7) in order to prevent double occupancy at the $`f`$-sites. Taking also into account the boson condensation, we obtain the superconducting part of the mean-field Hamiltonian from the substitution: $`f^{}f^{}ffzf^{}f^{}<zff>+h.c.`$. Following these ideas we write down the effective Hamiltonian as:
$`H_{eff}`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{k},m}{}}\left((ϵ_\stackrel{}{k}\mu )c_{\stackrel{}{k},m}^{}c_{\stackrel{}{k},m}+(ϵ_f\mu )f_{\stackrel{}{k},m}^{}f_{\stackrel{}{k},m}\right)`$ (12)
$`+`$ $`\sqrt{z}V{\displaystyle \underset{\stackrel{}{k},m}{}}(f_{\stackrel{}{k},m}^{}c_{\stackrel{}{k},m}+c_{\stackrel{}{k},m}^{}f_{\stackrel{}{k},m})`$ (13)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\stackrel{}{k},m}{}}\left(zf_{\stackrel{}{k},m}^{}f_{\stackrel{}{k},m^{}}^{}\mathrm{\Delta }_{\stackrel{}{k},m}+zf_{\stackrel{}{k},m^{}}f_{\stackrel{}{k},m}\mathrm{\Delta }_{}^{}{}_{\stackrel{}{k},m}{}^{}\right)`$ (14)
$``$ $`{\displaystyle \frac{N_s}{2J}}{\displaystyle \underset{m}{}}\mathrm{\Delta }_{}^{}{}_{m}{}^{}\mathrm{\Delta }_m+(ϵ_fϵ_0)(z1)N_s,`$ (15)
where $`N_s`$ denotes the number of lattice sites and $`ϵ_f`$ is the renormalized energy of the $`f`$ orbitals due to the on-site repulsion. The angular momentum projection $`m^{}=m`$ if electron pairing in a singlet state (s- or d-wave) is considered and $`m^{}=m`$ in the case of p-wave pairing. The gap function $`\mathrm{\Delta }_{\stackrel{}{k},m}=\eta _\stackrel{}{k}\mathrm{\Delta }_m`$ and the superconducting order parameter $`\mathrm{\Delta }_m`$ is given by
$$\mathrm{\Delta }_m=\frac{zJ}{N_s}\underset{\stackrel{}{k}}{}\eta _\stackrel{}{k}<f_{\stackrel{}{k},m}f_{\stackrel{}{k},m}>,$$
(16)
where $`\eta _\stackrel{}{k}`$ denotes any of the possible pairing symmetries considered in (11).
The density of the boson condensate $`z`$ minimizes the free energy of the system and $`ϵ_f`$ is obtained after imposing local particle (boson+fermion) conservation at the $`f`$-sites:
$`z`$ $`=`$ $`1{\displaystyle \frac{1}{N_s}}{\displaystyle \underset{\stackrel{}{k},m}{}}<f_{\stackrel{}{k},m}^{}f_{\stackrel{}{k},m}>,`$ (17)
$`ϵ_fϵ_0`$ $`=`$ $`{\displaystyle \frac{V}{2\sqrt{z}N_s}}{\displaystyle \underset{\stackrel{}{k},m}{}}\left(<f_{\stackrel{}{k},m}^{}c_{\stackrel{}{k},m}>+<c_{\stackrel{}{k},m}^{}f_{\stackrel{}{k},m}>\right)`$ (18)
$``$ $`{\displaystyle \frac{N_s}{zJ}}{\displaystyle \underset{m}{}}\mathrm{\Delta }_m^{},\mathrm{\Delta }_m`$ (19)
Equation (17) states that the mean number of electrons at an $`f`$-site is $`1z`$.
In order to derive the gap equation and the spectrum of elementary excitations we use the Gorkov Green’s function approach. The anomalous Green’s functions that we need to consider are:
$`_{f,m}^{}(\stackrel{}{k},\tau \tau ^{})`$ $`=`$ $`<T_\tau f_{\stackrel{}{k},m}^{}(\tau )f_{\stackrel{}{k},m}^{}(\tau ^{})>,`$ (20)
$`_{cf,m}^{}(\stackrel{}{k},\tau \tau ^{})`$ $`=`$ $`<T_\tau c_{\stackrel{}{k},m}^{}(\tau )f_{\stackrel{}{k},m}^{}(\tau ^{})>,`$ (21)
and we must also define three other Matsubara Green’s functions: one that is associated with the conduction electrons, another one for the $`f`$-electrons and the third one is related to the hybridization of the $`f`$ and $`c`$ bands:
$`𝒢_{c,m}(\stackrel{}{k},\tau \tau ^{})`$ $`=`$ $`<T_\tau c_{\stackrel{}{k},m}(\tau )c_{\stackrel{}{k},m}^{}(\tau ^{})>,`$ (22)
$`𝒢_{f,m}(\stackrel{}{k},\tau \tau ^{})`$ $`=`$ $`<T_\tau f_{\stackrel{}{k},m}(\tau )f_{\stackrel{}{k},m}^{}(\tau ^{})>,`$ (23)
$`𝒢_{cf,m}(\stackrel{}{k},\tau \tau ^{})`$ $`=`$ $`<T_\tau c_{\stackrel{}{k},m}(\tau )f_{\stackrel{}{k},m}^{}(\tau ^{})>,`$ (24)
After Fourier transforming these functions into frequency space, we may write down their equations of motion (Gorkov’s equations) according to the Hamiltonian (15):
$`(i\omega _n`$ $`+`$ $`ϵ_f\mu )𝒢_{f,m}(\stackrel{}{k},i\omega _n)+V\sqrt{z}𝒢_{cf,m}(\stackrel{}{k},i\omega _n)`$ (25)
$`+`$ $`Jz^2\mathrm{\Delta }_m(\stackrel{}{k})_{f,m}^{}(\stackrel{}{k},i\omega _n)=1,`$ (26)
$$(i\omega _n+ϵ_\stackrel{}{k}\mu )𝒢_{cf,m}(\stackrel{}{k},i\omega _n)+V\sqrt{z}𝒢_{f,m}(\stackrel{}{k},i\omega _n)=0,$$
(27)
$$(i\omega _nϵ_\stackrel{}{k}+\mu )_{cf,m}^{}(\stackrel{}{k},i\omega _n)V\sqrt{z}_{f,m}^{}(\stackrel{}{k},i\omega _n)=0,$$
(28)
$`(i\omega _n`$ $``$ $`ϵ_f+\mu )_{f,m}^{}(\stackrel{}{k},i\omega _n)V\sqrt{z}_{cf,m}^{}(\stackrel{}{k},i\omega _n)`$ (29)
$`+`$ $`Jz^2\mathrm{\Delta }_m^{}(\stackrel{}{k})𝒢_{f,m}(\stackrel{}{k},i\omega _n)=0.`$ (30)
Diagonalization of the above equations yields the energies of the poles of the Green’s functions (excitation energies) and the corresponding residues (coherence factors). The solutions are of the form
$$𝒢(\stackrel{}{k},i\omega )=\underset{i=1,2}{}\underset{\alpha =\pm }{}\frac{u_i^\alpha }{i\omega _n+\alpha E_i}.$$
(31)
The coherence factors, $`u_i^\alpha `$ and the excitation energies, $`E_i`$ are given in the Appendix.
The Green’s functions have to be determined self-consistently using the mean field equations (16)-(19). These equations can be rewritten in terms of Green’s functions as
$`z`$ $`=`$ $`1{\displaystyle \frac{T}{N_s}}{\displaystyle \underset{\stackrel{}{k},m}{}}{\displaystyle \underset{i\omega _n}{}}𝒢_{f,m}(\stackrel{}{k},i\omega _n),`$ (32)
$`ϵ_fϵ_0`$ $`=`$ $`{\displaystyle \frac{VT}{\sqrt{z}N_s}}{\displaystyle \underset{\stackrel{}{k},m}{}}{\displaystyle \underset{i\omega _n}{}}𝒢_{cf,m}(\stackrel{}{k},i\omega _n)`$ (33)
$``$ $`{\displaystyle \frac{N_s}{zJ}}{\displaystyle \underset{m}{}}\mathrm{\Delta }_m^{}\mathrm{\Delta }_m,`$ (34)
$`\mathrm{\Delta }_m`$ $`=`$ $`{\displaystyle \frac{zJT}{N_s}}{\displaystyle \underset{\stackrel{}{k}}{}}{\displaystyle \underset{i\omega _n}{}}\eta _\stackrel{}{k}_{f,m}(\stackrel{}{k},i\omega _n),`$ (35)
For a given number of particles per site, $`n`$, these equations must be supplemented with the particle conservation condition which yields the chemical potential $`\mu `$ for any temperature:
$$n=1z+\frac{T}{N_s}\underset{\stackrel{}{k},m}{}\underset{i\omega _n}{}𝒢_{c,m}(\stackrel{}{k},i\omega _n).$$
(36)
## III Results
In what follows we consider a cubic lattice in which the conduction band dispersion has the simple tight-binding form:
$$ϵ_\stackrel{}{k}=2t\underset{i=x,y,z}{}\mathrm{cos}(k_i),$$
so that $`D=6t`$ is half the bandwidth. We have used the subroutine hydrd.f from MINPACK in order to solve the four coupled equations (32)-(36).
The possible pairing symmetries expressed in Eq. (11) have been studied separately. The two $`\eta _\stackrel{}{k}`$ functions corresponding to the $`d`$-wave symmetry in equation (11) describe different spatial orientations of the angular momentum of the Cooper pairs and give degenerate solutions. This same remark also applies to the three $`p`$-wave $`\eta _\stackrel{}{k}`$ functions in Eq. (11).
The critical temperatures $`T_c`$ are obtained solving the mean-field equations using the normal state Green’s functions. On the other hand, the study of $`\mathrm{\Delta }(T)`$, $`z(T)`$, $`ϵ_f`$(T), and the specific heat requires solving the mean-field equations with the full Green’s functions. In the normal phase, the slave boson condensation temperature, $`T_z`$, above which $`z=0`$, is given by $`T_z=(ϵ_f\mu )/\mathrm{ln}(N1)`$. If $`N=2`$, $`z`$ is always finite. For larger values of $`N`$, and in particular in the limit $`N\mathrm{}`$, $`z0`$ as the temperature increases. Correspondingly, $`n_f1`$ and the $`f`$-electron superconductivity is inhibited. Therefore, for large values of $`N`$ it is expected that the mean-field theory will not yield superconductivity. One then has to take into account the boson fluctuations. We will focus our attention in the case $`N=2`$, relevant for instance for $`Ce`$ and $`Yb`$ materials, but we will return to this point later.
Figure 1 shows the behavior of the superconducting critical temperatures $`T_c`$ as function of the particle density per channel $`n/N`$, for each of the three pairing symmetries. It is readily seen that the critical temperatures associated with $`d`$ and $`p`$-wave pairing follow similar trends and that the $`d`$-wave symmetry exhibits the highest $`T_c`$ up to densities of about $`n/N0.6`$. At higher densities, a crossover occurs into a regime where the extended $`s`$-wave pairing becomes the most stable, for the parameters considered.
The value of $`T_c`$ vanishes at low densities because the $`f`$-level occupancy also becomes small in that limit ($`z1`$) and Cooper pairing occurs only between the $`f`$-electrons in the model under consideration. In the high density limit, $`T_c`$ vanishes because each $`f`$level is almost fully occupied with one electron ($`z0`$), and freezing of the charge fluctuations \[arising from the term $`f^{}f^{}`$ in (15)\] occurs because of the infinite on-site repulsion.
Heavy-fermion behavior in the normal phase occurs when the chemical potential $`\mu `$ lies close to the peak of the density of states (hence the strong effective mass). This peak is the equivalent of the Kondo resonance peak which appears in the single-impurity problem. For the lattice problem, two strong peaks appear due to hybridization between the conduction electron band and the dispersionless band of localized $`f`$-states, leading to the large electron’s effective mass. For densities above $`n/N0.7`$ the chemical potential becomes close to the density of states peak in the lower band.
In the superconducting phase the full solution of Eqs. (32)-(36) yields a renormalized excitation energy spectrum. In Fig. (2) we show the band structure for $`n/N=0.7`$ in the normal and superconducting phases. It is clear that for this density $`\mu `$ is in the flat region of the band in the normal state.
It is seen from Fig. 1 that as the density per channel $`n/N`$ approaches 1, the value of $`T_c`$ is strongly reduced until it eventually vanishes. From the same figure one can also see that the critical temperature of the $`s`$wave state at $`n/N0.7`$, for instance, is higher than that of the $`d`$ or $`p`$wave states. For the model parameters considered in Fig. 1 this means that as the temperature of a normal system is lowered, the system would first enter a superconducting state with extended s-wave symmetry. On lowering further the temperature, the nature of the superconducting state becomes a mixture of different symmetries. This sequence of phase transitions would be different had we chosen different parameters: our calculations show that if $`J/D`$ is less than about 0.4, then the critical superconducting temperature of a system with $`n/N0.7`$ would correspond to a $`d`$-wave order parameter (see left panel of Fig. 3).
The dependence of $`T_c`$ on the parameters $`V`$, $`ϵ_0`$ and $`J`$ shows interesting crossovers. If $`ϵ_0`$ is well below the chemical potential $`\mu `$ then the $`f`$-level is highly populated and the system cannot become superconducting unless the hybridization parameter $`V`$ is large enough. On the other hand, if $`ϵ_0`$ is not too low a superconducting ground-state is obtained even for small values of $`V`$. As can be seen from the right panel of Fig. 3, $`T_c`$ first increases with $`V`$ up to a maximum value, but as $`V`$ is further increased, large charge quantum fluctuations at the $`f`$ orbitals are induced and superconductivity is destroyed. Moreover, the $`d`$ and $`p`$wave superconductivity seem to be more stable than the $`s`$wave for large values of $`V`$. That this result is consistent with Fig. 1 can be easily understood as follows: upon increasing the hybridization between the $`f`$-orbitals and the conduction band, the electron occupation in the $`f`$-sites is reduced and Fig. 1 already showed that depletion of the $`f`$ band has the effect of reducing $`T_c`$ and increasing the stability of $`d`$-wave pairing relative to $`p`$ and $`s`$wave pairing.
The temperature dependence of the gap function in the superconducting phase is the standard one. In Fig. 4 we show a typical case. The crossing of the $`d`$ and $`p`$order parameters close to $`T_c`$ is related with the same crossing observed in $`z(T)`$. Since close to $`T_c`$, $`z(T)`$ for the superconducting $`d`$phase becomes slightly higher than for the superconducting $`p`$phase, the $`d`$wave phase has an effective superconducting coupling that is slightly higher than the $`p`$wave coupling, leading to an higher $`T_c`$.
In Fig. 5 we show the dependence of $`T_c`$ on the $`f`$-level position. It is seen that the $`d`$wave state has always a higher $`T_c`$ than the $`p`$wave over the range of $`ϵ_0`$ values considered. But the s-wave critical temperature exhibits a much stronger dependence on $`ϵ_0`$. In particular, $`s`$-wave pairing seems to be more strongly depressed for low $`ϵ_0`$.
In a normal system at zero temperature the renormalized $`f`$-level energy $`ϵ_f`$ is located above the chemical potential and $`ϵ_f\mu `$ is of the order of the Kondo temperature for the equivalent single-imputity problem. Keeping the particle density fixed, both $`ϵ_f`$ and $`\mu `$ increase with temperature but the difference $`ϵ_f\mu `$ decreases. Our calculations show that $`T_c`$ is smaller than $`ϵ_f\mu `$ by a factor of about 10 (see Fig 6) over almost the entire range of densities considered in Fig. 1. In Fig. 6 we present $`T_c`$, $`ϵ_f`$, and $`\mu `$ for the extended $`s`$wave order parameter (the curves for the other symmetries are qualitatively the same). The susceptibility $`dn_f/dϵ_f`$, in the region of densities characterized by $`n/N0.7`$ or larger, is very small since the $`f`$-level density of states is much larger than the $`c`$-level one, leading to a negative feedback changing the chemical potential in such a way as to keep $`ϵ_f`$ close to $`\mu `$. This is very clear from Fig. 6, where $`ϵ_f`$ is indeed close to $`\mu `$, for electronic densities where the density of states is large. This is consistent with the picture that the pairing is developed by the excitations of the system resulting from the Kondo compensated lattice.
Finally we calculate the specific heat for the various symmetries. The non-conventional pairing symmetry leads to a power law behavior at low $`T`$ in the superconducting phase. In Fig. 7 we show the specific heat for the various symmetries as a function of temperature. The specific heat jump at the transition is $`\mathrm{\Delta }C/C1.6,1.3,0.8`$ for the $`p,d,s`$symmetries respectively. We have found that the specific heat at low $`T`$ has a $`T^2`$ dependence for the $`p,d`$ symmetries, and has an exponential behavior for the $`s`$wave case.
Considering now the effect of increasing the number of channels $`N`$, we find that $`T_c`$ decreases by one order of magnitude or more, as $`N`$ changes from $`N=2`$ to $`N=4`$. For the parameters considered in the figures, the effect is most dramatic for $`s,p`$wave symmetries, where superconductivity is absent for $`N4`$. Furthermore, we found that the critical temperature of a system with a $`d`$wave order parameter is less sensitive to the number of channels $`N`$, as compared to the other symmetries.
## IV Summary
Heavy-fermions show a rich and complex behavior at low temperatures. In particular, the interplay between magnetic correlations, the Kondo effect and superconducting correlations is a difficult problem to solve. This is further complicated since neither the mechanism nor the pairing symmetry are fully established. In this paper we have focused on the superconducting order assuming that the superconducting correlations are the dominant ones. Using a generalized Anderson lattice model with nearest-neighbor attraction between the $`f`$-electrons and with infinite-$`U`$ local Coulomb repulsion, we studied the various pairing symmetries using a mean-field approach. In this way it is possible to compare the various solutions in contrast to an approach where, starting from the normal phase, the leading instabilities are identified.
The results show that there are several crossovers between the $`s`$, $`d`$ and $`p`$-wave pairing symmetries as the parameters of the model are varied. In contrast to a previous mean-field approach we find superconducting order, even though $`U=\mathrm{}`$. The reason is that we focus on a regime where $`0<n_f<1`$, while the previous work concentrated on a regime where $`1<n_f<2`$ (for finite $`U`$). Since we consider only the case $`U=\mathrm{}`$, $`n_f`$ has to be smaller than one due to the Coulomb repulsion. In the previous work as $`U`$ grows the density $`n_f1`$ the $`f`$-electrons become more localized inhibiting superconductivity. We find the same qualitative behavior as we approach the Kondo regime from the mixed valent regime. For small values of $`ϵ_0`$ we tend to a regime where $`n_f1`$ and superconductivity is suppressed.
In the mean-field approach if $`z0`$ ($`n_f1`$) the gap function $`\mathrm{\Delta }_m0`$. This happens for large densities $`n/N1`$. For $`n_f0`$ superconductivy is supressed, since the superconducting coupling is among the $`f`$electrons. Also, if $`N`$ is large $`z0`$ at lower temperatures. In particular, $`p`$-wave and extended $`s`$-wave symmetries are strongly suppressed. For $`N=2`$ $`z`$ is always finite. For larger values of $`N`$ in general it will be necessary to consider the boson fluctuations and a treatment beyond mean-field will be required. For systems where the spin degeneracy is low we expect the results to be qualitatively correct.
We have found that the $`d`$-wave and $`p`$-wave symmetries yield similar transition temperatures. For large nearest-neighbor attraction the extended $`s`$-wave pairing is preferred. Otherwise, the $`d`$-wave symmetry seems to be more robust, in particular as $`N`$ grows. Clearly, we are not considering magnetic correlations in our mean-field study and therefore the description applies to systems where there are no local moments (and therefore $`T_c<T_K`$) and where $`T_c>T_{RKKY}`$.
We found that superconductivity is preferred in a mixed valent regime (due to the infinite Coulomb repulsion). There are materials that are mixed valent and superconductors . In the framework of weak coupling BCS theory one would expect that the local magnetic character of the $`f`$-states should be pair breaking. However, the heavy fermion superconductivity in the Kondo limit (integer valent case) reveals that the pairing is of another nature that compensates the pair breaking effects of the local magnetic character. For materials such as $`CeRu_3Si_2`$ there is a considerable mixed valent character and accordingly the effective masses are not high. Also, the Wilson ratio is close to one indicating a conventional weak-coupling BCS superconductor. Other mixed valent superconductors are not conventional superconductors. It would be interesting to identify systems that by changing the mixed-valent character could change the superconducting temperature, $`T_c`$. In the framework of our model this would require $`f`$-states with large $`U`$ values. The nearest-neighbor attraction could be due to several mechanisms like spin fluctuations or slave boson fluctuations (Coulombic nature).
## V Acknowledgements
We would like to thank P. Estrela for bringing Refs. to our attention. This research was supported by PRAXIS under grant number 2/2.1/FIS/302/94.
## Poles and coherence functions for the Green’s functions
The algebraic solutions of Eqs. (26-30) for the Green’s functions $`𝒢_{f,m}(\stackrel{}{k},i\omega _n)`$, $`𝒢_{cf,m}(\stackrel{}{k},i\omega _n)`$, $`_{f,m}^{}(\stackrel{}{k},i\omega _n)`$, and $`𝒢_{c,m}(\stackrel{}{k},i\omega _n)`$ have the form
$$𝒢(\stackrel{}{k},i\omega _n)=\underset{i=1,2}{}\underset{\alpha =\pm }{}\frac{u_i^\alpha }{i\omega _n+\alpha E_i},$$
(37)
and the coherence factors $`u_i^\alpha `$ and the excitations energies $`E_i`$ are given below.
The energies $`E_i`$ have the form
$`E_1=\sqrt{\gamma /2\sqrt{\gamma ^2/4\beta }},`$ (38)
$`E_2=\sqrt{\gamma /2+\sqrt{\gamma ^2/4\beta }},`$ (39)
with $`\gamma `$ and $`\beta `$ given by
$`\gamma =(ϵ_f\mu )^2+(ϵ_\stackrel{}{k}\mu )^2+2V^2z+|Jz^2\mathrm{\Delta }(\stackrel{}{k})|^2,`$ (40)
$`\beta =[(ϵ_\stackrel{}{k}\mu )(ϵ_f\mu )V^2z]^2+|Jz^2\mathrm{\Delta }(\stackrel{}{k})|^2(ϵ_\stackrel{}{k}\mu )^2.`$ (41)
The $`u_i^\alpha `$ factors for $`𝒢_{f,m}(\stackrel{}{k},i\omega _n)`$ are given by
$`u_1^+=F(E_1+ϵ_\stackrel{}{k}\mu )X_1,`$ (42)
$`u_1^{}=F(E_1ϵ_\stackrel{}{k}+\mu )Y_1,`$ (43)
$`u_2^+=G(E_2+ϵ_\stackrel{}{k}\mu )X_2,`$ (44)
$`u_2^{}=G(E_2ϵ_\stackrel{}{k}+\mu )Y_2,`$ (45)
where the functions $`X_i`$ and $`Y_i`$ ($`i=1,2`$) are given by
$`X_i=(ϵ_\stackrel{}{k}\mu )(ϵ_f\mu )`$ $``$ $`(ϵ_\stackrel{}{k}ϵ_f2\mu )E_i`$ (46)
$`+`$ $`E_i^2zV^2,`$ (47)
$`Y_i=(ϵ_\stackrel{}{k}\mu )(ϵ_f\mu )`$ $`+`$ $`(ϵ_\stackrel{}{k}ϵ_f2\mu )E_i`$ (48)
$`+`$ $`E_i^2zV^2,`$ (49)
and the functions $`F`$ and $`G`$ are given by
$`F={\displaystyle \frac{1}{2E_1(E_2^2E_1^2)}},G={\displaystyle \frac{1}{2E_2(E_2^2E_1^2)}}.`$ (50)
The $`u_i^\alpha `$ factors for $`_{f,m}^{}(\stackrel{}{k},i\omega _n)`$ are given by
$`u_1^+=Jz^2\mathrm{\Delta }(\stackrel{}{k})F(E_1^2(ϵ_\stackrel{}{k}\mu )^2),u_1^+=u_1^{},`$ (51)
$`u_2^+=Jz^2\mathrm{\Delta }(\stackrel{}{k})G(E_2^2(ϵ_\stackrel{}{k}\mu )^2),u_2^+=u_2^{}.`$ (52)
The $`u_i^\alpha `$ factors for $`𝒢_{cf,m}^{}(\stackrel{}{k},i\omega _n)`$ are given by
$`u_1^+=V\sqrt{z}FX_1,u_1^{}=FY_1,`$ (53)
$`u_2^+=V\sqrt{z}GX_2,u_2^{}=GY_2.`$ (54)
The $`u_i^\alpha `$ factors for $`𝒢_{c,m}(\stackrel{}{k},i\omega _n)`$ are given by
$`u_1^+=FQ_1,u_1^{}=FR_1,`$ (55)
$`u_2^+=GQ_2,u_2^{}=GR_2,`$ (56)
where
$`Q_i=(ϵ_\stackrel{}{k}\mu E_i)|Jz^2\mathrm{\Delta }(\stackrel{}{k})|^2+(ϵ_f+E_i)X_i,`$ (57)
$`R_i=(ϵ_\stackrel{}{k}\mu +E_i)|Jz^2\mathrm{\Delta }(\stackrel{}{k})|^2(ϵ_f+E_i)Y_i.`$ (58)
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# Exact solutions for a family of discretely spiked harmonic oscillators
## 1 Introduction
The name *spiked oscillator* generally refers to Hamiltonians with potentials containing a parabolic term $`x^2`$ describing ordinary quantum oscillator, and a singular term $`1/x^\alpha `$ forming a sharp spike at $`x=0`$. A typical Hamiltonian of a spiked oscillator is of the form given by Harrell:
$$H=\frac{d}{dx^2}+x^2+\frac{\lambda }{x^\alpha }0x<\mathrm{}.$$
It has found variety of uses in atomic, molecular, nuclear and particle physics since it provides the simplest realistic model of interaction potentials due to its repulsive core $`x^\alpha `$.
Some forms of spiked oscillators have exact solutions, such as the one described by one dimensional Hamiltonian, examined four decades ago by Goldman and Krivchenkov
$$H=\frac{d}{d^2}+V_0(\frac{a}{x}\frac{x}{a})^2$$
whose exact solutions $`\varphi _n(x)[0,\mathrm{})`$ satisfy Dirichlet’s boundary condition $`\varphi _n(0)=0`$.
Models with exact solutions can serve as starting points for analyses of related cases where exact solutions do not exist and where a perturbative analysis is required. A good example can be found in the papers of Hall et al, who slightly modified the above Hamiltonian, computed a basis given by a set of exact orthonormal eigenfunctions $`\varphi _n(x)`$ for $`L^2(0,\mathrm{})`$, and then used this basis in variational analysis of the Hamiltonian
$$H=\frac{d}{dx^2}+Bx^2+\frac{A}{x^2}+\frac{\lambda }{x^\alpha }.$$
They argued that their singular basis forms much better starting point for perturbation techniques than the basis of the ordinary harmonic oscillator, as used in earlier work of Aguilera-Navarro et al, who had applied a perturbative variational analysis to the lowest eigenvalue of the Harrel’s Hamiltonian.
In this context, we present an analysis of family of following spiked Hamiltonians
$$H_l=\frac{d^2}{dx^2}+x^2+\frac{l(l+1)}{x^2},$$
(1)
parametrized by discrete values of $`l=0,1,2..`$. Depending on interpretation of variable $`x`$, this form can either describe a family of one dimensional linear oscillators, or a family of $`N`$ dimensional radial oscillators for odd $`N=1,3,5..`$ What is traditionally known as the three dimensional isotropic radial oscillator, with known exact solutions, can be reduced to a special case of $`N=1`$. We shall show that solutions for radial oscillators in even dimensions $`N=2,4,6..`$ do not form well defined orthogonal bases and, as such, they are of no practical interest.
The solutions we provide are exact and much simpler than any of those described in this introduction, and therefore they are good candidates for orthonormal bases that could be used in perturbative analyses of oscillators with spiked potentials other than $`x^2`$.
The factorization formalism developed here is completely algebraic and somehow related to methods of *Superstring Quantum Mechanics (SUSY QM)*. Modern algebraic factorization methods are overviewed in the paper of Rosas-Ortiz and closely related SUSY methods are described for example in introductory parts of the papers by Junker and Roy , and Lévai and Roy.
The common premise of all those methods is the use of algebraic recursive manipulations on two Hamiltonians: one with known spectrum and known eigenvectors and another which is isospectral to the former but whose eigenvectors are to be found. The methodology developed here uses the same premise, but goes a bit further by manipulating not just two Hamiltonians but the uncountable set of related Hamiltonians (1). On the other hand, the SUSY and similar methods rely on general solution of underlying *Ricatti* equation, which is completely ignored in the approach presented in this paper.
Quantum models with exact solutions play an important role and there is a renewed interest in enrichment of the traditional narrow set of exactly solvable models, such as harmonic oscillator, hydrogen-like potential, Morse potential and square well potential. It is only fair to mention other related methods, which admit certain type of non-hermitian Hamiltonians with complex-valued potentials giving rise to real energy spectrum. For example, Cannata et al use Darboux method to derive a general class of complex potentials with energy spectra identical to that of regular harmonic oscillator. Bender et al replaces the requirement of Hamiltonian hermiticity by weaker requirement of *PT-symmetry*, which is satisfactory condition for real energy spectra. They examine complex harmonic and anharmonic oscillators, complex square wells, etc. This avenue is being explored by many other researchers.
To complete this introduction, we just mention in passing that exactly solvable problems can be categorized in one of three classes: *Exactly solvable (ES), Quasi exactly solvable (QES) and Conditionally exactly solvable (CES)*.
In the next section we present an overview of our findings, while the remaining sections provide specific details.
## 2 Outline
The first Hamiltonian, $`H_0`$, from the family of discretely spiked oscillators (1) describes the well known harmonic oscillator, without the spike, of known discrete energy spectrum and known ladder operators $`a_0`$ and $`a_0^{}`$. Aside from this simple case with parabolic potential, the potentials of remaining members of this family have strong discontinuities at $`x=0`$ and do not resemble the first case at all. Yet they all are isospectral; that is, they all have the discrete energy spectra identical to that of the simple harmonic oscillator. The only difference is in the definition of their ground state energies, $`E_{l,0}=2l+1`$, in units of $`\frac{1}{2}\omega \mathrm{}`$.
We have chosen to denotate the coresponding ground states as $`|0,0`$, $`|1,0`$, $`|2,0`$, etc. for $`l=0,1,2..`$, respectively. Table 1 visualizes the relationship between energy states and energies for the entire family of discretely spiked oscillators.
For each Hamiltonian $`H_l`$ a pair of intertwining operators can be defined
$`b_{l+1}^{}|l,k`$ $``$ $`|l+1,k`$
$`b_{l+1}|l+1,k`$ $``$ $`|l,k`$
which connect eigenstates of two ‘neighbours’, $`H_l`$ and $`H_{l+1}`$. Both represent actions coresponding to diagonal traversal of the Table 1, in NE and SW directions, respectively.
In addition to intertwining operators, we can also define the ladder operators $`a_l`$ and $`a_l^{}`$ that allow for traversing of the Table 1 in vertical direction. The annihilation operator changes the state $`|l,k`$ with energy $`E_{l,k}`$ to the state $`|l,k1`$ of lower energy $`E_{l,k1}`$
$`a_l|l,0`$ $`=`$ $`0`$
$`a_l|l,k`$ $`=`$ $`\alpha _l|l,k1.`$
Similarly, the creation operator changes the state $`|l,k`$ with energy $`E_{l,k}`$ to the state $`|l,k+1`$ of higher energy $`E_{l,k+1}`$
$$a_l^{}|l,k=\alpha _l^{}|l,k+1$$
where $`\alpha _l`$ and $`\alpha _l^{}`$ are the normalization constants.
Known ladder operators $`a_0`$ and $`a_0^{}`$ for the classical harmonic oscillator are used to recursively prove the existence of similar operators for the remaining Hamiltonians $`H_l`$. The are given by formulas
$$a_{l+1}^{}=b_{l+1}^{}a_l^{}b_{l+1}$$
$$a_{l+1}=b_{l+1}^{}a_lb_{l+1}$$
The condition $`a_l|l,0=0`$ allows to generate explicit form of the state $`|l,0`$ in position representation; that is the vaweform function $`\varphi _{l,0}(x)`$ of the ground state. For $`l=0`$, this is the bell-shaped function, $`\varphi _{0,0}\mathrm{exp}(\frac{x^2}{2})`$, subject to normalization. Having computed the ground waveform $`\varphi _{l,0}`$, one can use the creation operator $`a_l^{}`$ to recusively generate all other eigenstates $`|l,k`$, or coresponding waveforms $`\varphi _{l,k}(x)`$.
This seemingly completes the scheme since, thanks to the pairs of ladder and intertwining operators, the entire Table 1 can be created – starting from the left lower corner and zigzagging one step a time, with or without the help of ladder operators $`a_1^{}`$, $`a_2^{}`$, etc.
However, we must also require that each set of state vectors is a basis in its own space; that is, they all are mutually orthogonal and normalized
$$l,j|l,k=\delta _{jk}.$$
But some of the wavefunctions are singular at $`x=0`$, since they contain factors $`\frac{1}{x}`$, $`\frac{1}{x^2}`$, etc. Depending on whether the Hamiltonians (1) describe linear or radial oscillators and, in the latter case, depending on the space dimension, some such functions are not square-integrable and therefore they must be rejected as unphysical. Those wave functions that are physical must additionally pass the test of orthogonality in order to qualify as members of useful bases. To take it all into account we therefore specialize the generic scheme and draw additional conclusions specific to dimensionality of the problem. All of this is supported by a Haskell program, which tests the theory and provides useful and flexible computational tool.
In the following sections we shall present a detailed account of this outline.
## 3 Operators
We shall define two pairs of state generation operators: intertwining operators, affecting the quantum number $`l`$, and ladder operators affecting the qunatum number $`k`$. The former correspond to diagonal traversal of Table 1 and the latter – for the vertical traversal.
### 3.1 Intertwining operators
Let’s define two operators $`b_l`$ and its adjoint $`b_l^{}`$
$`b_l`$ $`=`$ $`{\displaystyle \frac{d}{dx}}+\beta _l(x),`$ (2)
$`b_l^{}`$ $`=`$ $`{\displaystyle \frac{d}{dx}}+\beta _l(x)\text{where}`$ (3)
$`\beta _l(x)`$ $`=`$ $`x+{\displaystyle \frac{l}{x}}.`$ (4)
Function $`\beta (x)`$ and operator $`\frac{d}{dx}`$ do not commute
$$[\beta _l,\frac{d}{dx}]=\beta _l^{}(x)=\frac{l}{x^2}1,$$
and, as a consequence, the commutator $`[b_l,b_l^{}]`$ is not zero either, since
$`b_l^{}b_l`$ $`=`$ $`{\displaystyle \frac{d^2}{dx^2}}+\beta _l^2(x)\beta _l^{}(x)=H_l+(2l1),`$ (5)
$`b_lb_l^{}`$ $`=`$ $`{\displaystyle \frac{d^2}{dx^2}}+\beta _l^2(x)+\beta _l^{}(x)=H_{l1}+(2l+1),`$ (6)
where
$$H_1=H_0=\frac{d^2}{dx^2}+x^2,$$
and where the remaining Hamiltonians are defined by (1).
With the help of the above equations we can establish two very useful formulas, which we will use later:
$`H_lb_l^{}b_l^{}H_{l1}`$ $`=`$ $`2b_l^{},`$ (7)
$`H_{l1}b_lb_lH_l`$ $`=`$ $`2b_l.`$ (8)
These are quasi commutative relations, connecting two “neighbouring” Hamiltonians. Let us now consider specific cases of $`\beta _l(x)`$.
#### 3.1.1 Ordinary oscillator
The case $`l=0`$ is very special since
$`b_0^{}b_0`$ $`=`$ $`H_01,`$
$`b_0b_0^{}`$ $`=`$ $`H_1+1`$
$`=`$ $`H_0+1.`$
Both equations refer to the same Hamiltonian of ordinary harmonic oscillator, and the operators $`b_0`$ and $`b_0^{}`$ degenerate to the familiar ladder operators
$`a_0^{}`$ $`=`$ $`b_0^{}={\displaystyle \frac{d}{dx}}+x,`$ (9)
$`a_0`$ $`=`$ $`b_0={\displaystyle \frac{d}{dx}}+x,`$ (10)
which have these important properties:
$`[a_0,a_0^{}]`$ $`=`$ $`2,`$ (11)
$`[H_0,a_0]`$ $`=`$ $`2a_0,`$ (12)
$`[H_0,a_0^{}]`$ $`=`$ $`2a_0^{}.`$ (13)
Consequently, the energies are given by $`E_{0,k}=2k+1`$ for $`k=0,1,2`$…, in units of $`\frac{1}{2}\omega \mathrm{}`$. We will skip other details, since this is the very known case.
#### 3.1.2 Discretely spiked oscillators
The interesting part begins here. First, let us prove that $`b_{l+1}^{}|l,k`$ is an eigenvector of Hamiltonian $`H_{l+1}`$ by assuming that this eigen equation holds
$$H_l|l,k=E_{l,k}|l,k$$
and making use of formula (7):
$`H_{l+1}|b_{l+1}^{}(l,k)`$ $`=`$ $`H_{l+1}b_{l+1}^{}|l,k`$
$`=`$ $`b_{l+1}^{}H_l|l,k+2b_{l+1}^{}|l,k`$
$`=`$ $`(E_{l,k}+2)|b_{l+1}^{}(l,k).`$
The operator $`b_{l+1}^{}`$ transforms the state $`|l,k`$ with energy $`E_{l,k}`$ into the state $`|l+1,k`$ with energy $`E_{l+1,k}=E_{l,k}+2`$
$$|l+1,k=\frac{1}{\sqrt{2k+4l+4}}b_{l+1}^{}|l,k.$$
(14)
This coresponds to the diagonal move in NE direction in Table 1. One can verify the normalization factor in equation (14) by comparing the norm of its both sides, making use of (6) and noticing that the energy of l-th oscillator in the state $`|l,k`$ is given by
$$E_{l,k}=2(l+k)+1.$$
(15)
The latter can be easily established with the help of Table 1.
Similarly, we can prove that state $`b_{l+1}|l+1,k`$ is the eigenvector of Hamiltonian $`H_l`$:
$`H_l(b_{l+1}|l+1,k)`$ $`=`$ $`(b_{l+1}H_{l+1}2b_{l+1})|l+1,k)`$
$`=`$ $`(E_{l+1,k}2)(b_{l+1}|l+1,k).`$
The operator $`b_{l+1}`$ transforms the state $`|l+1,k`$ with energy $`E_{l+1,k}`$ into the state $`|l,k`$ with energy $`E_{l,k}=E_{l+1,k}2`$
$$|l,k=\frac{1}{\sqrt{2k+4l+4}}b_{l+1}|l+1,k.$$
(16)
This coresponds to the diagonal move in SW direction in Table 1.
### 3.2 Ladder operators
We shall now define two operators
$$a_{l+1}^{}=b_{l+1}^{}a_l^{}b_{l+1}$$
(17)
and
$$a_{l+1}=b_{l+1}^{}a_lb_{l+1}$$
(18)
and prove that they indeed have the expected properties of ladder operators. We shall first prove by induction, that
$$[H_l,a_l^{}]=2a_l^{}$$
(19)
This holds true for the case $`l=0`$, as seen in (13). Assuming now, that the above is true for $`l1`$, the proof is as follows:
$`H_la_l^{}`$ $`=`$ $`H_lb_l^{}a_{l1}^{}b_l`$
$`\text{from}\text{17}`$
$`=`$ $`(b_l^{}H_{l1}+2b_l^{})a_{l1}^{}b_l`$
$`\text{from}\text{7}`$
$`=`$ $`b_l^{}(a_{l1}^{}H_{l1}+2a_{l1}^{})b_l+2a_l^{}`$
$`\text{from}\text{19}\text{ by induction}`$
$`=`$ $`b_l^{}a_{l1}^{}H_{l1}b_l+4a_l^{}`$
$`=`$ $`b_l^{}a_{l1}^{}(b_lH_l2b_l)+4a_l^{}`$
$`\text{from}\text{8}`$
$`=`$ $`b_l^{}a_{l1}^{}b_lH_l+2a_l^{}`$
$`=`$ $`a_l^{}H_l+2a_l^{}`$
$`\text{from}\text{17}`$
Similarly, we can prove by induction that the following holds
$$[H_l,a_l]=2a_l.$$
(20)
It immediately follows from (19) and (20) that operators $`a_l^{}`$ and $`a_l`$ create and annihilate eigenstates of Hamiltonian $`H_l`$, respectively
$`H_la_l^{}|l,k`$ $`=`$ $`a_l^{}H_l|l,k+2a_l^{}|l,k`$
$`=`$ $`(E_{l,k}+2)a_l^{}|l,k,`$
$`H_la_l|l,k`$ $`=`$ $`a_lH_l|l,k2a_l|l,k`$
$`=`$ $`(E_{l,k}2)a_l|l,k.`$
Finally, we can prove by induction that if $`a_0|0,0=0`$ then similar relation holds for any other oscillator ($`l>0`$):
$$a_l|l,0=0.$$
## 4 Wavefunctions
### 4.1 Haskell module Spike
The remaining part of this paper refers to results obtained with the help of a little program,*Spike.hs*,<sup>1</sup><sup>1</sup>1The program is available at http://www.numeric-quest.com/haskell/Spike.hs. written in functional lazy language Haskell<sup>2</sup><sup>2</sup>2See home page of Haskell at http://www.haskell.org.. The program is used for generation of wavefunctions, for normalization of solutions in n-dimensional spaces and, generally, for testing the theory presented in the previous sections. The program performs some algebraic manipulation of Laurent series, such as computation of derivatives, integrals, sums and products, and directly implements raising and intertwining operators previously discussed. It is best used in interpreting environment<sup>3</sup><sup>3</sup>3Haskell interpreter Hugs is available for free at http://www.haskell.org/hugs/. It runs on any major platform., since one can evaluate any formula from module *Spike* in any order of one’s choice. Module *Spike* relies on other standard Haskell modules and on our module *Fraction.hs* <sup>4</sup><sup>4</sup>4Module Fraction is avauilable at http://www.numeric-quest.com/haskell/Fraction.hs.. For interesting introduction to Haskell as a coding tool for scientific applications see.
### 4.2 Laurent series
Any analytic function $`f(z)`$ in an open annulus $`r<|zz_0|<R`$ can be expressed as the sum of two series
$$f(z)=\underset{j=0}{\overset{\mathrm{}}{}}a_j(zz_0)^j+\underset{j=1}{\overset{\mathrm{}}{}}a_j(zz_0)^j.$$
Such an expansion, containing negative as well as positive powers is called *Laurent series* for $`f(z)`$ in this annulus. In Haskell we will represent it as a pair of two lists
$$f=L[a_0,a_1,a_2,\mathrm{}][a_1,a_2,\mathrm{}],$$
where coefficients of expansions can be any numbers belonging to class *Num*, such as *Integer, Double, Complex Double*, etc.
The following defines the type of data structure representing Laurent series
```
data Laurent a = L [a] [a],
```
where $`a`$ is a type variable of some numeric kind. It is convenient to choose it as $`Integer`$ (of unlimited size) since this leads to highly accurate results. For computations requiring normalization the data *Fraction* will be used instead of traditional *Double* since this guarantees very good accuracy, with relatively small performance loss.
As it will soon become clear, wavefunctions of spiked oscillators can be represented as products of two functions
$$\varphi _{l,k}(x)=f_{l,k}(x)e^{x^2/2}$$
(21)
where $`f_{l,k}(x)`$ is a real function of $`x`$ and can be expanded as Laurent series of type $`LaurentInteger`$.
### 4.3 One step operators
Application of the intertwining operator to function (21)
$`b_l^{}\varphi _{l,k}(x)`$ $`=`$ $`b_l^{}(f_{k,l}(x)e^{x^2/2})`$
$`=`$ $`(f_{k,l}^{}(x)+2xf_{k,l}(x)+{\displaystyle \frac{l}{x}}f_{k,l}(x))e^{x^2/2}`$
$`=`$ $`(f_{k,l}^{}(x)+p_l(x)f_{k,l}(x))e^{x^2/2},`$
results in manipulation of two Laurent series $`p_l`$ and $`f_{k,l}`$, where $`p_l(x)`$ is a short Laurent series, represented in Haskell as $`L[0,2][l]`$. The expression $`(f_{k,l}^{}(x)+p_l(x)f_{k,l}(x)`$ is implemented in Haskell as function *b’*. It takes as an input an integer $`l`$ and a Laurent series $`f_{l,k}`$ of type $`LaurentInteger`$, and then produces another Laurent series by invoking few primitives, such as multiplication, addition and differentiation defined for such series. Notice the *gcd normalization* of the results, which keeps the size of integers in check and also provides a canonical standard for comparison of different methods of computations.
```
b’ :: Integer -> Laurent Integer -> Laurent Integer
b’ l f = gcdNormalize (p * f - diff f)
where
p = L [0,2] [l]Ψ
```
Implementation $`b^{}l`$ of operator $`b_l^{}`$ is the only Haskell function needed to generate all eigenfunctions for all spiked harmonic oscillators. But to complement the theory developed in the previous sections, below are the remaining operators. Firstly, here is function $`bl`$ coresponding to intertwining operator $`b_l`$
```
b :: Integer -> Laurent Integer -> Laurent Integer
b l f = gcdNormalize (p * f + diff f)
where
p = L [] [l].
```
Next is the implementation of the rasing operator $`a^{}`$:
```
a’ :: Integer -> Laurent Integer -> Laurent Integer
a’ 0 = b’ 0
a’ l = b’ l . a’ (l - 1) . b l
```
Notice that function $`a^{}\mathrm{\hspace{0.17em}0}=b^{}\mathrm{\hspace{0.17em}0}`$ is a special case. It coresponds to the raising operator $`a_0^{}=b_0^{}`$ – exactly as specified in the theoretical sections. Finally, here is the implementation of the lowering ladder operator $`a_l`$:
```
a :: Integer -> Laurent Integer -> Laurent Integer
a 0 = b 0
a l = b’ l . a (l - 1) . b l.
```
### 4.4 Cumulative operators
Since the one-step operators from the previous section are to be applied recursively to some base functions it is convenient to define cumulative versions of these operators.
The cumulative intertwining operator,*twine’ l*, is a composition of (NE) operators $`b_lb_{l1}\mathrm{}b_1`$
```
twine’ :: Integer -> Laurent Integer -> Laurent Integer
twine’ 0 = id
twine’ 1 = b’ 1
twine’ l = b’ l . twine’ (l - 1)
```
The cumulative raising ladder operator, *ladder’ l k*, is a composition of $`k`$ raising ladder operators $`a_l^{}a_l^{}a_l^{}\mathrm{}`$
```
ladder’ :: Integer -> Integer -> Laurent Integer
-> Laurent Integer
ladder’ l 0 = id
ladder’ l 1 = a’ l
ladder’ l k = a’ l . ladder’ l (k - 1).
```
### 4.5 Generating wavefunctions
Using these tools we can easily generate sets of eigenfunctions for spiked oscillators. When $`l=0`$ the first wavefunction is given by $`e^{x^/2}`$ and that implies that $`f_{0,0}`$ is simply $`1`$, and its corresponding Laurent series is $`L[1][]`$. Here $`[]`$ represents an empty list – meaning no negative powers at all. By executing
```
iterate (a’ 0) (L [1] [])
```
we can generate infinite list of gcd-normalized eigenfunctions for ordinary harmonic oscillator. If we want to retain six, say, wave functions we need to cut this infinite list down to six elements
```
take 6 (iterate (a’ 0) (L [1] [])).
```
A slightly prettiefied output of the above looks like this:
```
[ -- Meaning:
L [1] [], -- 1
L [0,1] [], -- x
L [-1,0,2] [], -- -1 + 2x^2
L [0,-3,0,2] [], -- -3x + 2x^3
L [3,0,-12,0,4] [], -- 3 - 12x^2 + 4x^4
L [0,15,0,-20,0,4] [] -- 15x - 20x^3 + 4x^5.
]
```
It is evident that these solutions do not contain negative powers of $`x`$. In fact they represent (gcd normalized) Hermite polynomials in disguise, and this is exactly what one would expect from ordinary harmonic oscillator. To test implementation of the lowering ladder operator $`a_0`$, execute
```
take 6 (iterate (a 0) (L [0,15,0,-20,0,4])).
```
This should recreate the above list in reverse. In addition, the expression
```
a 0 (L [1] [])
```
generates a special kind of Laurent series
```
L [] [],
```
which is the Laurent series for *function zero*.
To find eigenfunctions of the first spiked oscillator ($`l=1`$) we have two choices. Firstly, we can apply operator $`b_1^{}`$, or function $`b^{}\mathrm{\hspace{0.17em}1}`$ to every eigenfunction of ordinary oscillator, by executing
```
map (b’ 1) (take 6 $ iterate (a’ 0) (L [1] [])).
```
Alternatively, we can iteratively apply function $`a^{}\mathrm{\hspace{0.17em}1}`$ to the first eigenfunction $`f_{1,0}=b\mathrm{\hspace{0.17em}1}(L[1][])`$ of the oscillator $`l=1`$
```
take 6 $ iterate (a’ 1) f where f = b’ 1 (L [1] []).
```
Either way, the output is as follows
```
[ -- Meaning:
L [0,2] [1], -- 2x + 1/x
L [0,0,1] [], -- x^2
L [0,-4,0,4] [-1], -- -4x + 4x^3 - 1/x
L [0,0,-5,0,2] [], -- -5x^2 + 2x^4
L [0,18,0,-36,0,8] [3], -- 18x - 36x^3 + 8x^5 + 3/x
L [0,0,35,0,-28,0,4] [] -- 35x^2 - 28x^4 + 4x^6
]
```
This time, every other solution contains term with $`1/x`$, and as such it might or might not be acceptable as physical solution, depending on the space dimension under consideration. More about it in later sections.
To generate similar list for the second spiked oscillator ($`l=2`$) we need to apply two functions $`b^{}\mathrm{\hspace{0.17em}1}`$ and $`b^{}\mathrm{\hspace{0.17em}2}`$ in succession to a list of eigenfunctions of ordinary harmonic oscillator
```
map (b’ 2 . b’ 1) (take 6 $ iterate (a’ 0) (L [1] [])),
```
or, alternatively, apply the raising ladder function $`a^{}\mathrm{\hspace{0.17em}2}`$ to the lowest eigenfunction of the oscillator $`l=2`$
```
take 6 $ iterate (a’ 2) f where f = (b’ 2 . b’ 1) (L [1] []).
```
Both methods lead to the same output:
```
[
L [4,0,4] [0,3],
L [0,0,0,1] [],
L [-6,0,-12,0,8] [0,-3],
L [0,0,0,-7,0,2] [],
L [24,0,72,0,-96,0,16] [0,9],
L [0,0,0,63,0,-36,0,4] []
].
```
Here is a general formula for generation of infinite lists of wavefunctions for spiked harmonic oscillators $`l=0,1,2..`$
```
wavefunctions :: Integer -> Laurent Integer
wavefunctions l = map (twine’ l) $ iterate (a’ 0) (L [1] [])
```
To generate just one eigenfunction $`f_{l,k}`$ the following Haskell function can be used
```
wavefunction :: Integer -> Integer -> Laurent Integer
wavefunction l k = twine’ l (ladder’ 0 k (L [1] []))
```
## 5 Normalization and orthogonalization
We have shown that the solutions for spiked oscillators are of the form $`\varphi _{l,k}=f_{l,k}(x)e^{x^2/2}`$, where $`f_{l,k}(x)`$ are real functions of $`x`$ and, specificly, they are some Laurent series – admitting both positive and negative powers of $`x`$. For a specific case of ordinary oscillator, these are the Hermite polynomials, with positive powers only.
The scheme presented so far accepts all solutions, but does not check whether they are square-integrable and, in addition, whether they obey the orthogonality relations
$$l,i|l,j=\delta _{i,j}.$$
Due to presence of the term $`\frac{1}{x}`$ in the definition of $`b_l^{}`$ some of the wavefunctions $`\varphi _{l,k}`$ are tainted by powers of $`\frac{1}{x}`$ and as such might not be admissible as physical solutions for a dimension under consideration. Terms such as
$$\frac{1}{x}e^{x^2/2},$$
when considered in one dimension, imply normalization integrals such as
$$_{\mathrm{}}^+\mathrm{}\frac{1}{x^2}e^{x^2}𝑑x$$
which are not integrable. On the other hand, if variable $`x`$ is treated as three dimensional radius then the normalization integral like this
$$_0^+\mathrm{}\frac{1}{x^2}e^{x^2}4\pi x^2𝑑x,$$
is integrable.
### 5.1 Admissible physical solutions
A general pattern of physically acceptable solutions is as follows:
* All odd eigenfunctions $`\varphi _{l,k}\text{for }k=1,3,5\mathrm{}`$ of any oscillator $`l=0,1,2\mathrm{}`$ are physically acceptable, regardless the space dimension.
* All even eigenfunctions $`\varphi _{l,k}\text{for }k=0,2,4\mathrm{}`$ of any oscillator $`l=0,1,2\mathrm{}`$ are acceptable as physical solutions only when $`N2l+1`$, where $`N`$ is the space dimension.
The above statements can be verified by Haskell function *physicalPattern*, which produces infinite list of booleans – each stating whether a solution $`k`$ for oscillator $`l`$ and space dimension $`N`$ is admissible as physical solution
```
physicalPattern :: Integer -> Integer -> [Bool]
physicalPattern = map (isPhysical n) (wavefunctions l)
isPhysical :: Integer -> Laurent Integer -> Bool
isPhysical n f
| length us’ == 0 = True
| otherwise = False
where
L us us’ = f * f * dV
dV = diff $ L ((take (fromIntegral n)
$ repeat 0)++[1]) []
```
For example, even solutions for oscillator $`l=2`$ are unphysical in three dimensional space, $`N=3`$, as seen below (Haskell lists are traditionally indexed from zero)
```
take 6 $ physicalPattern 3 2
==> [False,True,False,True,False,True].
```
### 5.2 Linear normalization
In order to normalize the eigenfunctions of one dimensional linear spiked oscillators and to test their orthogonality conditions we need algebraic formulae for the values of improper integrals
$$P_n=_{\mathrm{}}^+\mathrm{}x^ne^{x^2}𝑑x.$$
They are as follows
$`P_0`$ $`=`$ $`\sqrt{\pi }`$
$`P_1`$ $`=`$ $`0`$
$`P_n`$ $`=`$ $`{\displaystyle \frac{n1}{2}}P_{n2}.`$
Notice that integrals with odd powers of $`x^n`$ are all zero.
All solutions for the ordinary harmonic oscillator ($`l=0`$) are physically acceptable since their Laurent series do not contain coefficients of negative powers of $`x`$. What’s more, they form the orthogonal system of eigenfunctions
$$0,k|0,k^{}=\delta _{kk^{}}$$
Physically acceptable solutions for the remaining, spiked, oscillators ($`l>0`$) are those with odd values of $`k`$, $`k=1,3,5\mathrm{}`$ Each oscillator $`l`$ has its own orthogonal basis. This can be easily checked with the help of the Haskell function *bracket*, the scalar product of two (not normalized) wavefunctions – each annotated by a pair of integers $`(l,k)`$.
```
bracket :: Integer -> (Integer,Integer)
-> (Integer, Integer) -> Fraction
bracket n (l, k) (l’, k’)
= scalarProduct n f g
where
f = wavefunction l k
g = wavefunction l’ k’
```
The integer $`n`$ specifies the space dimension $`n=1,2,3..`$ when computing volume integrals for radial cases. But we can use the same function for computing linear integrals as well – by setting, somewhat artificially, this argument to zero, $`n=0`$. For example
```
decimal 8 $ bracket 0 (7,1) (7,1) ==> 14034.40729347
decimal 8 $ bracket 0 (7,1) (7,5) ==> 0
```
Notice the convertion from fractional to decimal representation, specified here with accuracy to eight decimal places. All internal computations are performed on integers of unlimited size and on fractions made of such integers. The accuracy can be extremely good, depending on the value of *eps*, set at the top of module *Spike*. It is only for presentation purposes that we convert fractions to decimal representation.
The following table summarizes the case of linear spiked oscillators
> > | $`l=0`$ | $`k=0,1,2\mathrm{}`$ | $`E_{0,k}=1,3,5\mathrm{}`$ |
> > | --- | --- | --- |
> > | $`l=1`$ | $`k=1,3,5\mathrm{}`$ | $`E_{1,k}=5,9,13\mathrm{}`$ |
> > | $`l=2`$ | $`k=1,3,5\mathrm{}`$ | $`E_{2,k}=7,11,15\mathrm{}`$ |
> > | | | |
### 5.3 Radial normalization
Scalar product of two basis vectors for radial oscillators can be defined by a generic form of improper *volume* integral
$$l,i|l,j=_0^+\mathrm{}f_{l,i}(x)f_{l,j}(x)e^{x^2}𝑑V,$$
(22)
where $`dV`$ represents volume of an elementary sphere in N-dimensional space. Using general formula for a volume of a sphere with radius $`x`$ in $`N`$ dimensions
$$V(N)=\frac{\pi ^{N/2}}{\mathrm{\Gamma }(1+N/2)}x^N,$$
where
$`\mathrm{\Gamma }(0)`$ $`=`$ $`1`$
$`\mathrm{\Gamma }(1/2)`$ $`=`$ $`\sqrt{\pi }`$
$`\mathrm{\Gamma }(1)`$ $`=`$ $`1`$
$`\mathrm{\Gamma }(N)`$ $`=`$ $`(N1)\mathrm{\Gamma }(N1),`$
the elementary volume $`dV`$, and thus volume integrals (22) can be easily computed. Specificly, the ‘volume’ of sphere with radius $`x`$ in dimensions $`N=1,2,3`$ are given by $`2x,\pi x^2,4/3\pi x^3`$, etc. To simplify Haskell computations we will ignore constant factors in definition of the N-dimensional volume and define $`dV`$ simply as $`dV=dx,xdx,x^2dx`$, etc., for $`N=1,2,3..`$, respectively.
#### 5.3.1 One dimensional
The ‘volume’ integral in one dimension is twice the value of the one dimensional integral over $`[0,+\mathrm{})`$. This is, generally, not the same as the improper integral over $`(\mathrm{},+\mathrm{})`$ – unless the function being integrated is even: $`g(x)=g(x)`$.
The recurrence formulae for integrals $`P_n=2_0^+\mathrm{}x^ne^{x^2}𝑑x`$, are given below.
$`P_0`$ $`=`$ $`\sqrt{\pi }`$
$`P_1`$ $`=`$ $`1`$
$`P_n`$ $`=`$ $`{\displaystyle \frac{n1}{2}}P_{n2}.`$
As opposed to the case of one dimensional linear normalization the integrals $`P_n`$ for odd powers of $`n`$ do not vanish.
The case $`l=0`$ admits all solutions $`k=0,1,2\mathrm{}`$, but further investigation reveals that they do not form the basis since scalar products of one odd and one even eigenfunctions is not equal to zero, as in this example,
$$0,0|0,1=0.797884.$$
However, the set of solutions for $`l=0`$ can be split into two sets: one even ($`k=0,2,4,6\mathrm{}`$) and one odd ($`k=1,3,5\mathrm{}`$) – each forming the basis
$$0,k|0,k^{}=\delta _{kk^{}}.$$
Solutions for the remaining oscillators ($`l>0`$) are restricted to odd eigenfunctions ($`k=1,3,5..`$), because even eigenfunctions are unphysical. As a result we can summarize this case as follows:
##### Even $`k`$ basis
> > | $`l=0`$ | $`k=0,2,4\mathrm{}`$ | $`E_{0,k}=1,5,9\mathrm{}`$ |
> > | --- | --- | --- |
##### Odd $`k`$ bases
> > | $`l=0`$ | $`k=1,3,5\mathrm{}`$ | $`E_{0,k}=3,7,11\mathrm{}`$ |
> > | --- | --- | --- |
> > | $`l=1`$ | $`k=1,3,5\mathrm{}`$ | $`E_{1,k}=5,9,13\mathrm{}`$ |
> > | $`l=2`$ | $`k=1,3,5\mathrm{}`$ | $`E_{2,k}=7,11,15\mathrm{}`$ |
> > | | | |
This table is particularly interesting because it summarizes what is traditionally known as *three dimensional isotropic harmonic oscillator*. When expressed in spherical coordinates and after separation of variables the radial part of its Hamiltonian becomes
$$H=\frac{\mathrm{}^2}{2m}\frac{1}{r}\frac{d^2}{dr^2}r+\frac{1}{2}m\omega ^2r^2+\frac{l(l+1)\mathrm{}^2}{2mr^2},$$
which can be further reduced to the form (1) after representing the radial solutions as
$$R_{l,k}=\frac{1}{r}\varphi _{l,k}(r).$$
Functions $`\varphi _{l,k}`$ corespond to $`f_{l,k}e^{r^2/2}`$ investigated in this paper, and the three dimensional normalization of solutions $`R_{l,k}`$ becomes one dimensional normalization of $`\varphi _{l,k}`$
$$_0^{\mathrm{}}R_{l,k}R_{l,k^{}}\mathrm{\hspace{0.17em}4}\pi r^2𝑑r=4\pi _0^{\mathrm{}}\varphi _{l,k}\varphi _{l,k^{}}𝑑r.$$
According to (15) energies are given by $`E_{l,k}=2(l+k)+1`$, in units of $`\frac{\mathrm{}\omega }{2}`$. It is customary to represent this as $`\epsilon _n=(n+3/2)\mathrm{}\omega `$, where $`n=l+k1`$ – taking one of the values $`n=0,1,2,3`$
With every $`n`$, we can associate either one or many pairs of $`(l,k)`$, which generate the states $`|l,k`$ coresponding to the same energy $`\epsilon _n`$
$`n=0`$ $``$ $`|0,1`$
$`n=1`$ $``$ $`|1,1`$
$`n=2`$ $``$ $`|0,3,|2,1`$
$`n=3`$ $``$ $`|1,3,|3,1`$
$`n=4`$ $``$ $`|0,5,|2,3,|4,1`$
$`\mathrm{}`$ $``$ $`\mathrm{}`$
The energy levels are degenerate; when $`n`$ is even then there are $`\frac{n}{2}+1`$ states
$$|0,n+1,|2,n1,|4,n3\mathrm{}$$
and when $`n`$ is odd then there are $`\frac{n+1}{2}`$ states
$$|1,n,|3,n2,|5,n4\mathrm{}$$
associated with each energy level. When $`2l+1`$ possible eigenvalues of $`L_3`$ are also considered for each $`l`$ then the total degeneracy of energy levels becomes $`\frac{1}{2}(n+1)(n+2)`$. The solutions obtained by other means ( see for example) are exactly the same as the ones presented in this paper.
#### 5.3.2 Two dimensional
According to the condition $`N2l+1`$, we can accept both even and odd solutions for oscillator $`l=0`$, while the even solutions for the remaining oscillators ($`l>0`$) must be rejected as unphysical. However, it seems that the two dimensional normalization does not lead to any well defined pattern of orthogonal bases and therefore we reject this case as ill-specified. But similar normalization in the next dimension ($`N=3`$) provides us with several usable bases, as shown below.
#### 5.3.3 Three dimensional
In accordance with the general pattern of physical admissibility of even wavefunctions, $`N2l+1`$, the first two oscillators, $`l=0`$ and $`l=1`$, admit both even and odd solutions $`k=0,1,2\mathrm{}`$. For the remaining oscillators, $`l2`$, only the odd solutions $`k=1,3,4\mathrm{}`$ are acceptable.
For each $`l`$, the following pattern of orthogonality relations emerges
$$l,k|l,k^{}=\delta _{kk^{}}\text{for }|k^{}k|=0,4,8,12\mathrm{}$$
Consequently, the eigenvectors of the first two oscillators can be partitioned into four bases - two even and two odd, while each of the remaining oscillators, $`l2`$, has two orthogonal bases. The following four tables summarize the three dimensional normalization.
##### Even $`k`$ bases
> > | $`l=0`$ | $`k=0,4,8\mathrm{}`$ | $`E_{0,k}=1,9,17\mathrm{}`$ |
> > | --- | --- | --- |
> > | $`l=1`$ | $`k=0,4,8\mathrm{}`$ | $`E_{1,k}=3,11,19\mathrm{}`$ |
> > | $`l=0`$ | $`k=2,6,10\mathrm{}`$ | $`E_{0,k}=5,13,21\mathrm{}`$ |
> > | --- | --- | --- |
> > | $`l=1`$ | $`k=2,6,10\mathrm{}`$ | $`E_{1,k}=7,15,23\mathrm{}`$ |
##### The first odd $`k`$ basis
> > | $`l=0`$ | $`k=1,5,9\mathrm{}`$ | $`E_{0,k}=3,11,19\mathrm{}`$ |
> > | --- | --- | --- |
> > | $`l=1`$ | $`k=1,5,9\mathrm{}`$ | $`E_{1,k}=5,13,21\mathrm{}`$ |
> > | $`l=2`$ | $`k=1,5,9\mathrm{}`$ | $`E_{2,k}=7,15,23\mathrm{}`$ |
> > | | | |
Its degenerated energy levels are given by
$$\epsilon _n=(n+\frac{3}{2})\mathrm{}\omega ,$$
where $`n=l+k1=0,1,2\mathrm{}`$
##### The second odd $`k`$ basis
> > | $`l=0`$ | $`k=3,7,11\mathrm{}`$ | $`E_{0,k}=7,15,23\mathrm{}`$ |
> > | --- | --- | --- |
> > | $`l=1`$ | $`k=3,7,11\mathrm{}`$ | $`E_{1,k}=9,17,25\mathrm{}`$ |
> > | $`l=2`$ | $`k=3,7,11\mathrm{}`$ | $`E_{2,k}=11,19,27\mathrm{}`$ |
> > | | | |
Its degenerated energy levels are given by
$$\epsilon _n=(n+\frac{7}{2})\mathrm{}\omega $$
where $`n=l+k3=0,1,2,3\mathrm{}`$
#### 5.3.4 N dimensional
Specific cases, examined in the last few sections, exhibit certain pattern, which can be generalized on $`N`$ dimensions and directly verified by Haskell function *bracket*.
In each $`N`$ dimensional space there exists a treshold integer number $`l^{}=(N1)/2`$, such that for all $`ll^{}`$ all wavefunctions $`k=0,1,2\mathrm{}`$ can be accepted as physical. Above this threshold the even solutions $`k=0,2,4`$ must be rejected as unphysical. However, it appears that radial normalization in any of the even dimensions $`N=2,4,6\mathrm{}`$ does not lead to properly defined basis or bases. In contrary, for each of the odd dimensions $`N=1,3,5..`$ there is a clear pattern of several staggered orthogonal bases. Half of them are even in $`k`$ and they are subjected to the limitation of physicality. The other half are orthogonal bases with odd vectors $`k`$ and they are all physical and well defined.
For $`N=1`$ the vectors belonging to the same basis are enumerated by $`\mathrm{\Delta }k=2`$. For $`N=3`$ the staggering of bases is given by $`\mathrm{\Delta }k=4`$, for $`N=5`$ it is $`\mathrm{\Delta }k=6`$, etc.
## 6 Conclusions
We have shown that spiked oscillators described by Hamiltonians (1) have simple and exact eigenfunctions, subject to further normalization restrictions, which are specific to a problem dimension. One such restriction is related to integrability of the solutions and the other to selection of orthogonal bases. The form (1) can be interpreted either as a family of one dimensional linear spiked oscillators, or families of radial spiked oscillators in odd $`N`$ dimensional spaces, where $`N=1,3,5..`$. Even dimensions $`N=2,4,6..`$ must be rejected because they do not give rise to orthogonal bases.
Each Hamiltonian,$`H_l`$, has an uncountable number of eigenvectors $`|l,k`$, where $`k=0,1,2..`$ Generally, all solutions indexed by odd $`k=1,3,5..`$ are integrable and constitute one basis or several interleaved bases. In contrary, only a limited number of oscillators $`l`$ admit even eigenfunctions $`k=0,2,4..`$ for a given odd dimension $`N=1,3,5..`$ – each forming a basis or a set of interleaved bases.
Each Hamiltonian (1) leads to an equidistant energy spectrum, isomorphic to a spectrum of ordinary quantum oscillator, but subjected to restrictions described above. Our factorization method which is akin to *SUSY* method, and which relies on two kinds of operators: *intertwing operators* $`b_l`$ and $`b_l`$, and *ladder operators* $`a_l`$ and $`a_l`$, is purely algebraic. But while the SUSY method correlates solutions of two Hamiltonians with isomorphic spectra, our method extends this approach on infinite set of Hamiltonians. The theory is augmented by very simple Haskell program, which directly implements these operators, generates the eigenfunctions, tests integrability of the solutions and verifies orthogonality conditions for wavefunctions enumerated by quantum numbers $`l`$ and $`k`$.
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# The ages of very cool hydrogen–rich white dwarfs
## 1. Introduction
White dwarfs are well studied objects and the physical processes that control their evolution are relatively well understood. In fact, most phases of white dwarf evolution can be succesfully characterized as a cooling process. That is, white dwarfs slowly radiate at the expense of the residual thermal energy of their ions. The release of thermal energy lasts for long time scales — of the order of the age of the galactic disk ($`10`$ Gyr). While their detailed energy budget is still today somehow controversial — the reason being basically the release of gravitational energy associated to phase separation upon crystallization (Mochkovitch 1983; García–Berro et al 1988a,b; Isern et al 1997) — their mechanical structures, which are largely supported by the pressure of the gas of degenerate electrons, are very well modeled except for the outer layers. These layers control the energy output and their correct modeling is necessary to properly understand the evolution of white dwarfs. This is especially true at very low luminosities when most of the white dwarf interior has already crystallized and the cooling process is controlled almost exclusively by the behavior of the very outer layers. The situation was quite unsatisfactory until very recently, when new developements in the physics of white dwarf atmospheres allowed the calculation of reliable hydrogen–dominated white dwarf atmospheres down to effective temperatures as low as 1500 K (Saumon & Jacobson 1999; Hansen 1999).
The study of very cool white dwarfs bears important consequences, since the recent results of the microlensing experiments carried out by the MACHO team (Alcock et al 1997) yield that perhaps a substantial fraction of the halo dark matter could be in the form of very cool white dwarfs. The search for this elusive white dwarfs has not been successful yet, although there are evidences that perhaps the observational counterparts of these white dwarfs could be the stellar objects recently reported in the Hubble Deep Field (Ibata et al 1999; Méndez & Minniti 1999). Most probably one of the reasons for this failure in detecting very cool white dwarfs was that their colors were expected to be redder than they really are (Hansen 1999). Moreover, the luminosity function of disk white dwarfs has been repeatedly used during the last decade to provide independent estimates of the age of galactic disk (Winget et al 1987; García–Berro et al 1988b; Hernanz et al 1994).
In view of the recent progresses in the physics of the dense hydrogen plasma we compute new cooling sequences which will be invaluable for the study of the structure and age of our Galaxy and its halo. The paper is organized as follows: in §2 we briefly describe the evolutionary code and the adopted physical inputs, in §3 we discuss in detail the cooling sequences, we compare them to other cooling sequences available so far, and we show their evolution in the color–magnitude diagram. Finally, our conclusions are drawn in §4.
## 2. The evolutionary code
The adopted cooling code is the same evolutionary code used by Salaris et al (1997) in which we have included an accurate treatment of the crystallization process of the carbon–oxygen (C/O) core, together with updated input physics suitable for computing white dwarf evolution, which will be described in the following paragraphs. Thus the cooling sequences described in the next section have been computed using a full evolutionary code instead of the approximate procedure used, for instance, in García–Berro et al (1996).
The equation of state (EOS) for the C/O binary mixture in the gaseous phase was taken from Straniero (1988), while for the liquid and solid phases the detailed EOS from Segretain et al (1994) was used. As for the pure H and He regions, we used the results of Saumon, Chabrier & Van Horn (1995), supplemented at the highest densities by an EOS for H and He using the physical prescriptions by Segretain et al (1994). Crystallization was considered to occur at $`\mathrm{\Gamma }=180`$, where $`\mathrm{\Gamma }`$ is the usual plasma ion coupling parameter. The associated release of latent heat was assumed to be equal to $`0.77k_\mathrm{B}T`$ per ion (see sections 3.1 and 3.2 below). The additional energy release due to phase separation of the C/O mixture upon crystallization was computed following closely Isern et al (1997, 2000).
Neutrino energy losses were taken from Itoh et al (1996). The conductive opacities for the liquid and solid phase of Itoh et al (1983) and Mitake, Ichimaru & Itoh (1984) were adopted; for the range of temperatures and densities not covered by the previous results, we used the conductivities by Hubbard & Lampe (1969). OPAL radiative opacities (Iglesias & Rogers 1993) with $`Z=0`$ were used for $`T6000`$ K in the He and H envelopes. In the H envelope, and for the temperatures and densities not covered by the OPAL tables, we computed Rosseland mean opacities from the monochromatic opacities of Saumon & Jacobson (1999), after adding the contribution of hydrogen lines. The surface boundary conditions needed to integrate the stellar structure ($`P`$ and $`T`$ at $`\tau =200`$, where the diffusion approximation is valid and one can safely start to integrate the full set of stellar structure equations using Rosseland mean opacities) were obtained from detailed non–grey model atmospheres: for $`T_{\mathrm{eff}}4000`$ K we used the results of Saumon & Jacobson (1999), whereas for higher temperatures, the results of Bergeron, Wesemael & Beauchamp (1995) were adopted. Bolometric luminosities and effective temperatures of the white dwarf models were transformed into $`V`$ magnitudes and colors by using the bolometric corrections and color–$`T_{\mathrm{eff}}`$ relations derived from the same model atmospheres. Superadiabatic convection in the envelope was treated according to the ML2 parametrization of the mixing length theory (see Bergeron, Wesemael & Fontaine 1992, and references therein), which is the same formalism used also in the Saumon & Jacobson (1999) and Bergeron, Wesemael & Beauchamp (1995) computations.
When computing the cooling sequences, for each white dwarf mass an initial model was converged at log$`(L/L_{})2.0`$ by considering a C/O core with the chemical composition profile taken from the evolutionary pre–white dwarf computations of Salaris et al (1997), together with “thick” H and He layers ($`M_\mathrm{H}=10^4M_{\mathrm{WD}}`$, $`M_{\mathrm{He}}=10^2M_{\mathrm{WD}}`$), and it was evolved down to luminosities $`\mathrm{log}(L/L_{})5.5`$, when all the model sequences already have large enough ages.
## 3. The cooling sequences
We have computed cooling sequences neglecting and including the release of gravitational energy associated to phase separation with the above described physical inputs for white dwarf masses $`M_{\mathrm{WD}}=0.538`$, 0.551, 0.606, 0.682, 0.768, 0.867 and 1.0 $`M_{\mathrm{}}`$, which are the core masses derived in Salaris et al (1997). The cooling sequences for a typical 0.606 $`M_{\mathrm{}}`$ white dwarf for both cases are shown in figure 1. The adopted release of latent heat for the calculations shown in figure 1 is $`k_\mathrm{B}T`$ per ion in order to facilitate the comparison with similar calculations. The solid line corresponds to the case in which the release of gravitational energy due to phase separation has been neglected, whereas the dotted line corresponds to the case in which the effects of phase separation have been fully taken into account. Also shown in figure 1 are the cooling sequences computed with $`l=0.77k_\mathrm{B}T`$ per ion for the two above mentioned cases in which phase separation has been included or neglected (long dashed line and dashed dotted line, respectively).
### 3.1. Comparison with previous evolutionary calculations
The most distinctive feature of the cooling sequences presented here is the duration of the cooling phase itself, which is larger than in other equivalent calculations (Hansen 1999, Althaus & Benvenuto 1998). For instance, for the case in which the effects of phase separation have been neglected and the adopted latent heat is $`k_\mathrm{B}T`$ per particle, at $`\mathrm{log}(L/L_{\mathrm{}})=4.5`$ the calculations for the 0.606 $`M_{\mathrm{}}`$ white dwarf reported here yield a cooling age of $`t10.7`$ Gyr, in contrast with the calculation of Hansen (1999) which, for the same core mass, the same adopted latent heat and the same adopted envelope, yields $`t8.9`$ Gyr. That is, the cooling age of a 0.606 $`M_{\mathrm{}}`$ white dwarf derived in the present work at this luminosity is $``$20% larger than the cooling age derived by Hansen (1999).
This difference can be easily accounted for from a detailed analysis of figure 2, where the core temperature–luminosity ($`T_\mathrm{c}L`$) relationship obtained in this work — solid line — and in Hansen (1999) — dotted line — are shown in the upper panel, whereas in the lower panel their relative difference is shown as a function of the core temperature. As it can be seen there, for a given $`T_\mathrm{c}`$ the luminosity is about 30% larger down to $`\mathrm{log}(L/L_{\mathrm{}})4.5`$ and, thus, the model envelopes of Hansen (1999) are systematically more transparent than our envelopes for the same $`T_\mathrm{c}`$, resulting in a more efficient cooling of the white dwarf interior. To be precise, let us quantify how this affects the cooling sequences. For luminosities smaller than $`L_0=10^2L_{\mathrm{}}`$ the contribution of thermal neutrinos and nuclear reactions are negligible and, thus, one can safely assume that the sole contribution to the cooling process is the release of binding energy. Therefore we can write:
$$L=\frac{dB}{dt}.$$
(1)
Accordingly, the difference in the cooling times between both cooling sequences can be easily estimated:
$$\mathrm{\Delta }t(L)=_{L_0}^{L(T_\mathrm{c})}\left(\frac{1}{L_{\mathrm{BH}}}\frac{1}{L_{\mathrm{TW}}}\right)\frac{dB}{dT_\mathrm{c}}𝑑T_\mathrm{c}$$
(2)
where $`L_{\mathrm{TW}}(T_\mathrm{c})`$ stands for the $`T_\mathrm{c}L`$ relationship derived in this paper and $`L_{\mathrm{BH}}(T_\mathrm{c})`$ is the one derived by Hansen (1999). We have independently computed a set of binding energies with the same equation of state described in the previous section and used equation (2) to obtain an estimate of the difference introduced by the differences in the transparency of the envelope. At $`\mathrm{log}(L/L_{\mathrm{}})=4.5`$ we got $`\mathrm{\Delta }t2.1`$ Gyr, which is in good agreement with the value derived from the evolutionary code (1.8 Gyr). Thus, this difference can be mainly ascribed to the differences in the transparency of the adopted model envelopes.
Now the question is, why is there a difference in the model envelopes? This question cannot be answered categorically. In fact, we are using the same thicknesses for the H and He layers, the same envelope EOS and the same OPAL ($`Z=0`$) opacities for $`T6000`$ K, but not the same model atmospheres for the boundary conditions. However we do have indirect ways of checking the consistency of our results.
We recomputed the cooling sequence of our 0.606 $`M_{\mathrm{}}`$ white dwarf using a grey $`T(\tau )`$ relation for deriving the boundary conditions. We found that the $`T_\mathrm{c}L`$ relationship derived in this way was coincident, as long as $`T_{\mathrm{eff}}6000`$ K, with the one obtained using the model atmospheres boundary conditions, in agreement with the findings by Hansen (1999). This allows us to directly compare our $`T_\mathrm{c}L`$ relationship (for $`T_{\mathrm{eff}}6000`$ K) with the independent results of Althaus & Benvenuto (1998) for their 0.600 $`M_{\mathrm{}}`$ white dwarf cooling track. This calculation adopts our same metallicity and thickness for the He and H layers. Moreover, Althaus & Benvenuto (1998) also used the same EOS in the envelope and the same OPAL opacities for $`T6000`$ K. The only differences of this calculation with respect to our calculation and that of Hansen (1999) are threefold. First, Althaus & Benvenuto (1998) used a grey $`T(\tau )`$ relation for deriving the boundary conditions. As we have shown, this procedure is well justified as long as $`T_{\mathrm{eff}}6000`$ K. Second, Althaus & Benvenuto (1998) adopted a slightly different C/O profile for the core. Since the adopted internal C/O stratification does not affect at all the derived $`T_\mathrm{c}L`$ relationship, the slightly different C/O profile adopted by Althaus & Benvenuto (1998) does not influence the result of the comparison. Third, they employed the full–spectrum turbulence theory of convection by Canuto, Goldman & Mazzitelli (1996) instead of the mixing length theory, for the computation of the convective superadiabatic gradient in the envelope. However, the treatment of superadiabatic convection does not affect the $`T_\mathrm{c}L`$ relationship.
In the upper panel of figure 2 we also show the core temperature–luminosity relationship derived by Althaus & Benvenuto (1998) as a dashed–dotted line, whereas in the lower panel of this figure the relative difference with respect to the present calculation is shown. Our result closely follows that of Althaus & Benvenuto (1998) in the core temperature range when $`T_{\mathrm{eff}}6000`$ K, and therefore it is quite apparent from this figure that for some reason the model envelopes of Hansen (1999) appear to be, at least in this temperature range, far more transparent than ours. It is however remarkable that the calculation reported in the present work and that of Hansen (1999) are parallel for almost the full range of luminosities studied here. Finally both our calculation and that of Hansen (1999) differ considerably at low $`T_\mathrm{c}`$ from the calculation of Althaus & Benvenuto (1998), as is expected because of the improved atmospheric treatment at very low luminosities. For the sake of completeness we also show the relationship obtained by Wood (1995) as a long dashed line (he also used a grey $`T(\tau )`$ relationship to derive the boundary conditions), although we refrain from doing a detailed comparison with our results because this cooling sequence was computed using a different EOS for the white dwarf envelope. Note, however, that at high central temperatures, where different treatments of non-ideal effects should not substantially affect the envelope EOS, the $`T_\mathrm{c}L`$ relationship of Wood (1995) is in good agreement with our results.
The ultimate reasons for the differences found in the $`T_\mathrm{c}L`$ relationship (and cooling times) with respect to Hansen (1999) results could be various. First, it may have to do, at least for $`T_{\mathrm{eff}}6000`$ K (which roughly corresponds to $`\mathrm{log}(L/L_{\mathrm{}})=3.7`$ and $`\mathrm{log}(T_\mathrm{c})=6.6`$ for our models) with the low temperature Rosseland mean opacities used for integrating the stellar structure equations, which are not the same in both calculations.
A contribution can arise also from the treatment of the boundary condition — see the discussion in §2.1 of Hansen (1998). We have compared the values of $`T`$ and $`P`$ at $`\tau `$=2/3 for $`\mathrm{log}(g)=8`$ using figure 3 of Hansen (1999), and we have found that for a given effective temperature the values of the pressure used by Hansen (1999) are systematically higher than ours when $`T_{\mathrm{eff}}`$ is lower than 5500 K. The maximum difference (in $`\mathrm{log}P`$) is $`0.62`$ at $`T_{\mathrm{eff}}=4000`$ K (which roughly correponds to $`\mathrm{log}(L/L_{\mathrm{}})=4.5`$) and steadily decreases to $`0.41`$ at $`T_{\mathrm{eff}}=3000`$ K and $`0.18`$ at $`T_{\mathrm{eff}}=2000`$ K. This means that the model atmospheres of Hansen (1998, 1999) are more transparent than ours. More transparent atmospheres means faster cooling which may explain the difference. We have thus performed a numerical experiment to test the influence of varying the boundary conditions. In the hypothesis that the mentioned differences of pressure at $`\tau `$=2/3 exist also at higher values of the optical depth, we have modified accordingly our boundary conditions and repeated the computation of the $`0.61`$ $`M_{\mathrm{}}`$ white dwarf (whose surface gravity is $`\mathrm{log}(g)=8`$). We have found that the cooling times at luminosities lower than $`\mathrm{log}(L/L_{\mathrm{}})4.4`$ are significantly shorter (up to $``$ 1.5 Gyr). However, at luminosities lower than about $`\mathrm{log}(L/L_{\mathrm{}})=4.5`$ Hansen (1999) $`T_\mathrm{c}L`$ relationship crosses ours, and at a fixed $`T_\mathrm{c}`$ the difference in luminosity is only of about 10% with respect to our results (our luminosities being higher in this case). This shows that the combination of different Rosseland mean opacities and different boundary conditions have a compensating effect, producing a $`T_\mathrm{c}L`$ relationship similar to ours for very dim white dwarfs.
Finally, we have investigated yet another possible source of the difference between the models of Hansen (1999) and our models, that is the treatment of the conductive opacity in the H and He envelope. Although we both use for the conductive opacities in the liquid phase a combination of the Itoh et al (1983) and Hubbard & Lampe (1969) results, we have found that our respective treatments are slightly different (Hansen, private communication; see also Hansen & Phinney 1998). More in detail, we use Itoh et al (1983) conductive opacities for the liquid phase whenever $`2\mathrm{\Gamma }\mathrm{\Gamma }_{\mathrm{crist}}`$ and $`0.0001r_\mathrm{s}0.5`$, where $`r_\mathrm{s}`$ is defined as in Itoh et al (1983). Additionally, for the parameter $`y`$ defined in Itoh et al (1983) we relax the constraint $`y<0.1`$ discussed by the authors, following the recommendations by Itoh (1994). At lower densities, Hubbard & Lampe (1969) tables are used. The first two conditions yield a lower limit for the density which depends on the temperature. Instead, in the models of Hansen (1999), when this limit is lower than $`\mathrm{log}(\rho )2`$, the tables of Hubbard & Lampe (1969) are used starting from $`\mathrm{log}(\rho )=2`$. As an example, for typical values of $`\mathrm{log}P`$ and $`\mathrm{log}T`$ at the base of the convective envelope for $`\mathrm{log}(L/L_{\mathrm{}})4.5`$ we are employing the results of Itoh et al (1983) whereas in Hansen (1998, 1999) the results of Hubbard & Lampe (1969) were used. However, in the region with $`\rho 10^2`$ g cm<sup>-3</sup>, where we are still using the Itoh et al (1983) opacities, the Hubbard & Lampe (1969) opacities are different by at most $`20`$%, which does not influence appreciably the evolution of the white dwarf.
### 3.2. The role of the latent heat
One of the most important sources of energy during the evolution of cooling white dwarfs is the release of latent heat upon crystallization. In fact this source of energy, together with the energy released by chemical fractionation, completely dominates the evolution at relatively low luminosities during a considerable fraction of time. The exact value of the latent heat is obtained from the thermodynamic properties of the plasma during the phase transition. In fact, both the temperature of solidification and the precise value of the latent heat are obtained from the following set of equations:
$`P_\mathrm{L}(\rho ,T)`$ $`=`$ $`P_\mathrm{S}(\rho ,T)`$
$`\mu _\mathrm{L}(\rho ,T)`$ $`=`$ $`\mu _\mathrm{S}(\rho ,T)`$
where $`T`$ is the crystallization temperature, $`P`$ is the pressure, $`\mu `$ is the chemical potential, and the subscripts S and L stand for the solid and the liquid phase, respectively. The phase transition is solved by providing the density in the liquid phase and solving equation (3) for the crystallization temperature and for the density in the solid phase. As a consequence, the solid usually has a different density and the latent heat is then the difference in entropy between both phases. It has been generally assumed that the latent heat is of the order of $`k_\mathrm{B}T`$ per particle but, as first pointed out by Lamb & Van Horn (1975), the latent heat can differ by up to 25% from this fiducial value, the exact value depending on the density and temperature of the crystallizing layer.
Since most of the computed cooling sequences adopt this fiducial value, in the previous subsection we have used $`k_\mathrm{B}T`$ per particle when comparing our models with previously published results. It is however important to realize that the latent heat depends critically on the adopted equation of state in both the liquid and solid phase. Therefore, in order to be consistent, one should compute the release of latent heat using equation (3) and the same prescriptions adopted for the thermodynamic quantities used to describe both phases. Moreover, this fiducial value was derived by Lamb & Van Horn (1975) using the best available physical inputs at that time. Since then, there have been several revisions of the equation of state of very dense plasmas, being the most significant ones those of Stringfellow, De Witt & Slattery (1990) and Iyetomi, Ogata & Ichimaru (1993). We have solved equation (3) using the chemical potential and pressures given by Stringfellow et al (1990) and Iyetomi et al (1993) and we have obtained for the range of densities and temperatures relevant for white dwarf cooling an almost constant value of $`l=0.77k_\mathrm{B}T`$ per particle in both cases, which is considerably smaller than the fiducial value adopted in most calculations. This is thus the value we adopt for the cooling sequences described below. The net result is, obviously, a decrease in the cooling times (see figure 1) which for a $`0.61M_{\mathrm{}}`$ white dwarf is almost 0.4 Gyr at $`\mathrm{log}(L/L_{\mathrm{}})=4.5`$.
### 3.3. The effects of phase separation
Now we turn our attention to the effect of phase separation upon crystallization of the C/O binary mixture. As already noted in Isern et al (2000), this effect is differential. That is, it depends crucially not only on the total amount of gravitational energy released by phase separation upon crystallization but also on the transparency of the envelope. The delay introduced by phase separation at $`\mathrm{log}(L/L_{\mathrm{}})=4.5`$ for our $`0.6`$ $`M_{\mathrm{}}`$ model is $`1`$ Gyr. Thus the differential effect of phase separation amounts to 9.6%, in good agreement with the results of Isern et al (2000).
For the sake of completeness, in figure 3 the delays introduced by phase separation upon crystallization are shown for the full range of masses. None of the delays is larger than 1.4 Gyr at $`\mathrm{log}(L/L_{\mathrm{}})=4.5`$, where the crystallization process is complete for all but the lightest white dwarfs. It should be noted as well that, for a fixed $`T_\mathrm{c}L`$ relationship, the delays depend on the profile of the C/O binary mixture, which varies considerably from one core mass to another, and which is extremely dependent on the adopted cross section for the <sup>12</sup>C$`(\alpha ,\gamma )^{16}`$O reaction and on the criterion for convective instability (Salaris et al 1997) used in the pre–white dwarf evolutionary phase. The results of Salaris et al (1997) indicate that the lowest central chemical abundance of <sup>16</sup>O occurs for a $`1.0`$ $`M_{\mathrm{}}`$ white dwarf, thus enhancing the effect of phase separation. Moreover, the size of the central oxygen–rich region reaches its maximum for this mass; this increases the size of the mixing region and, consequently also increases the effect of phase separation. Exactly the opposite occurs for a $`0.538`$ $`M_{\mathrm{}}`$ white dwarf. For intermediate core masses the situation is more complex since the two previously mentioned effects are not linear and, therefore, the analysis of their combination is more subtle, leading, in any case, to the distribution of delays shown in figure 3. However, it is important to realize that the delay is maximum for the central range of white dwarf core masses (around $`0.768`$ $`M_{\mathrm{}}`$). In table 1 a summary of the cooling ages for all the white dwarf masses when phase separation is neglected can be found<sup>1</sup><sup>1</sup>1Detailed cooling sequences are available upon request to the authors. (upper section), together with the accumulated time delays introduced by chemical fractionation at crystallization (bottom section).
Finally, in table 2 we show the total gravitational energy released by phase separation during the crystallization process and we compare it with the total release of latent heat. For the four cooling sequences corresponding to the low mass regime ($`M_{\mathrm{WD}}0.54,0.55,0.61`$ and $`0.68M_{\mathrm{}}`$), the release of gravitational energy (which depends mostly on the adopted phase diagram for the C/O binary mixture and on the pre–white dwarf chemical profile of the C/O mixture) is smaller than the total release of latent heat — but it is of the same order of magnitude — being almost identical for $`M_{\mathrm{WD}}0.77M_{\mathrm{}}`$ white dwarf. On the contrary, for the more massive white dwarfs (cooling sequences with masses $`M_{\mathrm{WD}}0.87`$ and 1.0 $`M_{\mathrm{}}`$) the total release of latent heat is smaller than the release of gravitational energy. Consequently, we stress that this additional source of energy cannot be neglected whatsoever in a realistic calculation.
At this point it should be recalled that the pre–white dwarf stratification of the C/O binary mixture plays a crucial role in determining the delays. If a high effective rate of the $`{}_{}{}^{12}\mathrm{C}(\alpha ,\gamma )^{16}\mathrm{O}`$ reaction is adopted — as was done in Salaris et al (1997) — the abundance of oxygen in the central layers is as high as 0.74 by mass for a typical 0.606 $`M_{\mathrm{}}`$ white dwarf and the degree of mixing in the liquid layers is strongly reduced by the very steep gradients of chemical composition, thus minimizing the effect of phase separation. However, the effective cross section of the <sup>12</sup>C$`(\alpha ,\gamma )^{16}`$O reaction rate is the subject of active debate and, consequently, the delays shown in figure 3 should be regarded as a conservative lower limit.
### 3.4. The colors of very cool DA white dwarfs
In figure 4 the color–magnitude diagrams for a 0.606 $`M_{\mathrm{}}`$ white dwarf are shown for four standard colors. Starting from the upper left corner and continuing clockwise the following color indices are shown: $`VI`$, $`JK`$, $`VR`$ and $`BV`$. The solid line corresponds to a 0.606 $`M_{\mathrm{}}`$ white dwarf and the dashed–dotted line corresponds to a 1.0 $`M_{\mathrm{}}`$ white dwarf. Selected evolutionary stages for $`\mathrm{log}t=7.0`$, 8.0, 9.0, 9.5, 10.0 and 10.2 are represented as dots. Both evolutionary tracks are for the case in which phase separation has been neglected. As it can be seen in this figure the first three color indices show a pronounced turn–off at low luminosities, whereas the $`BV`$ color index does not. In the infrared colors, intrinsically faint white dwarfs with hydrogen–dominated atmospheres should be bluer than previously thought, as first discussed by Hansen (1998) and Saumon & Jacoboson (1999). This behavior is due to the blocking effect in the infrared of the H<sub>2</sub> collision–induced absorption.
Finally, in figure 5 we show white dwarf isochrones for $`t=2`$, 4, 6, 8, 10 and 12 Gyr. They are particularly useful when studying the white dwarf populations in stellar clusters, for which the progenitors can be assumed to be coeval and with the same initial chemical composition. Solid lines correspond to the case in which phase separation has been neglected, whereas the dotted lines correspond to the case in which chemical fractionation upon crystallization has been fully taken into account. The calculation of isochrones involves two factors in addition to the cooling sequences. First and most important is the initial–final mass relation, second is the lifetime (up to the thermally pulsing phase) of the white dwarf progenitors. For both, we have adopted the results of Salaris et al (1997) for solar metallicity progenitors. The isochrones were thus computed self–consistently because the same evolutionary code was used for deriving all the required evolutionary data.
We would like to stress that, since the initial–final mass relationship is almost flat in the low mass regime and the mass–main sequence lifetime relationship is very steep, any attempt to derive ages of individual field white dwarfs from the position of low mass white dwarfs in the color–magnitude diagram is subject to potentially large uncertainties, since any small error in the determination of the white dwarf mass translates into a huge relative error on its total age. It should also be pointed out that there is a very common tendency to associate bright white dwarfs with young stars, which is inaccurate since they can be either bright massive white dwarfs — and, indeed, in this case their total age is small — or bright low mass white dwarfs with a low mass progenitor, which has a large main sequence lifetime, and therefore the reverse is true (Díaz–Pinto et al 1994, Isern et al 1999).
As it can be seen in figure 5, in all the color–magnitude diagrams the isochrones show a pronounced turn-off at their dimmer end which, as the age of the isochrone increases, it is located at increasingly larger magnitudes. The presence of this turn–off is due to the contribution of the most massive white dwarfs, while the upper portion of the isochrones closely resembles the cooling track of a $``$ 0.54 $`M_{\mathrm{}}`$ object. The shape of the turn-off is modulated, in the $`JK`$ colours (for ages larger than $`t=78`$ Gyr), by the intrinsic turn to the blue of the individual cooling tracks clearly seen in figure 4. For the $`VI`$ and $`VR`$ colors the turn to the blue of the individual cooling sequences begins to contribute significantly to the isochrones for ages $`t=1314`$ Gyr.
## 4. Conclusions
In this work we have computed the cooling sequences of very cool DA white dwarfs with C/O cores. These cooling sequences include the most accurate physical description of the thermodynamic quantities for both the core and the envelope and have been computed with a well tested, self–consistent evolutionary code. We have computed as well the release of latent heat using the most up to date physical inputs for the liquid–solid phase transition and we have found that for the most recent equations of state of the degenerate core the release of latent heat amounts to $`0.8k_\mathrm{B}T`$ per particle. Additionally, the release of gravitational energy associated to phase separation during crystallization has been also properly taken into account. Color indices for several bandpasses have also been computed from the most recent synthetic spectra. Our major findings can be summarized as follows. Firstly, we have found that the most recent cooling sequences of Hansen (1999) considerably underestimate the cooling age of C/O white dwarfs with hydrogen dominated atmospheres even when chemical fractionation is neglected. We have traced back the possible source of discrepancy and we have found that at a given core temperature the $`T_\mathrm{c}L`$ relationship of Hansen (1999) differs by roughly 30% with respect to other evolutionary calculations computed with the same physical inputs. This is true not only when the comparison is done with the cooling sequences reported here but also with the cooling sequences by other authors (namely Althaus & Benvenuto 1998). The ultimate reason for this discrepancy remains unidentified although we have explored several possible causes without much success. Secondly, we have also found that a conservative lower limit to the accumulated time delay introduced by the release of gravitational energy associated to phase separation is roughly 10% at $`\mathrm{log}(L/L_{\mathrm{}})=4.5`$, in good agreement with the results of Isern et al (2000).
We have transformed our cooling sequences from the $`\mathrm{log}(L/L_{\mathrm{}})`$-$`\mathrm{log}(T_{\mathrm{eff}})`$ plane into various color–magnitude diagrams, and found that intrinsically faint DA white dwarfs have a pronounced turn–off in the infrared colors, and therefore are bluer than previously thought, in good agreement with the results of Hansen (1998, 1999). Cooling isochrones taking into account the evolutionary time of the white dwarf progenitors — suitable for the analysis of white dwarfs in stellar clusters — have been produced as well. They show in all colors a turn–off at the fainter end, due to the contribution of the more massive objects; this turn–off is modulated by the intrinsic turn to the blue of the individual cooling tracks.
Finally we would like to stress the importance of having reliable models of the evolution of white dwarfs and, thus, given the substantial differences in the inferred ages found by different authors, and the complexity of the evolutionary codes needed to compute realistic cooling sequences, more independent calculations are highly desirable.
Acknowledgements This work has been supported by the DGES grants PB97–0983–C03–02 and PB97–0983–C03–03, by the NSF grant AST 97–31438 and by the CIRIT. We sincerely thank our referee, B. Hansen, for very valuable comments and criticism which have considerably improved the original manuscript. We also want to thank P. Bergeron for kindly providing us with model atmospheres. One of us, EGB, also acknowledges the support received from Sun MicroSystems under the Academic Equipment Grant AEG–7824–990325–SP.
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# An RVB phase in the triangular lattice quantum dimer model
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## Abstract
We study the quantum dimer model on the triangular lattice, which is expected to describe the singlet dynamics of frustrated Heisenberg models in phases where valence bond configurations dominate their physics. We find, in contrast to the square lattice, that there is a truly short ranged resonating valence bond (RVB) phase with no gapless collective excitations and with deconfined, gapped, spinons for a finite range of parameters. We also establish the presence of three crystalline phases in this system.
\]
The search for an RVB phase of frustrated magnets, inspired by ideas of Pauling and begun in earnest by Anderson , has been one of the recurrent themes in research on the cuprate superconductors over the past decade and somewhat more. Shortly after Anderson’s 1987 paper on the cuprates , RVB theory bifurcated—with a gapless RVB state (with valence bonds on many length scales) pursued largely by means of gauge theoretic treatments introduced by Baskaran and Anderson defining one track, while Kivelson, Rokhsar and Sethna hewed closely to the original vision and pursued the study of a short ranged RVB state. The latter proposal, of a state with exponentially decaying spin-spin correlations and no long range valence bond order, was expected to lead to a gapped collective excitation spectrum and deconfined, gapped, spinons. Unfortunately, it turned out that this was hard to arrange on the square lattice—the simplest implementation of RVB ideas, the quantum dimer model , exhibits crystalline order and confined spinons, except at a critical point . This fact is of a piece with an instability of the paramagnetic phase of the $`O(3)`$ non-linear sigma model to breaking translational symmetry due to Berry phase effects .
In this Letter, we report that the simplest quantum dimer model on the triangular lattice does possess a short ranged RVB phase with gapped collective modes, gapped deconfined spinons and spin-charge separation in its charged excitation spectrum. We establish the presence of nearby crystalline phases with confined spinons. We also suggest that a connection, made previously by P. Chandra and ourselves , between quantum dimer models and frustrated transverse field Ising models can, in this problem, be plausibly extended to track the transition out of the RVB phase. This conjecture is closely connected to the spinon deconfinement mechanism proposed by Read and Sachdev , Wen and the ideas of Senthil and Fisher .
We study the Rokhsar-Kivelson quantum dimer Hamiltonian generalized to the triangular lattice (as the contrasts are instructive, we will comment on the square lattice results along the way):
$`\widehat{H}`$ $`=`$ $`t\widehat{T}+v\widehat{V}={\displaystyle \underset{i=1}{\overset{N_p}{}}}\{t{\displaystyle \underset{\alpha =1}{\overset{3}{}}}(|\text{}\text{}|+h.c.)`$ (1)
$`+`$ $`v{\displaystyle \underset{\alpha =1}{\overset{3}{}}}(|\text{}\text{}|+|\text{}\text{}|)\}.`$ (2)
Here, the sum on $`i`$ runs over all of the $`N_p`$ plaquettes, and the sum on $`\alpha `$ over the three different orientations of the dimer plaquettes, namely rotated by 0 and $`\pm 60^o`$. We refer to the plaquettes with a parallel pair of dimers as flippable plaquettes. As a complete orthonormal basis set we use $`\left\{|c|c=1\mathrm{}N_c\right\}`$, where $`|c`$ stands for one of the $`N_c`$ possible hardcore dimer coverings of the triangular lattice. $`\widehat{V}`$ is diagonal in this basis, with $`\widehat{V}|cn_{fl}(c)|c`$ measuring the number, $`n_{fl}(c)`$, of flippable plaquettes in configuration $`c`$.
Rokhsar and Kivelson derived the square lattice version of $`\widehat{H}`$ as the leading effective Hamiltonian in the singlet manifold consisting of nearest neighbor valence bond coverings of the lattice by utilizing their overlaps as small parameters; subsequently, it was shown by Read and Sachdev that the purely kinetic energy ($`\widehat{T}`$) piece described the $`1/N`$ dynamics of the nearest-neighbor $`SU(N)`$ Heisenberg magnet in an extreme quantum limit. The former derivation is readily generalized to the triangular lattice and, crucially in this case, yields $`t>0`$ which we assume in the remainder; the latter is specific to bipartite lattices.
$`𝐓v,t`$: At high temperatures, but less than the gap to non-valence bond states, static properties are obtained by a classical sum over all dimer configurations. Most crisply, consider $`T=\mathrm{}`$ where equal time correlators are given by unweighted averages. The square lattice problem is critical in this limit, with algebraically decaying dimer-dimer correlations , whence it is not surprising that at $`T=0`$ it orders everywhere except at a point.
The important observation that motivated the present work, is that the triangular lattice problem is disordered in this limit with exponentially decaying dimer correlations. For example, the correlation function for two parallel dimers separated by distance $`x`$ along a column can be computed by standard Pfaffian/Grassman methods as,
$`n(x)n(0)`$ $`=`$ $`({\displaystyle \frac{1}{6}})^2+G_A^2(x)+G_B(x)G_B(x)`$ (3)
$`\mathrm{where}G_{\{A;B\}}(x)`$ $`=`$ $`{\displaystyle \frac{d^2k}{(2\pi )^2}e^{ik_xx}\frac{\{2i\mathrm{sin}k_x;\overline{g}(𝐤)\}}{4\mathrm{sin}^2k_x+|g(𝐤)|^2}}`$ (4)
with $`g(𝐤)=i[1e^{ik_x}e^{ik_y}e^{i(k_x+k_y)}]`$. It is easy to see that the integrands are analytic in a finite interval about real values of the momenta, whence the two Green functions, $`G_{\{A;B\}}`$, decay exponentially with $`x`$.
This feature already yields a high temperature RVB phase, and makes a liquid phase far more likely at zero temperature, to which we now turn.
Topological Sectors: An important ingredient of the analysis of the square lattice problem is the existence of two integer winding numbers, for periodic boundary conditions, that are conserved by any local dynamics in the dimer Hilbert space . In that problem it is believed that the simplest dynamics (the analog of our $`\widehat{T}`$) is ergodic within each of the $`O(L^2)`$ sectors. On the triangular lattice, the situation appears to be quite different. Repeating the analysis for the square case uncovers only four sectors corresponding to different combinations of even and odd windings which will be equivalent in the thermodynamic limit. With the minimal dynamics of $`\widehat{T}`$ we have found that large classes of configurations can be connected to each other with no further constraint. The notable exceptions are the “staggered” configuration in Fig. 2 and its symmetry related counterparts, which have no flippable plaquettes whatsoever. But even here a local four dimer rearrangement can be shown to allow a connection to other “generic” states. We conjecture therefore that a general local dimer dynamics (which may require as little as the inclusion of a four dimer move in addition to $`\widehat{T}`$) will be ergodic in each of the four topological sectors and that the dynamics in our model is nearly so with the exception of the dynamically disconnected non-flippable configurations. Note finally, that the Perron-Frobenius theorem implies that the ground state of $`\widehat{H}`$ in each sector is nodeless.
RVB Phase, $`𝐯_𝐜<𝐯𝐭:`$ As in the square lattice case, the triangular lattice dimer model has a “Rokhsar-Kivelson” (RK) point at $`v=t`$ at which the ground states are equal amplitude superpositions of all dimer coverings in a given sector. To see this, note that a lower bound for the ground state energy is obtained by considering each plaquette individually. A non-flippable plaquette is annihilated by $`\widehat{H}`$, whereas a flippable plaquette has a potential energy of $`v`$ and a kinetic energy of, at best, $`t`$, which implies $`E_0\mathrm{min}\{0,N_p(vt)\}`$. The equal amplitude state $`|RK=N_c^{1/2}_{c=1}^{N_c}|c`$ (in any sector) has energy $`RK\left|H\right|RK=(vt)n_{fl}`$ which vanishes and saturates the lower bound at $`v=t`$. Following our statements on the sectoral organization, we conclude that with the exception of the non-flippable configurations which trivially saturate this bound, there are four topologically degenerate ground states in the thermodynamic limit. The sum over all ground states is, for the purposes of computing correlations diagonal in the dimer basis, equivalent to the classical dimer problem. As the staggered configurations are irrelevant to this sum, we conclude that the four generic sector states have exponentially decaying dimer correlations—i.e. they are RVB states.
We next wish to argue that these states are representative of a phase at $`T=0`$. We will argue shortly that their coexistence with the staggered states is due to a first order transition out of the RVB phase exactly at the RK point—much as in the square lattice problem. Consequently, we wish to establish that there is a range $`v_c<vt`$ over which the RVB character of the ground state persists. Typically, one might anticipate that a disordered ground state goes along with a gap to local excitations (as opposed to the degeneracy with globally distinct states). We have examined candidates for collective modes in the single mode approximation and found that they are all gapped. We note that, in contrast, Rokhsar and Kivelson found that their critical RVB state supported gapless excitations they dubbed resonons . These together—the disordered character of the ground state and the gap in the local excitation spectrum—are strong evidence that the RK point is part of an RVB phase which is displaced to its right but will persist to its left (in Fig. 1).
To test this argument, we have carried out quantum Monte Carlo simulations on systems upto $`36\times 36`$ sites at temperatures as low as $`t/10`$ using the method described in Ref. . As expected, we find that the dimer correlations are very short ranged and practically those of the classical dimer problem, conservatively, for $`2/3<v/t<1`$. We also find a very weak temperature dependence, suggestive of a gap. As $`v`$ is reduced further, the correlations begin to exhibit crystalline structure (see below). At $`v=t`$ we find a hysteretic (first order) transition to the staggered phase described next.
Staggered Phase, $`𝐯>𝐭`$: For $`v>t`$ the lower bound derived previously implies $`E_00`$. As the staggered states are zero energy eigenstates of $`\widehat{H}`$ they are the ground states in this range. This is similar to what happens in the square lattice problem, but the degeneracy is much lower in our problem ($`O(L^0)`$ vs $`O(L)`$) and the exactness of the states is here a consequence of dynamics rather than topology as mentioned earlier.
An interesting consequence of this last observation is that the excitation gap is $`O(L^0)`$ in system size for the staggered phase . At $`v/t=\mathrm{}`$, the lowest energy excitations are the four dimer loop rearrangement of Fig. 2. At finite but large $`v/t`$, these get dressed by additional plaquette flips and acquire a dispersion. This branch of solitons is expected, by continuity, to be the relevant set of low energy excitations in the staggered phase even close to the RK point.
Columnar Phase, $`𝐯𝐭:`$ To complete the phase diagram, we move leftwards from the RVB phase in Fig. 1, turning first to the extreme case $`v/t=\mathrm{}`$. In this limit, the kinetic term $`\widehat{T}`$ is disregarded, and the ground states are the maximally flippable states, i.e. those states $`|c`$ with maximal $`n_{fl}(c)`$. To identify the maximally flippable states, we note that each dimer can be part of at most two flippable pairs. Since the total number of dimers is fixed, the maximally flippable states are those in which each dimer belongs to two pairs.
All maximally flippable states, whose number is exponential in $`L`$, can be obtained by carrying out the two operations A and B on a particular maximally flippable state as depicted in Fig. 2. Pairs of operations of either type as well as global translations and rotations can be shown to be generated by local plaquette flips. This illustrates the point made earlier about winding number sectors. All maximally flippable states differing by an even number of operations are in the same sector, whereas a single A or B operation generates a state in another.
Turning to the case of large but finite $`v/t`$, we construct a perturbation theory using the small parameter $`t/v`$. Since any two maximally flippable states differ in at least $`O(L)`$ dimers, the degenerate perturbation theory is diagonal at any finite order. The energy shift of a state $`|c`$ at order $`2n`$ in perturbation theory depends on the number of flippable plaquettes, $`n_{fl}(c^{})`$, of the states $`|c^{}`$ which can be reached from $`|c`$ by at most $`n`$ plaquette flips.
The result of this perturbation theory is a striking example of the phenomenon of quantum “order by disorder”— we find that it selects the columnar state depicted in Fig. 2. States obtained by operations of type A are disfavored at $`4^{\mathrm{th}}`$ order in perturbation theory, those generated by type B operations at $`6^{\mathrm{th}}`$ order.
Transverse Field Ising Point, $`𝐯=\mathrm{𝟎}:`$ Together with P. Chandra , we have recently shown that there exists an exact correspondence between the quantum dimer model at $`v=0`$, and the fully frustrated transverse field Ising model (FFTFIM) on the dual hexagonal lattice at fields $`\mathrm{\Gamma }`$ much smaller than the magnitude of the exchange $`J`$ . In this limit, the quantum ground state is constructed entirely out of the ground states of the classical frustrated model and the latter are (upto Ising degeneracy), in unique correspondence with dimer coverings of its dual, triangular, lattice. In our analysis we found a low temperature crystalline “$`\sqrt{12}\times \sqrt{12}`$” phase which exhibits, in dimer language, a triangular superlattice with a 12 site unit cell, that is consistent with a Landau-Ginzburg analysis of the Ising model. While the ordering observed at accessible system sizes is not conclusive regarding the fate of the model at $`T=0`$, it appears that the columnar phase gives way to a different crystalline phase in the proximity of $`v=0`$.
Spinons: As noted before, the RVB phase has a gap to collective excitations. That is also true of the crystalline phases. Also of interest is the question of confinement for spinons—the gapped spin $`1/2`$ excitations produced by breaking a valence bond. In the dimer model, these are represented by monomers or holons that carry a spin and so the questions of holon and spinon confinement are identical. This question is most easily addressed by considering the free energies of states in which two monomers are held a fixed distance apart. At high temperatures, this is again a classical computation and it is clear that the spinons are deconfined on the triangular lattice . The contrast with the square lattice is again instructive, for there the spinons are confined at high temperatures. (However the confinement is very weak, only logarithmic , which explains why the more disordered triangular lattice does not confine.) At $`T=0`$ one can readily show that the state with an equal amplitude sum over dimer configurations with two spinons localized a fixed distance apart is an eigenstate at the RK point with an energy independent of their separation—i.e. the spinons do not interact beyond one lattice constant . By continuity, we expect spinons to be deconfined in the entire RVB phase at $`T=0`$. From these considerations it also follows that the charged excitation created by removing an electron from the system, will decay into a spinon and holon. Evidently, the holons/spinons will be confined in the crystalline phases.
Spinon Confinement Transition: Several authors have suggested that a spinon confinement-deconfinement transition will be governed by an Ising ($`Z_2`$) gauge theory . In our own work we have found that frustrated transverse field Ising models (whose duals are precisely the Ising gauge theories) can provide a description of valence bond phases of Heisenberg antiferromagnets on their dual lattices via their connection to quantum dimer models, such as the one considered in this paper. In the current context, the geometrical identification used by us leads to an intriguing observation. As already noted, the quantum dimer model is equivalent to the $`\mathrm{\Gamma }J`$ limit of the FFTFIM at $`v=0`$. Interestingly, the paramagnetic ground state of the FFTFIM at $`\mathrm{\Gamma }J`$, when projected onto the dimer manifold, is the equal amplitude RVB sum that is the dimer model ground state at $`v=t`$. This suggests the conjecture that passing between these limits in the (projected) FFTFIM gives a description of the spinon confinement and translational symmetry breaking transition that takes place at the boundary of the RVB and $`\sqrt{12}\times \sqrt{12}`$ phases. We have argued in Ref. , that the transition in the unprojected model is in the $`O(4)`$ universality class. It remains to be seen whether this identification survives projection. (The remaining transition, between the columnar and $`\sqrt{12}\times \sqrt{12}`$ states must be first order on symmetry grounds.)
In closing, we note that there are two previous “sightings” of a spin liquid phase on the triangular lattice in the literature. First, there is the large $`N`$, $`Sp(N)`$ analysis of Sachdev which found a disordered ground state with gapped unconfined spinons at small “spin”. Second, there is the exact diagonalization work of Misguich et al. on an $`S=1/2`$ system with a ferromagnetic two spin exchange frustrated by an antiferromagnetic four spin exchange. They found a four-fold disordered ground state with a spin gap, but possibly confined spinons. We hope to clarify the connection between our results and these in the near future. Finally, we note that along the lines of the analysis in we expect doping to give rise, via holon condensation, to a superconducting phase on the triangular lattice.
As we were finishing this paper, there appeared Ref. , which reports neutron scattering evidence for deconfined spinons on an anisotropic triangular lattice. We note that the classical dimer problem in that case is also disordered and that one can construct an RK point with anisotropic terms that exhibits a disordered ground state. This suggests that the results of could be understood along the lines of our analysis in this paper.
Acknowledgements: We are grateful to S. Kivelson and D. Huse for very valuable discussions and to P. Chandra for collaboration on related work. This work was supported in part by grants from the Deutsche Forschungsgemeinschaft, the NSF (grant No. DMR-9978074), the A. P. Sloan Foundation and the David and Lucille Packard Foundation.
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# Interstellar Seeing. II. The Case of the Vela Pulsar: Source Unresolved
## 1 Introduction
Diffractive interstellar scintillation (DISS) is caused by multipath scattering of radio waves from small-scale irregularities in the ionized interstellar medium. In a previous paper (Cordes 2000; hereafter paper I) methodologies have been presented that exploit DISS as a superresolution phenomenon: one that can constrain or determine source sizes on scales that are orders of magnitude smaller than the inherent diffraction resolution of the largest apertures available, those involving space VLBI. In this paper we use some of the methods of Paper I to assess recent results reported on the Vela pulsar (Gwinn et al. 1997, 2000; hereafter G97 and G00, respectively), including the conclusion that the magnetosphere of that pulsar has been resolved.
G97 analyzed VLBI observations of the Vela pulsar using time and frequency resolutions that exploit the spatial resolving power of DISS. Through estimation of the probability density function (PDF) of the visibility function magnitude, they infer that DISS shows less modulation than expected from a point source and therefore conclude that the source must be extended. They estimate a transverse size $`500`$ km for the region responsible for the pulsed flux in a narrow range of pulse phase.
We reconsider the observations of G97 because effects other than source size can diminish the modulation index (the rms fractional modulation) from the unity value expected for a point source in the strong scattering regime (see Rickett 1990 and Cordes & Lazio 1991 for definitions of this regime). In particular, even modest averaging in time and in frequency over a narrow bandwidth — a feature of essentially all observations of DISS — can reduce the modulation index to below unity. We show explicitly that the averaging parameters used by G97 can account for the entire reduction of the modulation index. To interpret their data, we adopt a Bayesian method to place limits on the source size.
G00 presented another analysis of data on Vela that took into account some aspects of time-frequency averaging and other effects that can alter the visibility PDF. They revise downward the estimate of source size to $`440`$ km. Our treatment in this paper contests this result also. However, the biggest uncertainty in establishing a detection of or limit on source size is in the DISS parameters, the scintillation bandwidth and time scale. We discuss the role of the uncertainties in these quantities in some detail.
In §2 we summarize relevant observations of the Vela pulsar and in §3 we outline the use of DISS to resolve sources, including a discussion of isoplanatic scales and formulae for the second moment and the PDF of measured interferometer visibilities. In §4 we apply a Bayesian inference method to data on the Vela pulsar which yields an upper bound on the source size. §5 discusses the anticipated scale sizes of pulsar emission regions and interprets our upper limit on Vela pulsar. The paper is summarized in §6. Much of the paper refers to Paper I.
## 2 Observations of the Vela Pulsar
The Vela pulsar is ideal for DISS studies that aim at resolving its magnetosphere because it is very bright and, as shown below, the DISS isoplanatic scale is smaller than the magnetosphere of the pulsar. G97 report visibility measurements on the Vela pulsar made at a frequency $`2.5`$ GHz. They show that the scattering diameter of this strongly scattered pulsar is slightly elongated, with dimensions (FWHM) of $`3.3\pm 0.2\times 2.0\pm 0.1`$ mas. They also report values for the DISS time scale and bandwidth of $`\mathrm{\Delta }t_\mathrm{d}15`$ s and $`\mathrm{\Delta }\nu _\mathrm{d}39\pm 7`$ kHz, respectively. By comparing the angular diameter with that predicted from the DISS parameters for specific geometries, they conclude that scattering occurs in a thin scattering screen approximately 27% of the total pulsar-earth distance from the pulsar, i.e. $`D_s/d=0.27`$, where $`D_s`$ is the pulsar-screen distance and $`d`$ is the pulsar-Earth distance. The angular size, DISS bandwidth, and the inferred ratio $`D_s/d`$ differ significantly from those reported by Desai et al. (1992) ($`1.6\pm 0.2`$ mas, $`68\pm 5`$ kHz, and 0.81, respectively). The DISS bandwidth measured by G97 agrees, however, with that reported by Cordes (1986) when scaled to 2.5 GHz using a $`\nu ^4`$ scaling law. More recently, G00 quote different values for the DISS parameters at $`2.5`$ GHz, $`\mathrm{\Delta }t_\mathrm{d}26\pm 1`$s and $`\mathrm{\Delta }\nu _\mathrm{d}66\pm 1`$ kHz. The exact values of the DISS parameters have a strong influence on the visibility fluctuations and on any inference from those fluctuations on source size.
Some DISS observations on the Vela pulsar (Cordes, Weisberg & Boriakoff 1985; Johnston et al. 1998; Backer 1974) have suggested that the DISS bandwidth scales very nearly as $`\nu ^4`$, somewhat different from the $`\nu ^{22/5}`$ scaling commonly associated with the Kolmogorov wavenumber spectrum (Rickett 1990). However, a composite of all available data from the literature, shown in Figure 1, suggests otherwise. Simple, unweighted least-squares fits of log-log quantities yields $`\mathrm{\Delta }\nu _\mathrm{d}\nu ^{4.30\pm 0.12}`$, closer to the Kolmogorov scaling than to the $`\nu ^4`$ scaling. The increase in value of the exponent is due largely to the new high frequency measurements by Johnston et al. 1998. It is conceivable that there is not a single exponent value for the entire range of frequencies plotted, but the errors do not allow any break point to be established. Similarly, the DISS time scale varies as $`\mathrm{\Delta }t_\mathrm{d}\nu ^{1.17\pm 0.06}`$, very close to the $`\nu ^{6/5}`$ Kolmogorov scaling. Residuals from the fits imply that there are typically 9% errors on $`\mathrm{\Delta }t_\mathrm{d}`$ and 22% errors on $`\mathrm{\Delta }\nu _\mathrm{d}`$. We use these errors in our inference on source size rather than using the errors quoted by G00. We think the latter are underestimated.
As shown in Paper I and also as calculated by G97, the isoplanatic scale at the pulsar’s location can be calculated from the measured DISS parameters and scattering diameter. The length scale $`\mathrm{}_d`$ in the diffraction pattern in the observer’s plane is related to the measured angular diameter $`\theta _{\mathrm{FWHM}}`$ (assuming a circular, Gaussian brightness distribution and a thin screen) and to the DISS bandwidth as (c.f. Cordes & Rickett 1998, §3.)
$`\mathrm{}_d={\displaystyle \frac{\lambda \sqrt{2\mathrm{ln}2}}{\pi \theta _{\mathrm{FWHM}}}}=\left[\left({\displaystyle \frac{cd\mathrm{\Delta }\nu _\mathrm{d}}{2\pi \nu ^2C_1}}\right)\left({\displaystyle \frac{d}{D_s}}1\right)\right]^{1/2},`$ (1)
where $`C_1=0.957`$ for a thin screen with a Kolmogorov spectrum while $`C_1=1`$ for a thin screen with a square-law structure function. Using a nominal value for $`\theta _{\mathrm{FWHM}}`$, 2.6 mas, taken to be the geometric mean of the angular diameters measured by Gwinn et al., we obtain $`\mathrm{}_d10^{3.6}`$ km. Solving for $`D_s/d`$ using nominal values for $`\theta _{\mathrm{FWHM}}`$, the DISS bandwidth (from G97), and the distance $`d=0.5`$ kpc, we obtain $`D_s/d0.27`$, in agreement with G97. We note, however, that the true distance to the pulsar is not well known. Recent work (Cha, Sembach & Danks 1999) suggests that $`d`$ may be as small as 0.25 kpc, based on association of the pulsar with the Vela supernova remnant. This would yield $`D_s/d0.15`$. If we use the DISS bandwidth from G00, however, we get $`D_s/d0.39`$ for $`d=0.5`$ kpc and $`D_s/d0.24`$ using the smaller distance.
In Paper I, we define the isoplanatic scale at the pulsar, $`\delta r_{s,iso}`$, as the separation that two sources would have if their scintillations are correlated at a level of $`e^1`$. This scale is related to the diffraction scale $`\mathrm{}_d`$ as
$`\delta r_{s,iso}={\displaystyle \frac{(D_s/d)\mathrm{}_d}{1D_s/d}}.`$ (2)
Note that the expression is not valid for $`D_sd`$ because the small-angle approximation employed in all of our analysis breaks down. Using DISS parameters from G97, the isoplanatic scale $`\delta r_{s,iso}10^{3.1}`$ km for $`d=0.5`$ kpc and $`\delta r_{s,iso}10^{2.8}`$ km for half the distance, $`d=0.25`$ kpc. Using G00 parameters, we get $`\delta r_{s,iso}=10^{3.4}`$ and $`10^{3.1}`$ km for the two distances. For comparison, the light cylinder radius is $`r_{\mathrm{LC}}=cP/2\pi 10^{3.6}`$ km (where the spin period is $`P=0.089`$ sec), which is larger than the isoplanatic scale by a factor of 3 to 7. Figure 2 shows schematically the pulsar geometry and its relationship to the isoplanatic scale. We discuss the geometry further in §5 below. Later in the paper, we place an upper limit on the source size by choosing the largest of all estimates for the isoplanatic scale, $`\delta r_{s,iso}=10^{3.4}`$ km. This gives the least restrictive upper bound.
G97 present a histogram of visibility magnitudes on a baseline that does not resolve the scattering disk. The visibilities were calculated using time-bandwidth averaging intervals of $`T=10`$ s and $`B=25`$ kHz that resolve the DISS. They integrated over a pulse phase window of 1.16 ms, corresponding to $`\mathrm{\Delta }t/P=0.013`$ cycles. They also summed over $`N_p=112`$ pulse periods to yield their quoted integration time of 10 sec.
The reported histogram departs from the shape expected for a scintillating point source, which Gwinn et al. interpret as a signature for source extension, with source size $`460\pm 110`$ km. This result is based on the larger distance estimate for the pulsar (0.5 kpc). The quoted error is said to arise from uncertainty in the scintillation bandwidth, $`\mathrm{\Delta }\nu _\mathrm{d}`$.
G97 did not consider the effects of the time-bandwidth averaging on their results. Though $`T`$ and $`B`$ are respectively smaller than the nominal DISS time scale and bandwidth, $`\mathrm{\Delta }t_\mathrm{d}`$ and $`\mathrm{\Delta }\nu _\mathrm{d}`$, they are large enough to be important in any consideration of visibility fluctuations. In fact, we show that TB averaging can account for all the departure of the fluctuations from those expected for a point source. In our analysis below, we also consider contributions to visibility fluctuations from intrinsic variations in the pulsar. We discuss these in some detail in Paper I (Appendix A), where we present the scintillated, amplitude-modulated noise model. Amplitude fluctuations include the well-known pulse-to-pulse variations that typically represent 100% modulations of the pulsed flux at a fixed pulse phase.
G00 presented a new analysis of VLBI data on the Vela pulsar and considered TB averaging on their results. Also, their re-estimated scintillation parameters reduce the role of TB averaging from what it would have been using the G97 estimates so that their inferred source size is much the same as in G97. We argue that this new analysis is incorrect, though the DISS parameters are closer to those extrapolated from other measurements.
## 3 Resolving Sources with DISS
In Paper I we discuss quantitatively how measurements of interferometric visibilities (or single-aperture intensities) contain information about the intrinsic source size, even for sources unresolved by the interferometer baseline. The issue is whether the source is extended enough to reduce the scintillation modulation. Our discussion is based on the strong-scattering regime, which easily applies to observations of the Vela pulsar. A number of methods may be used to extract source-size information, including an investigation of the modulation index of the visibility and, more accurately, the full probability density function (PDF) of the visibility magnitude. For pulsars with multiple pulse components (and presumably multiple emission regions), cross-correlation analyses may also be conducted.
### 3.1 Second Moments & Isoplanatic Scales
To evaluate visibility statistics, we need the autocovariance (ACV) function $`\gamma _\mathrm{G}`$ of the ‘gain’ $`G`$ by which the intensity or visibility is modulated by DISS. The gain has unit mean, $`G=1`$. In the strong scattering (Rayleigh) limit the ACV for a spatial offset $`\delta 𝐫_𝐬`$ at the source, a temporal offset $`\tau `$ and an interferometer baseline $`b`$ is
$`\gamma _\mathrm{G}(𝐛,\tau ,\delta \nu ,\delta 𝐫_𝐬)=G(𝐫,t,\nu ,𝐫_𝐬)G(𝐫+𝐛,t+\tau ,\nu +\delta \nu ,𝐫_𝐬+\delta 𝐫_𝐬)1=\left|\gamma _\mathrm{g}(𝐫,t,\nu ,𝐫_𝐬)\right|^2.`$ (3)
The rightmost equality uses the second moment,
$`\gamma _\mathrm{g}(𝐛,\tau ,\delta \nu ,\delta 𝐫_𝐬)=g(𝐫,t,\nu ,𝐫_𝐬)g(𝐫+𝐛,t+\tau ,\nu +\delta \nu ,𝐫_𝐬+\delta 𝐫_𝐬),`$ (4)
where $`g`$ is the wave propagator defined in Paper I. For zero frequency lag $`\gamma _\mathrm{g}`$ is related to the phase structure function $`D_\varphi `$ by the well-known relation (e.g. Rickett 1990),
$`\gamma _\mathrm{g}(𝐛,\tau ,\delta \nu =0,\delta 𝐫_𝐬)=e^{D_\varphi (𝐛,\tau ,\delta 𝐫_𝐬)}.`$ (5)
In Paper I we give a general expression for $`D_\varphi `$ that applies to any distribution of scattering material along the line of sight. For the Vela pulsar, scattering evidently is dominated by a thin screen along the line of sight (Desai et al. 1992). For a thin screen at distance $`D_s`$ from the pulsar, we have (Paper I and references therein)
$`D_\varphi (𝐛,\tau ,\delta 𝐫_𝐬)`$ $``$ $`\left({\displaystyle \frac{|𝐛_{\mathrm{eff}}(s)|}{b_e}}\right)^\alpha `$ (6)
$`𝐛_{\mathrm{eff}}(𝐛,\tau ,\delta 𝐫_𝐬)`$ $`=`$ $`(D_s/d)𝐛+𝐕_{\mathrm{eff}}\tau +(1D_s/d)\delta 𝐫_𝐬`$ (7)
$`𝐕_{\mathrm{eff}}`$ $`=`$ $`(D_s/d)𝐕_{\mathrm{obs}}+(1D_s/d)𝐕_\mathrm{p}𝐕_\mathrm{m}(D_s),`$ (8)
where $`b_e`$ is the characteristic length scale at which $`D_\varphi =1`$ rad<sup>2</sup>. A square-law structure function with $`\alpha =2`$ may apply to the Vela pulsar’s line of sight, as discussed in §2, though the bulk of the evidence suggests that $`\alpha `$ is close to the Kolmogorov value of $`5/3`$. In Eq. 8 $`𝐕_\mathrm{p}`$ is the pulsar velocity, $`𝐕_{\mathrm{obs}}`$ is the observer’s velocity and $`𝐕_\mathrm{m}`$ is the velocity of the scattering material in the ISM. Note that the offset between sources, $`\delta 𝐫_𝐬`$, may be smaller or larger than the spatial offset associated with the pulsar velocity, $`𝐕_\mathrm{p}\tau `$.
Isoplanatic scales are defined using $`\gamma _\mathrm{G}=e^1`$, implying $`D_\varphi =1`$ and $`|𝐛_{\mathrm{eff}}|=b_e`$. The length scale for the diffraction pattern (at the observer’s location) is given by $`D_\varphi (\mathrm{}_d,0,0)=1`$ while the isoplanatic scale at the source is given by $`D_\varphi (0,0,\delta r_{s,iso})=1`$. Solving for $`\mathrm{}_d`$ and $`\delta r_{s,iso}`$ and eliminating $`b_e`$ yields Eq. 2.
### 3.2 Modulation Index
The modulation index of the visibility magnitude, defined as the rms value divided by the mean intensity, has four terms (c.f. Paper I): three associated with the scintillating source intensity and a fourth due to noise, from sky backgrounds and receiver noise, that adds to the wavefield. Here we are concerned with the leading term that is caused by DISS, $`m_{\mathrm{ISS}}^2`$. For sources with brightness distribution $`I_s(𝐫_𝐬)`$, $`m_{\mathrm{ISS}}^2(𝐛,\tau )`$ is given by
$`m_{\mathrm{ISS}}^2(𝐛,\overline{\tau })`$ $`=`$ $`I^2{\displaystyle 𝑑𝐫_{𝐬}^{}{}_{1}{}^{}𝑑𝐫_{𝐬}^{}{}_{2}{}^{}I_s(𝐫_{𝐬}^{}{}_{1}{}^{})I_s(𝐫_{𝐬}^{}{}_{2}{}^{})Q_{\mathrm{ISS}}(𝐛,\overline{\tau },𝐫_{𝐬}^{}{}_{2}{}^{}𝐫_{𝐬}^{}{}_{1}{}^{},T,B)},`$ (9)
where averaging over a time span $`T`$ and bandwidth $`B`$ is contained in
$`Q_{\mathrm{ISS}}(𝐛,\overline{\tau },\delta 𝐫_𝐬,T,B)`$ $`=`$ $`(TB)^1{\displaystyle _T^{+T}}𝑑\tau ^{}\left(1\left|{\displaystyle \frac{\tau ^{}}{T}}\right|\right){\displaystyle _B^{+B}}𝑑\delta \nu \left(1\left|{\displaystyle \frac{\delta \nu }{B}}\right|\right)`$ (10)
$`\times e^{ikd^1𝐛\delta 𝐫_𝐬}\gamma _\mathrm{G}(0,\tau ^{}+\overline{\tau },\delta \nu ,\delta 𝐫_𝐬).`$
In these expressions, $`\overline{\tau }`$ is the time lag that might be imposed in any data analysis where the visibilities from one site are lagged with respect to another. It is easy to show that $`Q_{\mathrm{ISS}}1`$, with equality only when $`T0`$, $`B0`$, and $`\delta 𝐫_𝐬=0`$.
Inspection of Eq. 10 shows that visibility fluctuations are independent of baseline $`𝐛`$ for unresolved sources, for which the complex exponential $`1`$. The baseline-independent property for visibility fluctuations is similar to the conclusion found by Goodman & Narayan (1989). In general, the integrand factor $`\gamma _\mathrm{G}`$ in Eq. 9-10 is not factorable. Practical cases, such as media with square-law phase structure functions and Kolmogorov media (Paper I), do not show factorability of $`\gamma _\mathrm{G}`$. Therefore, the effects of T-B averaging and source size must be considered simultaneously.
Additional contributions to the total modulation index from intrinsic source fluctuations and from additive noise are secondary to our discussion here. They are significant, however, in any practical application where the time-bandwidth product is low and where intrinsic source fluctuations are high. For pulsars, pulse-to-pulse amplitude variations are important when only a few pulses are included in any averaging. Appendices B and C of Paper I gives full expressions for all contributions to intensity variations.
### 3.3 Number of Degrees of Freedom in Fluctuations
As Eq. 9-10 indicate, TB averaging and extended structure represent integrals over $`\gamma _\mathrm{G}`$ that diminish scintillation fluctuations. The modulation index of the averaged intensity or visibility depends on time averaging and source extension in similar ways because both increase the number of degrees of freedom in the integrated intensity. The number of degrees of freedom is
$`N_{\mathrm{dof}}=2m_{\mathrm{ISS}}^2=2N_{\mathrm{ISS}}2,`$ (11)
where $`N_{\mathrm{ISS}}`$ is the number of independent DISS fluctuations (“scintles”) that are averaged. For observations in the speckle regime, where scintles are resolved in time and frequency, we expect $`1N_{\mathrm{ISS}}2`$.
If $`m_{\mathrm{ISS}}^2=1`$ (within errors), the source is unresolved by the DISS, the baseline $`𝐛`$ has not resolved the scattering disk, and the scintillations cannot have decorrelated over the averaging intervals $`T`$ and $`B`$. The DISS gain $`G`$ then has an exponential PDF associated with the two degrees of freedom in the scattered wavefield.
Alternatively, $`m_{\mathrm{ISS}}^2<1`$ can signify (1) variation of the DISS over the averaging time or averaging bandwidth; or that (2) the source has been resolved by the DISS, i.e. that it is comparable to or larger than the isoplanatic scale of the DISS. To discriminate between these possibilities, auxiliary information is needed that characterizes the dependence of $`\gamma _\mathrm{g}`$ on its four arguments, $`𝐛,\tau ,\delta \nu `$, and $`\delta 𝐫_𝐬`$. Such information is obtained by making DISS and angular broadening measurements over a wide range of frequencies (e.g. Rickett 1990). As discussed in §2, evidence suggests that a Kolmogorov medium distributed in a thin screen is relevant.
Complications in estimating $`m_{\mathrm{ISS}}^2`$ arise from the fact that scintillating sources fluctuate, on inverse-bandwidth time scales and on a variety of longer time scales, and there is additive noise in any real-world receiver system. We consider all such complications in Paper I (main text and Appendices B and C).
### 3.4 PDF of Visibility Fluctuations
The full histogram of visibility fluctuations is potentially much more sensitive to source structure than is the second moment, as suggested by Gwinn et al. (1997, 1998). Here we summarize our derivation in Paper I of the PDF that takes into account intrinsic source fluctuations, which are important for pulsars.
We assume that the interferometer baseline $`𝐛_{ij}`$ resolves neither the source nor the scattering disk. Then the time-average visibility can be written as
$`\overline{\mathrm{\Gamma }}GI+𝒩_i\delta _{ij}+X+C,`$ (12)
where $`I`$ is the mean source intensity, $`\delta _{ij}`$ is the Kronecker delta, $`𝒩_i`$ is the mean background noise intensity, $`X`$ is a real Gaussian random variable (RV) with zero mean, and $`C`$ is a complex Gaussian RV with zero mean. Source fluctuations are described by $`X`$ which includes the noise fluctuations and amplitude fluctuations in the amplitude modulated noise model (c.f. Appendix C of Paper I). $`C`$ includes additive radiometer noise combined with source noise fluctuations, but is uninfluenced by source amplitude fluctuations. Expressions for $`\sigma _X^2`$ and $`\sigma _C^2`$ are given in Appendix C of Paper I.
The PDF for the visibility magnitude is calculated by successively integrating over the PDFs for the different, independent terms in Eq. 12, as done in Appendix C of Paper I. The PDF for the scaled visibility magnitude, $`\gamma =|\overline{\mathrm{\Gamma }}|/\sigma _C`$, is
$`f_\gamma (\gamma )`$ $`=`$ $`{\displaystyle 𝑑Gf_G(G)𝑑Xf_X(X)\left[\gamma e^{\frac{1}{2}(\gamma ^2+G^2i^2)}I_0(\gamma i)\right]_{i=(I+X/G)/\sigma _C}}`$ (13)
where $`I_0`$ is the modified Bessel function. The integrand factor in square brackets is the Rice-Nakagami PDF of a signal phasor added to complex noise (e.g. Thompson, Moran & Swenson 1991, p. 260). In the absence of any source, the PDF is simply $`f_\gamma (\gamma )=\gamma e^{\gamma ^2/2}`$.
#### 3.4.1 Pulsar Noise Contributions
In Paper I we give detailed examples of the dependence of the visibility PDF on signal to noise ratio and on the $`X`$ and $`C`$ terms. In that paper, Eq. C20 gives the PDF of the complex visibility and Eq. C25 gives the PDF of the magnitude of the visibility. Our form for the visibility PDF differs significantly from the equivalent expression in Eq. 11 of G97. The main differences are that (1) we have unequal variances of the real and imaginary parts of $`\overline{\mathrm{\Gamma }}`$; (2) there is a superfluous factor of $`2\pi `$ in Eq. 11 of G97; and (3) the definition of the variances of the real and imaginary parts in G97 is a factor of two too large.
G97 argue that a pulsar’s contributions to visibility fluctuations become unimportant for large averaging times. While the pulsar fluctuations certainly decrease with increased averaging, it is also true for the additive radiometer noise, which contributes to the $`C`$ term in Eq. 12. Thus, if one is ignored, both should be ignored. However, neither should be ignored, as the shape of the PDF depends on both. In fact, the ratio of pulsar and radiometer contributions is independent of the averaging time and depends only on the system signal to noise ratio, $`I/𝒩_i`$. In practice, it takes a very strong source for the $`X`$ term and portions of the $`C`$ term to be significant. The Vela pulsar is marginally strong enough for this to be the case. We find that the PDF shape is altered by pulsar fluctuations at the level of about 1% at a radio frequency of 2.5 GHz for the Vela pulsar (c.f. Figure 11 of Paper I).
#### 3.4.2 Sensitivity of the PDF to Time-Bandwidth Averaging
As pointed out by G97 and G00, the mode of the visibility PDF is very sensitive to the number of degrees of freedom in the DISS fluctuations. To model the shape of the PDF correctly, $`N_{\mathrm{ISS}}`$ must be correct; but the level of noise fluctuations and the source intensity must also be correct. In paper I we demonstrate how the PDF varies with the number of ISS degrees of freedom, $`2N_{\mathrm{ISS}}`$. Increases in $`N_{\mathrm{ISS}}`$ can be due to time-frequency averaging or source extent or both. Statistically, the result is the same. A pure point source can have statistics that mimic those given by an extended source if there is sufficient TB averaging. Another effect is that if the channel bandwidth $`B`$ is varied, $`N_{\mathrm{ISS}}`$ is changed but so too is the signal to noise ratio. Decreasing the bandwidth reduces $`N_{\mathrm{ISS}}`$ but increases $`\sigma _X`$ and $`\sigma _C`$ relative to the $`GI`$ term in Eq. 12, thus masking source-size effects that might also contribute to the PDF shape. In addition, the PDF shape is sensitive to the nature of the scattering medium (thin screen vs. extended medium, Kolmogorov vs. square-law structure function). For reasons discussed in §2, we adopt a thin screen, Kolmogorov scattering medium. However, we also consider a thin screen with a square-law structure function.
## 4 Bayesian Analysis of Source Size for the Vela Pulsar
We represent a set of visibility magnitudes as $`\{|\overline{\mathrm{\Gamma }}|_i,i=1,N\}`$, where the bar denotes that each measurement has been averaged explicitly over time and implicitly over frequency. Given a model for the source and parameters for the scintillations, the likelihood function for the measurements is (using $`|\overline{\mathrm{\Gamma }}|\gamma \sigma _C`$)
$`={\displaystyle \underset{i}{}}f_\gamma (\gamma _i)/\sigma _C.`$ (14)
Following a Bayesian scheme presented in Paper I, we can infer the posterior PDF for source model parameters $`𝚯`$,
$`f_𝚯(𝚯)={\displaystyle \frac{}{𝑑𝚯}},`$ (15)
where we have assumed a flat prior for the model parameters (see Paper I).
To proceed, we assume that the only source parameter is the spatial scale $`\sigma _r`$ of a circular Gaussian brightness distribution,
$`I_s(𝐫_𝐬)=I_{s}^{}{}_{0}{}^{}\left(2\pi \sigma _r^2\right)^1\mathrm{exp}\left({\displaystyle \frac{|𝐫_𝐬|^2}{2\sigma _r^2}}\right).`$ (16)
We use this distribution to calculate the modulation index (squared) using Eq. 9 for $`𝐛0,\overline{\tau }0`$ and using appropriate values for $`T,B,\mathrm{\Delta }t_\mathrm{d}`$ and $`\mathrm{\Delta }\nu _\mathrm{d}`$ (c.f. §2). This yields $`m_{\mathrm{ISS}}^2`$ as a function of $`\sigma _r`$ from which we calculate the number of ISS fluctuations, $`N_{\mathrm{ISS}}`$, using Eq. 11. For Kolmogorov media, we have used tabulated values of the covariance function from Lambert & Rickett (1999). This, in turn, we use to calculate the PDF for the DISS gain, $`G`$, that is used in Eq. 13 to calculate $`f_\gamma `$. Finally, we calculate $`f_{\overline{\mathrm{\Gamma }}}`$ by appropriate scaling of $`f_\gamma `$.
To apply our method, we use the histogram presented in Figure 2 of G97, which we write as $`N_k=Nf_{\overline{\mathrm{\Gamma }}}(\overline{\mathrm{\Gamma }}_k),k=1,N_{\mathrm{bins}}`$, where $`N`$ is the total number of visibility measurements and $`N_{\mathrm{bins}}`$ is the number of bins in the histogram. Using the histogram, we rewrite the likelihood function as a product over bins,
$`={\displaystyle \underset{k}{}}\left[f_\gamma (\gamma _k)/\sigma _C\right]^{N_k}.`$ (17)
To calculate the PDF for a given source size, we need the parameters $`\mathrm{SNR}I_{s}^{}{}_{0}{}^{}/\sqrt{2}\sigma _+`$ (the signal to noise ratio), $`\sigma _+`$ (the rms noise $`\sigma _C`$ when there is no source signal), and the number of ISS fluctuations summed in each measurement, $`N_{\mathrm{ISS}}=N_{\mathrm{ISS}}(T/\mathrm{\Delta }t_\mathrm{d},B/\mathrm{\Delta }\nu _\mathrm{d},\sigma _r/\delta r_{s,iso})`$. The off-source rms is $`\sigma _+=\frac{1}{2}(\mathrm{\Delta }tBN_p)^{1/2}𝒩`$ where $`\mathrm{\Delta }t=1.12`$ ms and $`N_p=112`$ s $`=T/P`$, and $`B=25`$ kHz. In the absence of a source, the rms visibility is $`\sqrt{2}\sigma _+`$. We have used the geometric mean of the off-pulse system noise for the two sites, $`𝒩=\left(N_iN_j\right)^{1/2}`$, expressed in flux-density units, and the source strength $`I_{s}^{}{}_{0}{}^{}`$ is the flux density in the pulse-phase window used to calculate visibilities. This is much larger than the catalogued, period-averaged flux density. For Vela, the flux density at 2.5 GHz in the pulse phase window used by G97 is about 10 Jy and the intrinsic modulation index of the pulses (from activity in the neutron star magnetosphere) is about unity (Krishamohan & Downs 1983). Evidently, the flux density varies substantially on long time scales (e.g. Sieber 1973). This is not surprising given the occurrence of refractive interstellar scintillations (e.g. Kaspi & Stinebring 1992). We therefore take $`I_{s}^{}{}_{0}{}^{}`$ to be an unknown and consider SNR to be a variable to be fitted for.
### 4.1 Analysis of G97 Data
G97 report their results in arbitrary correlation units. We use the off-pulse PDF in their Figure 2 to estimate (from the mode of the PDF), $`\sigma _+550\pm 50`$ in these units. We then search over a grid in $`N_{\mathrm{ISS}}`$ and SNR to maximize the likelihood function. The expected mean flux density in the gating window can be estimated roughly from the shape of Gwinn et al.’s histogram and from theoretical PDFs. For the range of SNR and TB averaging relevant, the mean intensity is approximately the number of correlation units where the PDF has fallen to 50% of its peak value. From Figure 2 of Gwinn et al., this is roughly 3600 to 4400 correlation units, implying SNR $`7.5\pm 0.5`$. We use these coarse estimates solely to define a search grid for $`N_{\mathrm{ISS}}`$ and SNR.
Figure 3 shows (log) likelihood contours plotted against $`N_{\mathrm{ISS}}`$ and SNR. The plus sign in the figure designates the location of maximum likelihood while the vertical lines indicate ranges of $`N_{\mathrm{ISS}}`$ expected using different estimates for the DISS parameters and using either a scattering screen with a square-law structure function or a screen with Kolmogorov electron density fluctuations. The figure caption gives details. The best-fit SNR $`7.6`$ is consistent with our crude estimate above.
The contours indicate that $`N_{\mathrm{ISS}}`$ and $`\mathrm{SNR}`$ are anticovariant: the best fit $`N_{\mathrm{ISS}}`$ increases as $`\mathrm{SNR}`$ decreases. This occurs because, as noted by G00, the peak of the visibility histogram is the feature most sensitive to model parameters. As either $`N_{\mathrm{ISS}}`$ or $`\mathrm{SNR}`$ increases the peak moves to the right. (For example, the variation with $`N_{\mathrm{ISS}}`$ is shown in Figure 12 of Paper I.) Therefore the two quantities must compensate each other in order to match the peak and are thus negatively correlated. Based on contours that we do not show, we note that smaller values of $`\sigma _+`$ move the plotted likelihood contours upward and toward the right. This also is consistent with the trend just mentioned. Therefore acceptable fits are found from a family of values for $`\mathrm{SNR}`$, $`N_{\mathrm{ISS}}`$ and $`\sigma _+`$.
Another source of uncertainty results from the assumed scattering medium. The choice of medium doesn’t alter the location of the contours<sup>1</sup><sup>1</sup>1The reason the contours do not change is because we have used the chi-square PDF with specified $`N_{\mathrm{ISS}}`$ to calculate the PDF of the scintillation gain, $`G`$. but it does alter the mapping of $`T`$, $`B`$ and source size to $`N_{\mathrm{ISS}}`$. Usage of a square-law medium rather than one with Kolmogorov statistics moves $`N_{\mathrm{ISS}}`$ to lower values.. Considering uncertainties in the type of medium relevant and in the DISS parameters, a conservative conclusion is that the maximum likelihood solution for $`N_{\mathrm{ISS}}`$ can be accounted for fully by time-bandwidth averaging without invoking any contributions from a finite source size.
The observed histogram and the PDF obtained using our best fit parameters are shown in Figure 4. The agreement of the histogram with our PDF, based on a point source and accounting for TB averaging, is as good as G97’s PDF (their Figure 2), which was based on an extended source but ignored TB averaging. This shows that not only is our best fit consistent with $`N_{\mathrm{ISS}}`$ expected for a point source (and TB averaging), but also that it is as good a fit as can be expected.
We conclude that the measurements of G97 imply only an upper bound on the source size of the Vela pulsar using the DISS method. To constrain the allowed source size, we recalculate the likelihood by varying $`\sigma _r/\delta r_{s,iso}`$, calculating $`m_{\mathrm{ISS}}^2`$ for a thin-screen, Kolmogorov medium while accounting for TB averaging and, hence, calculating $`N_{\mathrm{ISS}}`$ as a function of $`\sigma _r/\delta r_{s,iso}`$. We then use $`N_{\mathrm{ISS}}`$ to calculate, in turn, $`f_G(G)`$, the visibility PDF, the likelihood, and the posterior PDF for $`\sigma _r/\delta r_{s,iso}`$. These calculations were made by integrating over a grid of values for $`\mathrm{\Delta }t_\mathrm{d},\mathrm{\Delta }\nu _\mathrm{d},\mathrm{SNR}`$, and $`\sigma _+`$ to take into account their uncertainties. The resultant posterior PDF and CDF for the source size are shown in Figure 5 as dashed lines. The solid lines are the PDF and CDF obtained when we fix $`\mathrm{SNR}`$ and $`\sigma _+`$ at their best-fit values. We use the DISS parameter values from G00 rather than from G97 because we then obtain a larger upper bound on source size. That is, we derive the least restrictive upper bound. If we use G97 parameter values and a Kolmogorov screen, the visibility fluctuations are actually quite small compared to what is predicted. If we assume a square-law screen, however, then the $`N_{\mathrm{ISS}}`$ value expected from TB averaging is again consistent; we do not believe, however, that the square-law screen is consistent with DISS observations (c.f. Figure 1 and previous discussion in text). In the PDF, the most probable source size is formally $`\sigma _r/\delta r_{s,iso}0.04`$ but the PDF amplitude is nearly the same for zero source size. We therefore interpret the PDF and CDF in terms of an upper bound on source size. From the CDF we find that $`\sigma _r/\delta r_{s,iso}\mathrm{¡}\mathrm{}0.096`$ at the 68% confidence level and 0.16 at 95% confidence for the case where we marginalize over all parameters. Using the largest of the isoplanatic scales estimated in §2, we obtain an upper limit on source size at 95% confidence of 400 km if the pulsar is at $`d=0.5`$ kpc and 200 km if $`d=0.25`$ kpc. Use of alternative values of the DISS parameters yield even smaller upper bounds. These upper bounds are very conservative. If we fix SNR and $`\sigma _+`$ at their best fit values, we get an upper bound of 150 km at the 95% confidence level.
Clearly, our results imply an upper limit that is potentially much smaller than the source size estimated by G97. The difference is due to our taking into account the time-bandwidth averaging in the analysis. To show this, if we pretend that time-bandwidth averaging is negligible and repeat the calculation of the PDF for source size, we find a peak that excludes zero source size corresponding to a size determination similar to that of Gwinn et al., $`\sigma _r0.33\delta r_{s,iso}415`$ km. This inference is, of course, erroneous.
### 4.2 Analysis of G00 Data
G00 present a new analysis of data on the Vela pulsar that takes into account some aspects of time-bandwidth averaging. They first fit the visibility histogram using a finite source size and correct for T-B effects only after the fact. Moreover, the scintillation time scale and bandwidth are estimated simultaneously in the fitting process, yielding values that are significantly larger than values presented in G97. Larger DISS parameters cause T-B averaging to be less important so any apparent diminuition of the DISS modulation is attributable mostly to source size effects. For these reasons, the source sizes derived by G00 are only slightly less than those found in G97.
The analysis in G00 appears flawed for the following reasons. First of all, the model visibility PDF used in their fits excludes contributions from pulsar noise. In practice, this is only about a 1% error in the PDF for the first pulse gate considered by G00; but then, the PDF difference they identify between a point-source model and the data is only about 4% (G00, Figure 7). Secondly, they first fit the PDF to find the source size and then, post facto, correct the source size for time-bandwidth averaging effects. Time-bandwidth averaging is handled by calculating one-dimensional integrals (Equations 6,8 in G00) and combining them. This procedure is not mathematically correct because, as noted in §3.2, the scintillation autocovariance $`\gamma _\mathrm{G}`$ is not factorable.
To illustrate, consider the fact that during the 10 s averaging time, the pulsar moves approximately 1400 km given its proper motion of 140 km s<sup>-1</sup>, calculated for a distance of 0.5 kpc (Bailes et al. 1989). This is about three times the size inferred by G00 for the emission region in the first pulse gate. By calculating only a one-dimensional temporal integral corresponding to the averaging time and a separate source-size integral, the combined effects of source size and averaging are mis-estimated. This statement follows by considering the area swept out by the emission region as it is translated spatially by the pulsar’s velocity. For a circular region of radius $`\sigma _r`$, the area swept out $`\pi \sigma _r^2(1+2V_{p}^{}{}_{}{}^{}/\pi \sigma _r)`$, or about three times the area of the emission region itself. As shown in Paper I (c.f. Figures 1-6), the two-dimensional source-size integral causes the scintillation variance to decline faster than does the one-dimensional temporal integral. The product of the factors is less than the proper multiple integral, therefore overestimating $`m_{\mathrm{ISS}}^2`$ and thus underestimating $`N_{\mathrm{ISS}}`$.
Lastly, it is unclear exactly how the scintillation parameters are fitted for in the analysis of G00. Apparently they are solved for while also fitting the visibility histogram for the size and flux density of a source component. This differs from the standard procedure of determining the DISS parameters using the intensity correlation function, so the systematic errors in the procedure are not known. The DISS parameters are quoted to 1.5% and 4% precisions for the scintillation bandwidth and time scale, respectively. Given the number of visibility measurements used in the histogram (Table 2 of G00) and the implied number of independent ISS samples used (Eq. 30 of Paper I), the quoted measurement error for the scintillation bandwidth seems too small. Rather, it should be comparable to the error on the DISS time scale. Moreover, the scintillation parameters are significantly larger than presented by G97. It is therefore reasonable to suspect that there are systematic errors on the DISS parameters that are not included in the quoted uncertainties.
In Figure 3 we designate the range of $`N_{\mathrm{ISS}}`$ that is consistent with the DISS parameters quoted by G00. For nominal values (26 s and 66 kHz) and assuming a point source, we obtain $`N_{\mathrm{ISS}}1.125`$ when adopting a Kolmogorov screen. This value is consistent with the best fit $`N_{\mathrm{ISS}}`$ needed to account for the shape of the histogram. We therefore conclude, as before, that the data of G97 are consistent with a point source or, at least, a source whose size is below the level of detection in the scintillations.
## 5 Pulsar Geometries
To interpret the empirical constraints on source size that we have derived, we now turn to pulsar models and consider the physical quantities that determine source sizes.
Viable pulsar models associate most radio emission with relativistic particle flow along those magnetic field lines that extend through the velocity of light cylinder. Emission is beamed in the directions of particles’ velocity vectors, which are combinations of flow along field lines and corotation velocities. In the following, we assume that beaming is predominantly tangential to the field lines, with only modest corrections from rotational aberration. The total extent of the emission region in rotational latitude and longitude is small, so corotation aberration is roughly constant over the emission region. We also assume for now that any refraction of radiation in the magnetosphere is negligible. Considerable evidence supports the view that radio emission radii are much less than $`r_{\mathrm{LC}}`$ (Rankin 1990; Cordes 1992 and references therein).
### 5.1 Size of the Overall Emission Region
We distinguish between the extent of the overall emission region — where all pulse components originate — from the instantaneous size responsible for emission seen at a specific pulse phase. The latter is expected to be smaller than the former because of relativistic beaming and pulsar rotation.
What do we expect for the transverse extent of the overall emission regions? Consider the open field-line region of the pulsar to be filled with emitting material. The natural length scale is, of course, the light-cylinder radius, $`r_{\mathrm{LC}}`$. But detailed estimates depend on relativistic beaming, the magnetic-field topology, and the radial depth and location of emission regions. CWB83 estimate the transverse separation of emission regions responsible for different pulse components separated by pulse phase $`\mathrm{\Delta }\eta `$ to be
$`\mathrm{\Delta }r_s={\displaystyle \frac{1}{3}}\mathrm{\Delta }\eta \mathrm{sin}\alpha _\mu r_{em},`$ (18)
where $`\alpha _\mu `$ is the angle between the spin and magnetic axes and $`r_{em}`$ is the emission radius. Figure 2 shows schematically the locations of emission regions for two pulse components, along with the light cylinder and the isoplanatic patch. If we consider a small span of pulse phase $`\mathrm{\Delta }\eta `$ in a single pulse component (as we consider here), the same equation (18) applies. This estimate assumes constant-altitude, highly-beamed emission (with essentially infinite Lorentz factor) in directions tangential to dipolar magnetic field lines.
### 5.2 Instantaneous Sizes of Emission Regions
To predict the transverse extent appropriate for a single pulse component, we need alternative estimates. We again consider highly-beamed radiation along tangents to dipolar magnetic field lines, modified slightly by rotational aberration. For simplicity, we assume that the magnetic axis and the line of sight are both orthogonal to the spin axis. The pulse phase $`\eta `$ is, for given magnetic polar angle $`\theta `$ and radius $`r`$,
$`\eta \pm {\displaystyle \frac{3}{2}}\theta {\displaystyle \frac{2(r\overline{r})}{r_{\mathrm{LC}}}};`$ (19)
the second term accounts for propagation retardation and rotational aberration (Phillips 1992) which, together, introduce a time perturbation $`2c^1(r\overline{r})`$ (for a nonorthogonal rotator, the two would be replaced by $`[1+\mathrm{sin}\alpha _\mu ]`$). The dual signs for $`\theta `$ account for emission from either side of the magnetic axis, and we define a reference radius $`\overline{r}`$ that may be a weak function of frequency (Blaskiewicz, Cordes & Wasserman 1991; Phillips 1992).
It is clear from Eq. 19 that different $`(\theta ,r)`$ combinations can correspond to the same pulse phase $`\eta `$. Thus, a contrived emissivity distribution could produce a very narrow pulse even for a large radial depth. However, we think it more likely that the radial extent $`\mathrm{\Delta }r_{em}`$ and the total angular extent $`\theta _{\mathrm{max}}`$ will not be related in this way and so both will contribute to the overall width of the observed pulse. In Figure 6 we show the geometry and definitions of the radil depth, $`\mathrm{\Delta }r_{em}`$, and the angular width, $`\theta _{\mathrm{max}}/2`$. If the radial depth $`\mathrm{\Delta }r_{em}`$ is the same at all $`\theta \theta _{\mathrm{max}}`$, then
$`\mathrm{\Delta }\eta 3\theta _{\mathrm{max}}+2\mathrm{\Delta }r_{em}/r_{\mathrm{LC}}.`$ (20)
We can place limits on both $`\mathrm{\Delta }r_{em}`$ and $`\theta _{\mathrm{max}}`$ by assuming that either one can dominate the observed pulse width. For an observed pulse width $`\mathrm{\Delta }\eta /2\pi 0.05`$ cycles (typical of many pulsars, including the Vela pulsar), we have the separate limits
$`\theta _{\mathrm{max}}`$ $``$ $`{\displaystyle \frac{1}{3}}\mathrm{\Delta }\eta `$ (21)
$`\mathrm{\Delta }r_{em}`$ $``$ $`{\displaystyle \frac{1}{2}}\mathrm{\Delta }\eta r_{\mathrm{LC}}.`$ (22)
Another contribution to the observed pulse width may come from the radiation beam width of an individual particle. If we relax our assumption of highly beamed radiation, and allow the pulse width to be determined by the relativistic beam width, then Lorentz factors
$`\gamma _p{\displaystyle \frac{3}{2\mathrm{\Delta }\eta }}`$ (23)
are allowed.
We now wish to calculate the instantaneous source size $`\delta r_s`$ relevant to an observation over a vanishingly small range of pulse phase. At a fixed pulse phase, the transverse extent of the emission region is $`\delta r_s=0`$ as $`\mathrm{\Delta }r_{em}0`$ and $`\gamma _p\mathrm{}`$. Over a small range of pulse phase, the transverse extent is given by Eq. 18. For finite depth $`\mathrm{\Delta }r_{em}`$, the transverse extent (again for $`\gamma _p\mathrm{}`$) is
$`\delta r_s{\displaystyle \frac{1}{2}}\mathrm{\Delta }r_{em,}\mathrm{\Delta }r_{em}\theta ,`$ (24)
where $`\theta `$ is the mean $`\theta `$ over those $`(r,\theta )`$ in Eq. 19 that satisfy $`\eta =`$ constant. Taking $`\theta _{\mathrm{max}}`$ as an upper limit on $`\theta `$, we obtain
$`\delta r_s{\displaystyle \frac{1}{2}}\mathrm{\Delta }r_{em}\theta _{\mathrm{max}}.`$ (25)
If relativistic beaming is finite, the transverse extent due to finite Lorentz factor $`\gamma _p`$ is
$`\delta r_sr_{em}\gamma _p^1.`$ (26)
Expressing these results in terms of $`r_{\mathrm{LC}}`$ we have
$`{\displaystyle \frac{\delta r_s}{r_{\mathrm{LC}}}}\{\begin{array}{cc}\frac{1}{2}ϵ_{em}\theta _{\mathrm{max}}\left(\frac{\mathrm{\Delta }r_{em}}{r_{em}}\right)\frac{1}{12}\left(\mathrm{\Delta }\eta \right)^2\hfill & \text{finite depth}\hfill \\ & \\ ϵ_{em}\gamma _p^1\frac{2}{3}ϵ_{em}\mathrm{\Delta }\eta \hfill & \text{finite Lorentz factor},\hfill \end{array}`$ (30)
where $`ϵ_{em}r_{em}/r_{\mathrm{LC}}`$ and the limits are based on Eq. 21-22.
Pulsar phenomenology suggests that $`ϵ_{em}0.1`$ and $`\mathrm{\Delta }\eta /2\pi 0.1`$ cycles, implying nearly identical upper limits in Eq. 30 that are $`3`$% of $`r_{\mathrm{LC}}`$. However, the true limits are probably much smaller. Rankin (1990, 1993) has shown that pulse widths (in cycles) scale $`P^{1/2}`$, as expected if they are determined by the angular extent of the open field-line region near the magnetic axis. (A slightly different scaling has been described by Lyne & Manchester .) This implies that any contributions from relativistic beaming and radial depth are small. If so, the upper limits in Eq. 21-22 are probably a factor of ten smaller and the lower limit on $`\gamma _p`$ in Eq. 23 is a factor of 10 larger.
### 5.3 Length Scales in the Vela Pulsar’s Magnetosphere
Using the pulse width (4.5 ms at 10%) and period (89 ms) of the Vela pulsar, we have $`\mathrm{\Delta }\eta /2\pi 0.05`$ cycles and Eq. 30 yields upper bounds on $`\delta r_s/r_{\mathrm{LC}}`$ of 0.008 and $`0.02(ϵ_{em}/0.1)`$ for radial depth and Lorentz-factor limits, respectively. These correspond to 34 and 89 km (for $`ϵ_{em}=0.1)`$. Rankin (1990) has classified the radio pulse for the Vela pulsar as a ‘core’ component whose width is consistent with the angular size of the open field-line region expected for a dipolar field and for an emission radius $`r_{em}2R_{\mathrm{NS}}`$, where $`R_{\mathrm{NS}}=10`$ km is the assumed neutron star radius. By contrast, Krisnamohan & Downs define four separate components in the radio emission. Thus $`r_{em}20`$ km, corresponding to $`ϵ_{em}10^{2.3}`$. With $`\mathrm{\Delta }r_{em}r_{em}`$, we obtain limits $`\delta r_s/r_{\mathrm{LC}}10^{3.61}`$ and $`10^{3.0}`$, or $`\delta r_s1`$ and 4 km, respectively. The largest relevant transverse extent may in fact derive from the finite gating window used by G97 and G00, which is $`\mathrm{\Delta }\eta /2\pi =0.013`$ cycles, yielding
$`\delta r_s{\displaystyle \frac{1}{3}}\mathrm{\Delta }\eta r_{em}118ϵ_{em}\mathrm{km}.`$ (31)
Simple pulsar geometries therefore suggest that the transverse scales of pulsar emission region(s) in the Vela pulsar should be less than the upper bound from the DISS observations and, in fact, may be very much less. In addition, using Eq. 23 and $`\mathrm{\Delta }\eta /2\pi 0.01`$ cycle for the component on the leading edge of Vela’s radio pulse (Krishnamohan & Downs 1983), a lower bound on the Lorentz factor is $`\gamma _p150`$.
More complicated geometries may ensue if refraction within the pulsar magnetosphere is important, as has been suggested by a number of authors (Melrose 1979; Barnard & Arons 1986; Cordes & Wolszccan 1988; G97; G00). Refraction can duct wavemodes along magnetic field lines and then convert energy to propagating electromagnetic waves at some altitude that effectively could be defined as the ‘emission’ altitude. Any such ducting may alter the scalings of radio beam width with period from those expected from the Goldreich-Julian polar-cap size and in the absence of refraction. Differential refraction might cause radiation from disparate regions of the magnetosphere to reach the observer simultaneously. However, like the simple geometry considered here, if refraction is significant, one might also expect pulse widths to be much larger than observed, or at least larger than predicted from the size of the open-field-line region. That statement, along with the remarkable consistency of scaling laws for pulse widths with the Goldreich-Julian polar cap size and relatively small emission radii, suggests that refraction does not enlargen the transverse scales from which radio radiation emerges from pulsar magnetospheres.
We conclude that the magnetosphere of the Vela pulsar is likely to have radio emission regions with transverse scales too small to have been detected in the observations of Gwinn et al. It is also possible that they will never be detected using DISS or any other technique, for that matter, short of traveling to the vicinity of the pulsar.
### 5.4 Comparison with Gamma-ray Emission
High-energy emission from the Vela pulsar, such as the $`>100`$ MeV pulsed gamma-rays seen with EGRET, shows a double pulse that is offset to later pulse phases from the radio pulse. Rankin (1990) classifies the radio pulse as a core component which, if similar to core emission from other pulsars, is consistent with an emission altitude close to the neutron star surface. Gamma-ray emission may be high-altitude “polar-cap” radiation that derives from the flow along open field lines originating near the magnetic polar cap (e.g. Harding & Muslimov 1998). Alternatively, it may be “outer-gap” emission from near the light cylinder whose beaming is highly affected by rotation (e.g. Romani 1996). Our conclusion that the radio emission’s transverse extent is small, at least in a narrow range of pulse phase, is consistent with both of these pictures of $`\gamma `$-ray emission.
## 6 Summary and Conclusions
In this paper we applied our superresolution methodology to the recent VLBI observations of the Vela pulsar by Gwinn et al. (1997, 2000) and find that the scintillation statistics may be accounted for fully by time-bandwidth averaging and a finite signal-to-noise ratio. Any contribution from extended source structure is less than an upper limit of about 400 km at the 95% confidence interval. This limit is larger than the size expected from conventional models that place radio emission well within the light cylinder of the pulsar and therefore is not restrictive on those models. Scintillation observations at frequencies lower than 2.5 GHz have better resolution — because the isoplanatic scale at the pulsar’s position scales as $`\nu ^{1.2}`$ — and thus may be able to detect the finite source size. Recent work by Macquart et al. (2000) failed to find any effects of finite source size at 0.66 GHz.
Our results differ from those of Gwinn et al. for several reasons. We have used a more complete signal model combined with a rigorous treatment of time-bandwidth averaging and consideration of the substantial uncertainties in the measured scintillation parameters. These parameters, the diffraction time scale and bandwidth, determine the number of degrees of freedom encompassed by scintillation fluctuations when time-bandwidth averaging and source-size effects are considered. This number, in turn, has a strong influence on the shape of the visibility histogram, from which the source size is inferred or constrained. Another source of uncertainty is the nature of the scattering medium, viz. the wavenumber spectrum and phase structure function. We have found that the medium is close to being Kolmogorov in form by investigating the scaling laws of the scintillation parameters with frequency. We note also that if the shape of the visibility histogram is analyzed incorrectly, the inferred source size will have a systematic error that will scale with the isoplanatic length scale referenced to the location of the source (Eq. 2). The isoplanatic scale, $`\delta r_{s,iso}`$, scales with frequency approximately as $`\nu ^{1.2}`$, so one would expect the mis-estimated source size to scale in the same way.
Our methodology can be applied to any radio source in the strong scattering regime, including compact active galactic nuclei and gamma-ray burst afterglows, though it is not clear that these sources are compact enough to show diffractive scintillations in this regime. In another paper, we will address sources of these types and we will also consider scintillations in the weak and transition scattering regimes.
I thank Z. Arzoumanian, S. Chatterjee, C. R. Gwinn, H. Lambert, T. J. W. Lazio, M. McLaughlin, and B. J. Rickett for useful discussions and H. Lambert and B. J. Rickett for making available their numerically-derived autocovariance functions for Kolmogorov media. This research was supported by NSF grant 9819931 to Cornell University and by NAIC, which is managed by Cornell University under a cooperative agreement with the NSF.
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# 1 Introduction
## 1 Introduction
The $`SL(2,R)`$ WZNW model is an old and interesting problem. The case of WZNW models based on a compact Lie group such as $`SU(2)`$ was completely analysed in . In this case the unitary representations are finite dimensional and correlation functions can be written as a finite sum of conformal blocks by the bootstrap approach. The non-compact groups gives a much more complicated situation and the general solution is not known. There has been much work done on analysing this .
The group $`SL(2,R)`$, which is one of the simplest non-compact groups, is important in many key areas.
The first is in two dimensional gravity . The gravitational Ward identities for correlation functions of the metric are precisely the same as the $`SL(2,R)`$ KZ equations . If matter is coupled to gravity then the gravitationally dressed matter correlation functions also obey identities of the same form as $`SL(2,R)`$ KZ equations . The case of string theory on $`AdS_3`$ is described by an exact conformal field theory - the $`SL(2,R)`$ WZNW theory. There has recently been much renewed interest in this due to the AdS/CFT correspondence although so far correlation functions have only been computed in the leading supergravity approximation. The final application is with respect to the Quantum Hall plateau transition where the $`SL(2,R)`$ WZNW model has recently been proposed as a low energy effective field theory (See also where a more general $`PSL(2|2)`$ group was discussed, and earlier references .)
Correlation functions in WZNW model satisfy Knizhnik-Zamolodchikov (KZ) equation which is a set of coupled partial differential equations coming from constraints due to null states. For compact groups like $`SU(N)`$ it can be solved exactly for four point functions by the bootstrap approach which reduces it to a set of ordinary differential equations. This is due to the fact that only a finite number of conformal blocks occur. In the $`SL(2,R)`$ case one particular four point function was found that had a logarithmic behaviour. In this paper we find some more exact solutions of KZ equation which are logarithmic in nature and show that this cannot be understood if the only fields present are primary operators.
These logarithmic singularities in four point functions have a natural interpretation in terms of logarithmic conformal field theory (LCFT) . Correlation functions in LCFT have a special form. In particular operator product expansion (OPE) of two operators can receive contribution from logarithmic operators whenever the dimensions of the two or more operators become degenerate. This gives logarithms of cross ratios in four point functions. LCFTs are also discussed earlier in connection with the WZNW model . The presence of logarithmic singularities in the four point functions in $`SL(2,R)`$ was discussed in but only in some asymptotic limits.
The plan of the paper is as follows. In next section we discuss correlation functions in the $`SL(2,R)`$ WZNW model. In the following section we present some exact solutions of KZ equation and also some asymptotic solutions which are logarithmic in nature. Finally we make some comments about the OPE structure in $`SL(2,R)`$ WZNW theory and the appearance of several symmetries.
## 2 Correlation functions
The WZNW theory has affine $`SL(2,R)`$ left and right moving symmetries whose modes generate the Kac-Moody algebra:
$`[J_n^3,J_m^3]`$ $`=`$ $`{\displaystyle \frac{k}{2}}n\delta _{n+m,0}`$
$`[J_n^3,J_m^\pm ]`$ $`=`$ $`\pm J_{n+m}^\pm `$ (1)
$`[J_n^+,J_m^{}]`$ $`=`$ $`2J_{n+m}^3+kn\delta _{n+m,0}`$
We have similar expressions for the $`\overline{J}_n^a`$. The zero mode algebra is just the group $`SL(2,R)`$. We introduce the following representation for the $`SL(2,R)`$ generators :
$`J^+=x^2{\displaystyle \frac{}{x}}2jx,J^{}={\displaystyle \frac{}{x}},J^3=x{\displaystyle \frac{}{x}}+j`$ (2)
There is also a similar algebra in terms of $`\overline{x}`$ for the antiholomorphic part. From now on we drop the antiholomorphic dependence in our notation although it can easily be restored. They obey the standard relations:
$`[J^+,J^{}]=2J^3[J^3,J^\pm ]=\pm J^\pm `$ (3)
The quadratic Casimir is:
$`C_2=\eta _{ab}J^aJ^b={\displaystyle \frac{1}{2}}J^+J^{}+{\displaystyle \frac{1}{2}}J^{}J^+J^3J^3=j(j1)`$ (4)
The stress-energy tensor in the case of the ungauged WZNW model is given in terms of the normal Sugawara construction:
$$T(z)=\frac{1}{k2}\eta _{ab}:J^a(z)J^b(z):$$
(5)
where : : denotes normal ordering. In terms of a mode expansion:
$$T(z)=\underset{n}{}\frac{L_n}{z^{n+2}}L_n=\frac{1}{k2}\underset{m}{}\eta _{ab}:J_m^aJ_{nm}^b:$$
(6)
The modes $`L_n`$ obey the standard Virasoro algebra with central charge $`c=\frac{3k}{k2}`$.
We introduce primary fields, $`\varphi _j(x,z)`$ of the KM algebra:
$$J^a(z)\varphi _j(x,w)=\frac{1}{zw}J^a\varphi _j(x,w)$$
(7)
where $`J^a`$ is given by (2). The fields $`\varphi _j(x,z)`$ are also primary with respect to the Virasoro algebra with $`L_0`$ eigenvalue:
$$h=\frac{j(j1)}{k2}$$
(8)
The two point functions are fully determined using global $`SL(2,R)`$ and conformal transformations and can be normalised in the standard way:
$$\varphi _{j_1}(x_1,z_1)\varphi _{j_2}(x_2,z_2)=\delta _{j_2}^{j_1}x_{12}^{2j_1}z_{12}^{2h}$$
(9)
The general form of the three point function is:
$`\varphi _{j_1}(x_1,z_1)\varphi _{j_2}(x_2,z_2)\varphi _{j_3}(x_3,z_3)=C(j_1,j_2,j_3)x_{12}^{j_1j_2+j_3}x_{13}^{j_1j_3+j_2}x_{23}^{j_2j_3+j_1}`$ (10)
$`z_{12}^{h_1h_2+h_3}z_{13}^{h_1h_3+h_2}z_{23}^{h_2h_3+h_1}`$
The $`C(j_1,j_2,j_3)`$ are the structure constants which in principle completely determine the entire theory. For the case of $`SL(2,R)`$ these were found in . In the case of the closely related Liouville theory these were found in .
For the case of the four point correlation functions of $`SL(2,R)`$ primaries the form is determined by global conformal and $`SL(2,R)`$ transformations up to a function of the cross ratios.
$`\varphi _{j_1}(x_1,z_1)\varphi _{j_2}(x_2,z_2)\varphi _{j_3}(x_3,z_3)\varphi _{j_4}(x_4,z_4)`$ $`=`$ $`z_{43}^{h_2+h_1h_4h_3}z_{42}^{2h_2}z_{41}^{h_3+h_2h_4h_1}`$ (11)
$`z_{31}^{h_4h_1h_2h_3}x_{43}^{j_2+j_1j_4j_3}x_{42}^{2j_2}`$
$`x_{41}^{j_3+j_2j_4j_1}x_{31}^{j_4j_1j_2j_3}F(x,z)`$
Here the invariant cross ratios are:
$$x=\frac{x_{21}x_{43}}{x_{31}x_{42}}z=\frac{z_{21}z_{43}}{z_{31}z_{42}}$$
(12)
In a normal CFT we expect the OPE of two primary fields to take the general form:
$`\varphi _{j_1}(x_1,z_1)\varphi _{j_2}(x_2,z_2)={\displaystyle \underset{J}{}}C(j_1,j_2,J)z_{12}^{h_1h_2+h_J}x_{12}^{j_1j_2+J}[\varphi _J(x_2,z_2)]`$ (13)
where we have by denoted $`[\varphi _J]`$ all fields that can be produced from the given primary field $`\varphi _j`$. In principle given $`C(j_1,j_2,j_3)`$, we know the entire operator content of the theory and should be able to determine all higher point correlation functions using the OPE (13) and the crossing symmetries. This is called conformal bootstrap and has only so far been solved for the minimal models. We will see however that our solutions require more operators to be included in the OPE.
## 3 The Knizhnik-Zamolodchikov Equation
Correlation functions of the WZNW model satisfy a set of partial differential equations known as Knizhnik-Zamolodchikov (KZ) equation due to constraints from the null states following from (6). These are:
$$|\chi =(L_1\frac{1}{k2}\eta _{ab}J_1^aJ_0^b)|\varphi $$
(14)
For two and three point functions this gives no new information. However for the four point function (11) it becomes a partial differential equation for $`F(x,z)`$. For a compact Lie group this equation can be solved as it reduces to a set of ordinary differential equations. For a non-compact group the situation is much more complicated. Below we specialise to the case of the group $`SL(2,R)`$ in which case the KZ equation is similar to that for $`SU(2)`$ which was considered in :
$$\left[(k2)\frac{}{z_i}+\underset{ji}{}\frac{\eta _{ab}J_i^aJ_j^b}{z_iz_j}\right]\varphi _{j_1}(z_1)\mathrm{}\varphi _{j_n}(z_n)=0$$
(15)
where $`k`$ is the level of $`SL(2,R)`$ WZNW model.
If we now use our representation (2) we find the KZ equation for four point functions is:
$$(k2)\frac{}{z}F(x,z)=\left[\frac{𝒫}{z}+\frac{𝒬}{z1}\right]F(x,z)$$
(16)
Explicitly these are:
$`𝒫`$ $`=`$ $`x^2(1x){\displaystyle \frac{^2}{x^2}}+((j_1+j_2+j_3j_4+1)x^22j_1x2j_2x(1x)){\displaystyle \frac{}{x}}`$ (17)
$`+2j_2(j_1+j_2+j_3j_4)x2j_1j_2`$
$`𝒬`$ $`=`$ $`(1x)^2x{\displaystyle \frac{^2}{x^2}}((j_1+j_2+j_3j_4+1)(1x)^22j_3(1x)2j_2x(1x)){\displaystyle \frac{}{x}}`$ (18)
$`+2j_2(j_1+j_2+j_3j_4)(1x)2j_2j_3`$
These are the same expressions obtained in except we differ in notation by $`jj`$.
## 4 Solutions of the KZ equation
### 4.1 Exact solutions
We seek solutions of the KZ equation which are of the following form:
$$F(x,z)=z^{\frac{\alpha }{k2}}(1z)^{\frac{\beta }{k2}}x^p(1x)^q[X(x)+Z(z)]$$
(19)
This ansatz allows us to see the various crossing symmetries in the solution nicely. Here we summarize the solutions leaving the details to the Appendix. Also we restrict ourselves to real values of $`j`$ (Discrete series; see Appendix for details about $`SL(2,R)`$ representation theory) in this section and in the rest of the paper.
* $`j_1=j_2=j_3=j_4=0`$ This is the solution found previously in .
$$F(x,z)=A\left[\frac{\mathrm{ln}z}{k2}\mathrm{ln}x\right]+B\left[\frac{\mathrm{ln}(1z)}{k2}\mathrm{ln}(1x)\right]+C$$
(20)
* Logarithms in $`x`$ and $`z`$.
1. $`j_1=j_2=0jj_3=j_4\frac{1}{2}`$.
$$F(x,z)=A\left[\frac{\mathrm{ln}(z1)}{k2}+\frac{\mathrm{ln}(x1)}{2j1}\right]+B$$
(21)
This obeys the crossing symmetry $`12x\frac{x}{x1},z\frac{z}{z1}`$.
2. $`j_1=j_3=0jj_2=j_4\frac{1}{2}`$.
$$F(x,z)=A\left[\frac{\mathrm{ln}z\mathrm{ln}(z1)}{k2}+\frac{\mathrm{ln}x\mathrm{ln}(x1)}{2j1}\right]+B$$
(22)
There is now no reason to expect symmetry under $`12`$ but instead we have $`13x1x,z1z`$. This solution is also well behaved as $`x,z\mathrm{}`$
3. $`j_1=j_4=0jj_2=j_3\frac{1}{2}`$.
$$F(x,z)=A\left[\frac{\mathrm{ln}z}{k2}+\frac{\mathrm{ln}x}{2j1}\right]+B$$
(23)
Here we see the symmetry $`14x\frac{1}{x},z\frac{1}{z}`$.
* $`j_2=0j_1+j_3+j_41=0`$ Logarithms in $`z`$ disappear.
$`\varphi _{1j_3j_4}(x_1,z_1)\varphi _0(x_2,z_2)\varphi _{j_3}(x_3,z_3)\varphi _{j_4}(x_4,z_4)=z_{43}^{h_1h_4h_3}z_{41}^{h_3h_4h_1}`$
$`z_{31}^{h_4h_1h_3}x_{43}^{12j_32j_4}`$
$`z_{41}^{2j_31}x_{31}^{2j_41}{\displaystyle (1x)^{2j_3}x^{2j_1}𝑑x}`$ (24)
This has logarithms if $`j_1`$ or $`j_3Z`$.
There are also solutions related to the above by various symmetries which we mention later.
### 4.2 Asymptotic solutions
In some asymptotic solutions were found by considering the $`x1`$ asymptotics. However it was not clear, without having some exact solutions, that these persisted to all orders and that the logarithmic behaviour was genuine. Here we look at the $`z0`$ leading asymptotics. If we assume a power series in $`z`$ then the leading form of $`F(x,z)`$ obeys:
$$(k2)\frac{F(x,z)}{z}=\frac{𝒫F(x,z)}{z}$$
(25)
We can easily solve this by separation of variables to get solutions:
$$F_\alpha (x,z)=z^{\frac{\alpha }{k2}}G_\alpha (x)$$
(26)
with:
$`G_\alpha (x)`$ $`=`$ $`A_2F_1(j_2j_1+{\displaystyle \frac{1}{2}}+{\displaystyle \frac{\eta }{2}},j_3j_4+{\displaystyle \frac{1}{2}}+{\displaystyle \frac{\eta }{2}};1+\eta ;x)x^{j_1j_2+\frac{1}{2}+\frac{\eta }{2}}`$ (27)
$`+B_2F_1(j_2j_1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\eta }{2}},j_3j_4+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\eta }{2}};1\eta ;x)x^{j_1j_2+\frac{1}{2}\frac{\eta }{2}}`$
$`\eta `$ $`=`$ $`\sqrt{4j_{1}^{}{}_{}{}^{2}4j_1+4j_{2}^{}{}_{}{}^{2}4j_2+14\alpha }`$
The value of $`\alpha `$ is arbitrary so far but if we wish to have agreement between the OPE and the tensor product of two representations then we should have $`\alpha =0`$ as this leads to:
$$\varphi _{j_1}(x,z)\varphi _{j_2}(0,0)=\underset{J=j_1+j_2}{}|z|^{h_1h_2+h_J}|x|^{j_1j_2+J}\left[O_J(0,0)\right]$$
(28)
where $`[O_J]`$ denotes the conformal family of operators, which are not primaries in general, that have conformal dimension $`J`$.
We thus generically have logarithms in the Kac-Moody states if $`\eta Z`$.
## 5 Structure of the OPE
### 5.1 Logarithmic Conformal Field Theory
The occurrence of logarithmic singularities in conformal blocks signals the appearance of logarithmic operators. In a conformally invariant theory the OPE of primary operators is of the form
$$O_i(z_1)O_j(z_2)=\underset{k}{}\frac{f_{ij}^k}{|z_{12}|^{h_i+h_jh_k}}O_k(z_1)+\mathrm{}$$
(29)
Two and three point functions are fixed by global conformal invariance up to a numerical constant. Four point functions are also fixed up to an arbitrary function of conformally invariant cross-ratios.
$$<O_i(z_1)O_j(z_2)O_k(z_3)O_l(z_4)>=\underset{m,n}{}\frac{f_{ij}^mf_{kl}^n}{|z_{12}|^{h_i+h_j}|z_{34}|^{h_k+h_l}}F(z)$$
(30)
where $`z`$ is the conformally invariant cross-ratio defined as above. If the OPE is of the above form then $`F(z)`$ is a power series in $`z`$ and there can be no logarithmic singularities in four point functions. However we have explicit solutions which are logarithmic in nature. In that case it becomes necessary to modify the OPE and include the logarithmic operators:
$$O_i(z_1)O_j(z_2)=\frac{1}{|z_{12}|^{h_i+h_jh}}(C\mathrm{ln}|z_{12}|^2+D)+\mathrm{}$$
(31)
where operators $`C`$ and $`D`$ form a logarithmic pair with the degenerate dimension $`h`$. We have the following form of two point functions for $`C`$ and $`D`$
$`<C(z_1)C(z_2)>=0`$
$`<C(z_1)D(z_2)>={\displaystyle \frac{1}{|z_{12}|^{2h}}}`$ (32)
$`<D(z_1)D(z_2)>={\displaystyle \frac{2\mathrm{ln}|z_{12}|^2+\delta }{|z_{12}|^{2h}}}`$
They have the following form of OPE with the stress-tensor $`T(z)`$:
$$T(w)C(z)=\frac{h}{(wz)^2}C(z)+\frac{1}{(wz)}\frac{C}{z}+\mathrm{}$$
(33)
$$T(w)D(z)=\frac{h}{(wz)^2}D(z)+\frac{1}{(wz)^2}C(z)+\frac{1}{(wz)}\frac{D}{z}+\mathrm{}$$
(34)
Or equivalently:
$$L_0|C>=h|C>L_0|D>=h|D>+|C>$$
(35)
In this case $`L_0`$ cannot be diagonalised and $`C,D`$ are degenerate. It is immediately apparent, from the form of the two point functions, that such theories cannot be unitary as they possess extra zero norm states which do not decouple from the physical spectrum. The hamiltonian $`L_0`$ forms an indecomposable Virasoro module. These extra null states can reproduce the structure of a full gauge theory despite the fact we only have scalars involved due to the symmetry of the OPE under $`DD+\lambda C`$. In this was discussed from the point of view of singletons. The relation between singletons and LCFT will be discussed later.
Correlation functions with logarithmic behaviour were found in the $`GL(1,1)`$ WZNW model in . There has been a lot of work on LCFTs , in particular the role of extended algebras in the general $`c_{p,q}`$ models and for the special $`c=2`$ case. The indecomposable representations of the Virasoro algebra have also been investigated . As with a general CFT these are only well defined if correlators obey crossing symmetry and single-valuedness as well as modular invariance on the torus. This has been shown to hold for some classes of LCFTs so far.
We can also get a similar structure arising if we have logarithms in $`x`$. In this case it is $`J_0^3`$ that is of Jordan normal form :
$$J_0^3|C>=m|C>J_0^3|D>=m|D>+|C>$$
(36)
In a more general case if we have degeneracy with the descendants of another primary then we will also have logarithms in the OPE.
We now wish to show how this general structure arises from our solutions.
### 5.2 Exact Logarithmic Solutions and OPE
* $`j=j_1=j_3j_2=j_4=0`$
$`\varphi _j(x_1,z_1)\varphi _0(x_2,z_2)\varphi _j(x_3,z_3)\varphi _0(x_4,z_4)=|z_{31}|^{2h}|x_{31}|^{2j}[A\left(\mathrm{ln}\right|{\displaystyle \frac{1x}{x}}|+`$ (37)
$`{\displaystyle \frac{2j1}{k2}}\mathrm{ln}\left|{\displaystyle \frac{1z}{z}}\right|)+B]`$
If we take the limit $`13`$ and $`24`$ which corresponds to $`x,z\mathrm{}`$ we can expand our solution:
$$\varphi _j(x_1,z_1)\varphi _0(x_2,z_2)\varphi _j(x_3,z_3)\varphi _0(x_4,z_4)=|z_{31}|^{2h}|x_{31}|^{2j}\left[1+\mathrm{}\right]$$
(38)
where … stands for subleading terms that vanish in the above limit. Then we find that:
$`\varphi _j(x_1,z_1)\varphi _j(x_3,z_3)`$ $`=`$ $`|z_{31}|^{2h}|x_{31}|^{2j}\left[1+\mathrm{}\right]`$ (39)
$`\varphi _0(x_2,z_2)\varphi _0(x_4,z_4)`$ $`=`$ $`1+\mathrm{}`$ (40)
We also find the correct two point functions if we take the correlation function of the above. Note we can always rescale operators so that the leading contribution has coefficient unity. The proper coefficients can only be obtained by consistency with higher order terms in the expansions. We do not attempt this here.
Now taking the $`12`$ and $`34`$ limit which corresponds to $`x,z0`$ we can again expand our solution:
$`\varphi _j(x_1,z_1)\varphi _0(x_2,z_2)\varphi _j(x_3,z_3)\varphi _0(x_4,z_4)=|z_{42}|^{2h}|x_{42}|^{2j}[\mathrm{ln}|x_{21}|+\mathrm{ln}|x_{43}|`$
$`2\mathrm{ln}|x_{42}|+{\displaystyle \frac{2j1}{k2}}(\mathrm{ln}|z_{21}|+\mathrm{ln}|z_{43}|2\mathrm{ln}|z_{42}|)\mathrm{}]`$ (41)
where we use $`\mathrm{ln}x=\mathrm{ln}x_{21}+\mathrm{ln}x_{43}\mathrm{ln}x_{31}\mathrm{ln}x_{42}=\mathrm{ln}x_{21}+\mathrm{ln}x_{43}2\mathrm{ln}x_{42}+\mathrm{}`$ and similarly for $`z`$. Again the subleading terms are ignored. We now find that:
$`\varphi _j(x_1,z_1)\varphi _0(x_2,z_2)`$ $`=`$ $`E_j^x(x_2,z_2)\mathrm{ln}|x_{12}|+E_j^z(x_2,z_2)\mathrm{ln}|z_{12}|+F_j(x_2,z_2)+\mathrm{}`$ (42)
$`\varphi _j(x_3,z_3)\varphi _0(x_4,z_4)`$ $`=`$ $`E_j^x(x_4,z_4)\mathrm{ln}|x_{34}|+E_j^z(x_4,z_4)\mathrm{ln}|z_{34}|+F_j(x_4,z_4)+\mathrm{}`$
Using our exact solution leads to the following 2-pt functions when $`j0`$:
$`E_j^aE_j^b=0,E_j^z(x,z)F_j(0,0)={\displaystyle \frac{2j1}{(k2)}}|z|^{2h}|x|^{2j}`$ (43)
$`E_j^x(x,z)F_j(0,0)=|z|^{2h}|x|^{2j}`$
$`F_j(x,z)F_j(0,0)=2|z|^{2h}|x|^{2j}\left({\displaystyle \frac{2j1}{k2}}\mathrm{ln}|z|+\mathrm{ln}|x|\right)`$
The fields $`E_j^z,F_j`$ form a Virasoro logarithmic pair which mix under the conformal transformations $`z\lambda z`$:
$$F_jF_j\mathrm{ln}\lambda E_j^z,E_j^zE_j^z$$
(44)
$`E_j^x,F_j`$ form a Kac-Moody logarithmic pair which mix under the $`SL(2,R)`$ transformations $`xϵx`$:
$$F_jF_j\mathrm{ln}ϵE_j^x,E_j^xE_j^x$$
(45)
If $`j=\frac{1}{2}`$ it is consistent at this level to set $`E_j^z=0`$.
* $`j_2=\frac{1}{2}j_1+j_3+j_4=\frac{1}{2}`$
We find
$`F(x,z)=|x|^{2j_1}|1x|^{2j_3}|z|^{j_1}|1z|^{2j_3}\left|{\displaystyle x^{2j_11}(x1)^{2j_31}𝑑x}\right|`$ (46)
The correlation function is then obtained using (11).
$`j_10`$
$`\varphi _{j_1}(x,z)\varphi _{\frac{1}{2}}(0,0)`$ $`=`$ $`|z|^{h_1h_2+h_J}|x|^{j_1j_2+J}O_J(0,0)+\mathrm{}`$ (47)
$`\varphi _{j_3}(x,z)\varphi _{j_4}(0,0)`$ $`=`$ $`|z|^{h_3h_4+h_J}|x|^{j_3j_4+J}O_J^{^{}}(0,0)+\mathrm{}`$
where $`J=j_1+\frac{1}{2}`$. Then:
$$O_J(x,z)O_J^{^{}}(0,0)=|z|^{2h_J}|x|^{2J}$$
(48)
$`j_1=0`$
$`\varphi _0(x,z)\varphi _{\frac{1}{2}}(0,0)`$ $`=`$ $`|z|^{h_1h_2+h_J}|x|^{j_1j_2+J}[C_J(0,0)\mathrm{ln}|x|+D_J(0,0)+\mathrm{}`$ (49)
$`\varphi _{j_3}(x,z)\varphi _{j_4}(0,0)`$ $`=`$ $`|z|^{h_3h_4+h_J}|x|^{j_3j_4+J}[E_J(0,0)\mathrm{ln}|x|+F_J(0,0)+\mathrm{}`$
where again $`J=j_1+\frac{1}{2}=\frac{1}{2}`$. Then:
$`C_J^aE_J=0,C_J(x,z)F_J(0,0)=D_J(x,z)E_J(0,0)=|z|^{2h_J}|x|^{2J}`$ (50)
$`D_J(x,z)F_J(0,0)=2|z|^{2h_J}|x|^{2J}\mathrm{ln}|x|`$
We find these OPEs and the ones from our other solutions are consistent. In particular we find:
$$\varphi _j(x,z)\varphi _0(0,0)=\varphi _j(x,z)\varphi _0(0,0)+\alpha E_j^x\mathrm{ln}|x|+\beta (2j1)E_j^z\mathrm{ln}|z|+\gamma F_j+\mathrm{}$$
(51)
Also for $`j_3,j_40`$
$`\varphi _{j_3}(x,z)\varphi _{j_4}(0,0)=`$ $`\varphi _{j_3}(x,z)\varphi _{j_4}(0,0)+`$ (53)
$`{\displaystyle \underset{J}{}}|z|^{h_3h_4+h_J}|x|^{j_3j_4+J}(C_{Jj_3j_4}O_J(0)`$
$`+\delta _{J,\frac{1}{2}}D_{j_3j_4}O_J^{^{}}(0)\mathrm{ln}|x|)`$
The relative normalisation of the coefficients $`C_{Jj_3j_4},D_{j_3j_4},\alpha ,\beta ,\gamma `$ should be determined by the consistency of the OPE with the $`SL(2,R)`$ Kac-Moody algebra. Although these OPEs are not complete they at least give us a hint of the possible structure and in particular some interesting behaviour at $`j=\frac{1}{2}`$. In $`SL(2,R)`$ WZNW $`j=\frac{1}{2}`$ is the pre-logarithmic operator that is the analogue of the puncture operator in Liouville. In $`AdS_3`$ it is also the critical point below which solutions in the bulk are not square integrable.
## 6 Symmetries of the Correlation function
The four point correlation function is fully determined by $`F(x,z)`$ and we now wish to comment on a few of the discrete symmetries of the $`SL(2,R)`$ model that are manifest in the KZ equation.
### 6.1 Reflection $`j1j`$
Under $`j1j`$ we find the Casimir $`j(j1)`$ is unchanged and thus fields with the same other quantum numbers may be expected to mix. It is interesting that we can also see this symmetry from the point of view of the KZ equation. To do this we differentiate the KZ equation $`12j_2`$ times then we find that
$$\frac{^{12j_2}}{x^{12j_2}}F_{j_1,j_2,j_3,j_4}(x,z)andF_{j_1,1j_2,j_3,j_4}(x,z)$$
(54)
satisfy the same KZ equation. This is due to the fact that, up to normalisation:
$$\frac{^{12j}}{x^{12j}}\varphi _j(x)=\varphi _{1j}(x)$$
(55)
This only makes sense if $`12jZ_+`$ which is exactly the case in which the series starting at $`j`$ and $`1j`$ can be indecomposable.
### 6.2 $`j\frac{k}{2}j`$
It can also be easily verified that:
$$(xz)^{kj_1j_2j_3j_4}z^\alpha (1z)^\beta F_{\frac{k}{2}j_3,\frac{k}{2}j_4,\frac{k}{2}j_2,\frac{k}{2}j_1}$$
(56)
and $`F_{j_1,j_2,j_3,j_4}(x,z)`$ obey the same KZ with appropriately chosen $`\alpha ,\beta `$. As was discussed in precisely this symmetry is related to spectral flow.
### 6.3 Strange $`Z_2`$ symmetry
It is very easy to see that under the transformation:
$`j_1`$ $``$ $`\stackrel{~}{j_1}={\displaystyle \frac{j_1}{2}}+{\displaystyle \frac{j_2}{2}}{\displaystyle \frac{j_3}{2}}+{\displaystyle \frac{j_4}{2}}`$
$`j_2`$ $``$ $`\stackrel{~}{j_2}={\displaystyle \frac{j_1}{2}}+{\displaystyle \frac{j_2}{2}}+{\displaystyle \frac{j_3}{2}}{\displaystyle \frac{j_4}{2}}`$
$`j_3`$ $``$ $`\stackrel{~}{j_3}={\displaystyle \frac{j_1}{2}}+{\displaystyle \frac{j_2}{2}}+{\displaystyle \frac{j_3}{2}}+{\displaystyle \frac{j_4}{2}}`$
$`j_4`$ $``$ $`\stackrel{~}{j_4}={\displaystyle \frac{j_1}{2}}{\displaystyle \frac{j_2}{2}}+{\displaystyle \frac{j_3}{2}}+{\displaystyle \frac{j_4}{2}}`$ (57)
$`𝒫`$ and $`𝒬`$ shift by a constant. If we consider:
$$z^{\frac{\alpha }{k2}}(z1)^{\frac{\beta }{k2}}F_{\stackrel{~}{j_1},\stackrel{~}{j_2},\stackrel{~}{j_3},\stackrel{~}{j_4}}$$
(58)
where $`\alpha =\frac{1}{2}(j_1j_2)^2+\frac{1}{2}(j_3j_4)^2,\beta =\frac{1}{2}(j_2j_3)^2+\frac{1}{2}(j_1j_4)^2`$.
Then we find that this satisfies exactly the same KZ equation as $`F_{j_1,j_2,j_3,j_4}(x,z)`$ and thus can, presumably by uniqueness of the solutions, be identified with it; at least up to some overall scale. This strange symmetry was noted earlier in in the Dotsenko-Fateev model for specific values of spin $`J_0^3`$. It was shown to follow from an integral identity in in the case of $`SL(2)`$ degenerate representations.
## 7 String theory on $`AdS_3`$
String theory in $`AdS_3`$, which is the $`SL(2,R)`$ group manifold, is described by the $`SL(2,R)`$ WZNW theory. There has been much early work discussing the spectrum and consistency of such a theory . There has been renewed interest recently due to the AdS/CFT correspondence . The duality is between string theory in the bulk of Anti-de-Sitter space (AdS) and a conformal field theory (CFT) on the boundary of the space-time. Bulk fields naturally couple to local operators in the boundary CFT. Correlation functions on the boundary can be computed from the bulk theory in the supergravity approximation by taking the classical tree level graphs for the bulk interactions. A precise recipe to compute this is given in . For the $`AdS_5`$ case for example see .
So far most work has been conducted in this supergravity approximation. For general $`AdS_n`$, this is because we do not understand how to describe the full string theory in such a background. For the case of $`AdS_3`$ however the worldsheet theory is described by the $`SL(2,R)`$ WZNW model and the boundary theory is a two dimensional CFT and so there is the possibility to go beyond this leading supergravity approximation. Fields are naturally classified according to the representation theory of $`SL(2,R)`$. It is also an extremely interesting case in which we have an exactly solvable string theory and thus we can study physics in this background, for instance the issues of unitarity in such a curved space.
In it was suggested that logarithmic operators in the worldsheet generate zero modes in the target space which restore the symmetries that are broken by the background. This was used in to describe D-brane recoil. For the latest developments in this area see and references therein. In the case of $`AdS_3`$ geometry, which is the near horizon limit of $`D1D5`$ system, these zero modes should restore the full Poincare invariance broken by the position of the branes. It would be interesting to see if our solutions could be interpreted in this context.
Short distance logarithmic singularities in four point correlation functions in the boundary CFT came as a surprise. Short distances in boundary theory are mapped to long distances in the bulk and vice versa via the familiar UV/IR relation . In the $`AdS_5/CFT_4`$ case, in which the boundary theory is $`𝒩=4`$ supersymmetric Yang-Mills, these singularities are interpreted as an outcome of perturbative expansion in terms of anomalous dimensions of unprotected multi-trace operators which correspond to multiparticle states in the bulk . In this case it is expected that as the boundary theory in this case is reasonably well understood that these logarithms are not due to logarithmic structure in $`𝒩=4`$ SYM. This also seems to be the case in $`AdS_3/CFT_2`$ as it has been shown that the spectrum of multiparticle states in $`AdS_3`$ agrees precisely with the multi-trace operators in boundary CFT . However in $`AdS_3`$ the grey body spectrum has a logarithmic singularity which cannot be explained by an anomalous dimension. In it was shown that this could not be due to a LCFT on the boundary with degenerate primary states. This paradox was solved in which suggested that a LCFT with degeneracy between a primary state and the descendant of another could give the calculated spectrum. It is thus still an open question as to the role of LCFT in the boundary theory particularly beyond the supergravity approximation.
The $`j=0`$ operators on the world sheet correspond to the non-normalisable scalar singletons in the bulk of $`AdS_3`$. In general $`AdS_n`$ these are special finite dimensional representations which lie at the limit of unitarity . It was shown in that if one considers singletons in the bulk then they would naturally induce a LCFT on the boundary with correlation functions given by (5.1). This form of the bulk action was also considered in but from a different point of view. As we have found logarithmic correlation functions for $`j=0`$ operators on the worldsheet there might be a relation between the possible logarithmic structure of the worldsheet and boundary LCFTs.
## 8 Conclusions and Discussion
We have found several solutions exhibiting logarithmic behaviour in both the Kac-Moody and Virasoro parts. These clearly show that the set of primaries is not complete and that extra fields can be generated in the OPE. At the free level no-ghost theorems have been proved for $`0<j<\frac{k}{2}`$, and null states are expected to decouple, but general statements about the fully interacting case cannot be made until the full field content has been established. There has been several attempts to compute correlation functions using other methods . An analysis of two point functions shows that $`j`$ should obey a more stringent bound $`\frac{1}{2}<j<\frac{k1}{2}`$ (See also for a discussion of this bound from a different perspective). All the exact logarithmic solutions found by us involve some operators outside this range. Correlation functions involving primaries with $`j<1/2`$ can be mapped into correlation functions with $`j>1/2`$ using reflection symmetry (54) and logarithms disappear. The special case $`j=\frac{1}{2}`$, which is at the boundary of discrete and continuous series, deserves further study . In the free field approximation was used, valid near to the boundary of $`AdS_3`$, and it was found that zero and negative norm states could be produced by interactions. It would be interesting to compare these different approaches. Our results seem to suggest that at the interacting level it is not possible to remove all the zero norm states as a logarithmic pair is created.
There are many aspects which we feel deserve further attention. It is not obvious which indecomposable representations are allowed in the $`SL(2,R)`$ model. The KZ equation for a non-compact group, unlike the compact case, does not reduce to an ordinary differential equation and so there are issues of uniqueness and boundary conditions. The bulk degrees of freedom corresponding to the logarithmic operators are presumably singletons and it would be interesting to see how these enter in the bulk action. It is also clear that these are perhaps not the most general logarithms that could occur as the Casimir $`J^2`$ is always diagonal in our representation which is not required for indecomposable representations. We hope to study some of these in future work.
## 9 Acknowledgements
We would like to thank I. Kogan, A.M. Tsvelik, and M.J. Bhaseen for interesting and helpful discussion. We would also like to thank J. Rasmussen for drawing our attention to one of his papers which also discussed the strange symmetry. The research of A.N is funded by PPARC studentship number PPA/S/S/1998/02610. Sanjay is funded by the Felix Scholarship (University of Oxford).
## Appendix A $`SL(2,R)`$ representation theory
For completeness we include the representations of $`SL(2,R)`$ for the primary fields. The representations are characterised by the eigenvalues of the quadratic Casimir
$$C_2=\frac{1}{2}(J_0^+J_0^{}+J_0^{}J_0^+)J_0^3J_0^3=j(j1)$$
(59)
and also the eigenvalue of $`J_0^3`$. This set includes the unitary representations as well as the non-unitary ones and also importantly the non-decomposable representations as all of these could contribute to the operator product expansion. The fact that not all representations can be written in a block diagonal (decomposable) form is due to the fact that more than one field may possess the same quantum numbers $`j`$ and $`m`$ and so in general $`J_0^3`$ can only be reduced to Jordan normal form rather than fully diagonal. The representations are all built starting from a primary field and then acting with $`J_0^\pm `$. They are
* Continuous representations:
This is an irreducible, infinite dimensional representation that does not have a highest or lowest weight state.The $`j`$ and $`m`$ eigenvalues are unrelated. A representation of this type has the states
$$\{|j,\alpha ;m:m=\alpha ,\alpha \pm 1,\alpha \pm 2,\mathrm{}\}$$
(60)
where $`J_0^3|j,\alpha ;m=m|j,\alpha ;m`$ and also $`j\pm \alpha \mathrm{ZZ}`$ so that the representation is unbounded in either direction. We shall only consider real valued of $`\alpha `$ and so without loss of generality we can take $`0\alpha <1`$ . For fixed $`C_2`$ and $`\alpha `$, the representations for each branch of (59) are equivalent. Imposing unitarity restricts us to two types of representations:
1. $`𝒞_j^\alpha `$ –“Principal continuous series” occurs for $`C_2<\frac{1}{4}`$; thus $`j=\frac{1}{2}+is`$ with $`s\mathrm{IR}`$
2. $`_j^\alpha `$–“Complementary series” occurs for $`C_2>\frac{1}{4}`$. Here $`j`$ is real with $`\frac{1}{2}<j<1`$ and $`j\frac{1}{2}<|\frac{1}{2}\alpha |`$.
* Discrete lowest weight representation.
This is an irreducible, infinite-dimensional lowest weight representation and exists for $`2j\mathrm{ZZ}^{}\{0\}`$. A representation of this type has the states
$$\{|j;m:m=j,j+1,j+2,\mathrm{}\}$$
(61)
where $`J_0^{}\{|j;j=0`$ and $`J_0^3|j;m=m|j;m`$
This representation will be unitary if $`j>0`$. For representations of the group $`SL(2,R)`$ then $`2j\mathrm{ZZ}^+`$ but for the universal cover $`j\mathrm{IR}^+`$.
The derivation of this representation in shows that it can be embedded in a reducible, nondecomposable representation where $`mj`$ is an arbitrary integer. We can act on these states with $`m<j`$ repeatedly using $`J_0^+`$ but once we get into the normal irreducible representation we cannot leave it with actions of $`J_0^{}`$ because of the highest weight state.
* Discrete highest weight representation:
This is an irreducible, infinite-dimensional highest weight representation and again exists for $`2j\mathrm{ZZ}^{}\{0\}`$. A representation of this type has the states
$$\{|j;m:m=j,j1,j2,\mathrm{}\}$$
(62)
where $`J_0^+\{|j;j=0`$ and $`J_0^3|j;m=m|j;m`$
This representation will be unitary if $`j>0`$. For representations of the group $`SL(2,R)`$ then $`2j\mathrm{ZZ}^+`$ but for universal cover again $`j\mathrm{IR}^+`$.
In this is embedded in a reducible, nondecomposable representation which contains $`m>j`$ similar to that for the lowest weight representations.
* Finite dimensional representations:
This is an irreducible, finite-dimensional representation. It occurs when $`2j\mathrm{ZZ}^{}\{0\}`$. A representation of this type has the states
$$\{|j;m:m=|j|,|j|1,|j|2,\mathrm{}|j|\}$$
(63)
where $`J_0^{}\{|j;|j|=0`$ , $`J_0^+\{|j;|j|=0`$ and $`J_0^3|j;m=m|j;m`$. The representation is only unitary in the case $`j=0`$, i.e. for the identity representation, also known as the singleton . It is contained in a reducible nondecomposable representation for which $`m`$ is arbitrary.
## Appendix B Exact solutions of the KZ equation
We can find particularly simple solutions to (16) if the terms in $`𝒫`$ and $`𝒬`$ with no derivatives in $`x`$ are zero. We thus look for solutions in which these terms vanish. However rather than do this naively we first make the substitution
$$F(x,z)=z^{\frac{\alpha }{k2}}(1z)^{\frac{\beta }{k2}}x^p(1x)^qG(x,z)$$
(64)
This is mostly redundant but does allow us to see the various crossing symmetries in the solutions.
The KZ equation (16) then becomes
$$(k2)\frac{G(x,z)}{z}=\frac{\stackrel{~}{𝒫}G(x,z)}{z}+\frac{\stackrel{~}{𝒬}G(x,z)}{z1}$$
(65)
We write $`\stackrel{~}{𝒫}=\stackrel{~}{𝒫}_{}+\stackrel{~}{𝒫}_0`$ where $`\stackrel{~}{𝒫}_{}`$ denotes the part of $`\stackrel{~}{𝒫}`$ with derivative terms in $`x`$ and $`\stackrel{~}{𝒫}_0`$ is the remainder. Of course $`\stackrel{~}{𝒫}_0`$ is just an algebraic expression whereas $`\stackrel{~}{𝒫}_{}`$ is a differential operator. Similarly we decompose $`\stackrel{~}{𝒬}=\stackrel{~}{𝒬}_{}+\stackrel{~}{𝒬}_0`$. Explicitly these are:
$`\stackrel{~}{𝒫}_{}`$ $`=`$ $`x^2(1x){\displaystyle \frac{^2}{x^2}}+x((j_1+j_2+j_3j_4+1+2q)x2(1x)(j_2+p)2j_1){\displaystyle \frac{}{x}}`$ (66)
$`\stackrel{~}{𝒫}_0`$ $`=`$ $`x(2j_2+p+q)(j_1+j_2+j_3j_4+p+q)\alpha +{\displaystyle \frac{q(j_1j_2j_3+j_4q)}{1x}}+const`$
$`\stackrel{~}{𝒬}_{}`$ $`=`$ $`x(1x)^2{\displaystyle \frac{^2}{x^2}}(1x)((j_1+j_2+j_3j_4+1+2p)(1x)2x(j_2+q)2j_3){\displaystyle \frac{}{x}}`$
$`\stackrel{~}{𝒬}_0`$ $`=`$ $`x(2j_2+p+q)(j_1+j_2+j_3j_4+p+q)\beta +{\displaystyle \frac{p(j_1j_2+j_3+j_4p)}{x}}+const`$
If we choose $`\alpha `$ and $`\beta `$ correctly to absorb the constant terms then we need not be concerned about these.
Although these expressions for general $`j_i`$ are quite complicated we can search for simple solutions using the ansatz $`G(x,z)=X(x)+Z(z)`$. Then
$$(k2)\frac{dZ(z)}{dz}=\frac{\stackrel{~}{𝒫}_{}X(x)}{z}+\frac{\stackrel{~}{𝒫}_0(X(x)+Z(z))}{z}+\frac{\stackrel{~}{𝒬}_{}X(x)}{z1}+\frac{\stackrel{~}{𝒬}_0(X(x)+Z(z))}{z1}$$
(67)
If we impose the vanishing of $`\stackrel{~}{𝒫}_0`$ and $`\stackrel{~}{𝒬}_0`$ then we have:
$$(k2)\frac{dZ(z)}{dz}=\frac{\stackrel{~}{𝒫}_{}X(x)}{z}+\frac{\stackrel{~}{𝒬}_{}X(x)}{z1}$$
(68)
As this is to be true for all $`z`$ we must have
$$\stackrel{~}{𝒫}_{}X=r\stackrel{~}{𝒬}_{}X=s$$
(69)
Thus if we seek solutions with this ansatz then we must impose the constraints $`\stackrel{~}{𝒫}_0=0`$ and $`\stackrel{~}{𝒬}_0=0`$ and then find solutions $`X(x)`$ that satisfy (69). Of course $`G(x,z)constant`$ is always a solution to (65). In the next section we search for less trivial ones.
### B.1 Solution to the constraints $`\stackrel{~}{𝒫}_0=0`$ and $`\stackrel{~}{𝒬}_0=0`$
Solving the constraints we find the following cases:
$`\begin{array}{ccccccc}Case& p& q& j_1& j_2& j_3& j_4\\ & & & & & & \\ 1& 0& 0& & 0& & \\ 2& 0& 2j_2& j_3j_2j_4& & & \\ 3& 0& j_2j_3+j_4& 0& & & \\ 4& 0& 0& & & j_4j_1j_2& \\ 5& j_3j_1j_2& j_1j_2j_3& & & & 0\\ 6& 2j_2& 0& & j_1j_3j_4& & \\ 7& j_1j_2+j_4& 0& & & 0& \\ 8& 2j_1& 2j_3& & & & j_2j_1j_3\end{array}`$
Also
$$\alpha =p(12j_12j_2p)+q(j_1+j_2+j_3j_4+q)2j_1j_2$$
(70)
$$\beta =q(j_1+j_2j_3j_4+1+2p)+p(2j_1+4j_22j_4+2p)+2j_2(j_1+j_2j_4)$$
(71)
Although each of these gives different solutions for $`F(x,z)`$ we find that the actual physical quantity, the correlation function (11), is the same in cases 1,3,5,7 and in 2,4,6,8 due to the prefactor multiplying $`F(x,z)`$ and thus these give us no new information. Therefore we concentrate on the first and fourth cases as these are simpler due to $`p=q=0`$ condition.
Thus we have
* $`j_2=0F(x,z)=X(x)+Z(z)`$
* $`j_3=j_4j_1j_2F(x,z)=z^{\frac{2j_1j_2}{k2}}(1z)^{\frac{2j_2j_3}{k2}}(X(x)+Z(z))`$
#### B.1.1 Case I: $`j_2=0`$
Then we have:
$`\stackrel{~}{𝒫}_{}X(x)`$ $`=`$ $`x^2(x1){\displaystyle \frac{^2X(x)}{x^2}}+((j_1+j_3j_4+1)x^22j_1x){\displaystyle \frac{X(x)}{x}}`$ (72)
$`\stackrel{~}{𝒬}_{}X(x)`$ $`=`$ $`x(1x)^2{\displaystyle \frac{^2X(x)}{x^2}}((j_1+j_3j_4+1)(1x)^22j_3(1x)){\displaystyle \frac{X(x)}{x}}`$ (73)
Rather than immediately solving $`\stackrel{~}{𝒫}_{}X=r`$,$`\stackrel{~}{𝒫}_{}X=s`$ we note that
$`(x1)\stackrel{~}{𝒫}_{}X(x)+x\stackrel{~}{𝒬}_{}X(x)`$ $`=`$ $`r(x1)+sx`$ (74)
$`(j_1+j_3+j_41)x(1x){\displaystyle \frac{X(x)}{x}}`$ $`=`$ $`r(x1)+sx`$ (75)
If $`j_1+j_3+j_410`$ then
$$\frac{X(x)}{x}=\frac{1}{j_1+j_3+j_41}\left(\frac{r}{x}+\frac{s}{x1}\right)$$
(76)
It now remains to find if there are any extra constraints if we force this to be a solution of (69). We find that for non-trivial solutions (i.e $`r,s0`$)
1. $`j_1=j_3=j_4=0=j`$
2. $`r=sj=j_1=j_3\frac{1}{2}j_4=0`$
3. $`r=0j=j_3=j_4\frac{1}{2}j_1=0`$
4. $`s=0j=j_1=j_4\frac{1}{2}j_3=0`$
The case in which $`j_i0`$ was the one previously found in .
The solution of KZ equation in all these case follows from (76) is of the form:
$$F(x,z)=r\left[\frac{\mathrm{ln}z}{k2}+\frac{1}{2j1}\mathrm{ln}x\right]+s\left[\frac{\mathrm{ln}(z1)}{k2}+\frac{1}{2j1}\mathrm{ln}(x1)\right]+C$$
(77)
If $`j_1+j_3+j_41=0`$ then we have $`r=s=0`$ and $`F(x,z)`$ has no $`z`$ dependence.
$`\stackrel{~}{𝒫}_{}X(x)`$ $`=`$ $`0\stackrel{~}{𝒬}_{}X(x)=0`$ (78)
$`{\displaystyle \frac{X(x)}{x}}`$ $`=`$ $`(1x)^{2j_3}x^{2j_1}`$ (79)
$`F(x,z)`$ $`=`$ $`{\displaystyle (1x)^{2j_3}x^{2j_1}𝑑x}`$ (80)
This has logarithmic behaviour in $`x`$ if $`2j_1`$ or $`2j_3Z^+`$. However we see that the logarithmic behaviour in $`z`$ has vanished. If say $`j_4=0`$ this can be thought of as the limiting case of the above solutions when we reach the special values $`j_1=j_3=\frac{1}{2}`$.
#### B.1.2 Case II: $`j_3=j_4j_1j_2`$
We can repeat the same arguments to find non-trivial solutions. We find only the same solutions as before with logarithms in both $`x`$ and $`z`$. However we can find solutions with no logarithms in $`z`$ similar to (80) when $`j_4=\frac{1}{2}`$:
$$F_{j_1,\frac{1}{2}j_1j_3,j_3,\frac{1}{2}}=z^{2j_1j_2}(1z)^{2j_2j_3}x^{2j_31}(1x)^{2j_11}𝑑x$$
(81)
A rearrangement of this solution is analysed in (46). It is simpler however to observe that this is related to the previous case B.1.1 by the strange symmetry (6.3). If $`j_4=j_2j_1j_3`$ then we find this gives the same $`F(x,z)`$ as
$`j_1`$ $``$ $`j_1+j_2`$ (82)
$`j_2`$ $``$ $`0`$ (83)
$`j_3`$ $``$ $`j_1+j_4`$ (84)
$`j_4`$ $``$ $`j_2+j_4`$ (85)
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# Fate stochastic management and policy benchmark in 421, a popular game
## 1 Aim and interest of the study
Following , I look for strategies, maximizing the probability of win, or some expected utility, in the game 421 combining chance and decision (see appendix A).
By the way, or, indeed, by *serendipity*, I encounter the problem of fate *stochastic management*: optimizing today’s decisions, with respect to a future utility, and in spite of tomorrow’s odds. Such issues as what are the optimal policies, in what circumstances, and how much they demand on intellectual resources, can be resolved mathematically, suggesting that management could be an exact science (as part of operations research).
A lottery is not a game, in the sense of game theory, but a stochastic process (a sequence of random variables). Game theory treats classically multiple player decision games, the archetype of which is chess. A game in which players’ fates depend on both chance and their decisions, like 421 and backgammon , is a *stochastic game* . Chance makes decision more complex. For example, consider a variation on chess: the player proposes a list of $`n`$ moves, the actual one being determined by casting dice. $`n=1`$ yields the standard pure decision game; $`n=n_1`$, where $`n_1`$ is the number of possible moves, yields a pure chance game; $`nn_1/2`$ yields a game of chance and decision, more complex than the former and the latter.
Game theory primarily focuses on the existence proofs for optimal strategies. However, “usable techniques for obtaining practical answers” also matter \[4, §1.1\]. Indeed, little can be done from existence without construction: this is the old debate around Zermelo’s axiom of choice. Hence the interest of investigating, as in Church’s thesis , calculability, the existence of an algorithmic solution. But even calculability may not be sufficient for actual computation. For example, consider again chess, a finite but very large game: the algorithmic solution provided by the Zermelo theorem \[6, ch. 6\] is of no practical use (until the final moves), as noticed by \[7, §11.4\], as it exceeds the capacity of any computer. The study of finite games does not stop with Zermelo theorem, and this is because of complexity boundedness. Algorithms shall be compared not only with respect to optimality (degree of completion of the task) but also complexity, using a bit of *complexity theory* .
An algorithm is characterized by its optimality, size and computing time on a given computer, specialized by high-level functions and data. The algorithm may be good or bad, short or long, fast or slow. The three qualities and quantities are not independent: the exchange of computing time against size is the principle of data compression, the exchange of computing time against optimality is the principle of heuristics. When setting two quantities, the minimum of the third one, as a function of other implicit parameters, can be defined: the minimum size and the minimum computing time are respectively related with the Kolmogorov complexity and the Bennett logical depth .
For strategy-generating algorithms, or deciding algorithms, or *policies*, optimality is the expected utility. For a perfect information stage game, it is interesting, for possible extensions, to characterize the asymptotic behavior of the computing time, when the depth tends to infinity, i. e. to know whether the algorithm is linear or polynomial, rather than exponential, as feared from the tree structure of the game.
The present study is thus an occasion to relate with each other, on a live case, various tools and concepts attached to games, processes, probabilities, control, programming, algorithmics and complexity, with applications in management and game practice.
## 2 Backward induction optimal policy
A stochastic game reduces to a pure decision game, by considering providence as a particular player \[8, ch. 4\], whose mixed strategy, known a priori, results from usual statistical postulates (independence, stationarity) and cannot be optimized.<sup>1</sup><sup>1</sup>1Probability theory began as the study of the providential strategy in chance games, at the time of Bayes or the Bernoullis. 421, thus reduced, and with some precautions on the rules (appendix A), is a perfect information finite game and the Zermelo theorem applies.
I will solve only a sub-game, the player’s round against providence (while other players stand still), a stochastic management problem, featuring a martingale problem and, for the first player, a stopping time problem . The analogy with Brownian motion provides statistical mechanics tools.
### 2.1 Alea
Let $`D`$ be the number of dice, normally 3, and $`F^{}`$ the number of faces of every dice, normally 6. Dice are discernible<sup>2</sup><sup>2</sup>2Discernibility is not an innocuous hypothesis, as shown by Gibbs’ paradox ., so that the probability space is the set of face sequences, or arrangements. The class of arrangements corresponding to each other by a permutation is a combination, e. g., “nénette”, 221, is the subset of arrangements $`\{(1,2,2),(2,1,2),(2,2,1)\}`$, of redundancy three. In 421, which dice produced which face does not matter, because ranking depends only on combination.
I describe the die system as in statistical mechanics: each die is a particle, with only one *phase* variable, *face*. The laws of mechanics are replaced by usual statistical hypothesis, abstracting chance from any specific random generator. The system is described, in Lagrangian notation, by a face combination, or, in Eulerian notation, by the sequence $`d_f`$ of occupation numbers of every face $`f=1\mathrm{}F`$, e. g. the Lagrangian notation 421 translates into the Eulerian notation $`(1,1,0,1,0,0)`$ ($`F=6`$).
The interest of Eulerian notation lies in that the set of Eulerian combinations
$$𝐝=(d_f)_{f=1\mathrm{}F}.$$
is the partially ordered normed vector space $`^F`$. The canonic basis $`𝐞_g=(\delta _{f,g})_f`$ is aligned with “brelans”, combinations with all faces of a kind. I define the ball
$$B(D)\{𝐝^F,|𝐝|D\}$$
and similarly (replacing $``$ above by $`=`$ or $`<`$ ) the sphere $`B(D)`$ and the open ball $`\stackrel{˘}{B}(D)`$. The intersections with the positive cone are represented by $`+`$ exponents; the set of actual combinations is $`B^+(D)`$. The norm of a combination is the sum of Eulerian component absolute values. The norm of a positive combination is just its number of dice. The canonic order $``$, partial on $`^F`$, differs from the hierarchic order $``$ (54), total on $`B^+(D)`$.
Distinct casts are independent and the probability of any face to be on top is $`1/F`$ (unloaded dice). The arrangements of one combination are thus equiprobable, and the probability of a combination is just that of any of its arrangements, times the combination redundancy. For example, the probability of obtaining the combination 21 is $`2/F^2`$, while the probability of obtaining the combination 11 is $`1/F^2`$. More generally, the probability of obtaining the combination $`𝐝`$, after one cast, is given by the multinomial law, with usual notations generalizing power and factorial to integer vectors:
$$p(𝐝)=𝐩^𝐝\frac{|𝐝|!}{𝐝!},𝐩=\frac{1}{F}(1\mathrm{}1)^F,\underset{𝐝B^+(D)}{}p(𝐝)=1.$$
(1)
### 2.2 Fate
For all $`j`$, let the state $`𝐝_j`$ be the combination, accumulated after $`j`$ casts, and the event $`𝐝_{j+1/2}`$ be the combination, obtained from the $`j+1`$-th cast. Fate is the infinite state and event alternate sequence
$$\phi (𝐝_0,𝐝_{1/2},𝐝_1,𝐝_{3/2}\mathrm{}).$$
(2)
The integer or half-integer index is used as a discrete time, integer time for states, half-integer for events. The set of possible fates is described by the fate tree, where branching represents chance (from integer time to half-integer time) or decision (conversely).
The rules of 421 imply:
$`𝐝_0`$ $``$ $`\mathrm{𝟎},`$ (3)
$`j,𝐝_{j+1/2}`$ $``$ $`B^+(D_jD|𝐝_j|),`$ (4)
$`0𝐝_{j+1}𝐝_j`$ $``$ $`𝐝_{j+1/2},`$ (5)
$`(j,jJ),𝐝_j`$ $``$ $`B^+(D),`$ (6)
where $`J`$ is the maximum round duration, normally 3. $`D_j`$ is the number of live dice, which have not been accumulated after $`j`$ events and one state. From (3, 6), $`D_0=D,D_J=0`$.
From (4, 5),
$`D_{j+1}`$ $``$ $`D_j,`$ (7)
$`(𝐝_{j+1}𝐝_j=𝐝_{j+1/2})`$ $``$ $`𝐝_{j+1}B^+(D).`$ (8)
The effective round duration $`J_1`$ is the minimum of $`j`$ in (6). The next players’ effective round durations must equal the first player’s. Therefore, for all players,
$`(j,j<J_1),𝐝_j\stackrel{˘}{B}^+(D)`$ , $`𝐝_{j+1/2}0,`$ (9)
$`𝐝_{J_1}B^+(D)`$ , (10)
$`(j,j>J_1),𝐝_{j1/2}=0`$ , $`𝐝_j=𝐝_{J_1}.`$ (11)
(9, 10) are used, firstly, after the first player’s end of round, to determine $`J_1`$, subsequently, as additional rules for next players. When $`j`$ increases, the state vector $`𝐝_j`$ moves in the positive ball, off the origin, towards its boundary where it gets stuck at $`𝐝_{J_1}`$, the round *result*. Fate is virtually continued by an infinite sequence, asymptotically alternating the result and the null event.
### 2.3 Utility
Following von Neumann and Morgenstern \[7, ch. 27\], a player’s utility is a number, given by a causal function, i. e. a function of history (past fate), compatible with the player’s preferences, and such that the utility before a random event is just the expected utility, i. e. the probability-weighted utility average, over possible outcomes. Thus, expected utility is anti-causal, i. e. prescribing utilities at some future time determines its expectation at all prior times.
One never knows when a game actually stops, as a it is often embedded in a larger game. Tennis is a familiar example: a tennis “game” is a sub-game of a set, itself a sub-game of a match, tournament, ranking system…this cascade does not even stop with a player’s life, because of cooperation between individuals. But, if we want to obtain any result, we must stop somewhere in the game cascade, and judge utility more or less empirically. (Quite similarly, in mechanics or thermodynamics, the studied system is coupled with the rest of the world, by an often delicate boundary or cut-off condition.)
The study of 421 should stop at end of game, by setting players’ utilities, for example, a binary utility: one for win, zero for loss, or incorporating economy, à la Bernoulli, the logarithm of earning divided by wealth . However, I treat only the round. At end of round, the Bernoulli formula does not make sense and utility is not given directly by the rules (in particular, the transfer function of table 5). By examining the rules, a few properties of utility are obtained; for example, at constant time, for a rational player, utility must be compatible with the hierarchic order (54), etc.
But I will not further characterize utility. On the contrary, I will consider the round independently of the rest of the game, with arbitrary utilities, in order to treat the problem of fate stochastic management in a rather general way.
For all fate $`\phi `$ (2), utility is judged at some time $`J_\phi `$, either integer or half-integer in general (in the round, $`J_\phi \{1/2,3/2,5/2\}`$), as a causal function:
$$u(𝐝_0,𝐝_{1/2}\mathrm{}𝐝_{J_\phi }).$$
(12)
The function $`u`$ has a variable number of arguments, formally, it is defined on $`_j^{}B^+(D)^j`$. Utility is judged forever:
$$u(\mathrm{}𝐝_{J_\phi },.)u(\mathrm{}𝐝_{J_\phi }).$$
(13)
The rules (4, 5, 6, 9, 10) are superseded by $`\mathrm{}`$ utilities for rule breaking histories (excluding cheating). In particular, the utilitarian version of (4) is
$$(j,𝐝_{j+1/2}B^+(D_j)),u(\mathrm{}𝐝_j,𝐝_{j+1/2})=\mathrm{}$$
(14)
and the next players’ round duration conditions (9, 10) become
$$(j,j<J_1,𝐝_jB^+(D)),u(\mathrm{}𝐝_j)=\mathrm{}.$$
(15)
### 2.4 Optimal strategies
The greatest utility, drawn from any event-terminated history, is
$`j,u(\mathrm{}𝐝_j,𝐝_{j+1/2})`$ $`=`$ $`\underset{𝐝_{j+1}}{\mathrm{max}}u(\mathrm{}𝐝_j,𝐝_{j+1/2},𝐝_{j+1}),`$ (16)
$`u(.)`$ $`=`$ $`\underset{𝐝_1}{\mathrm{max}}u(.,𝐝_1).`$ (17)
The latter equation, where $`𝐝_1`$ is a dummy variable, is a more formal expression of the former. The nature of the dummy variable is shown by its index (state for integer, event for half-integer). The set of states, corresponding to optimal decisions, is
$$S_u(.)\underset{𝐝_1}{argmax}u(.,𝐝_1).$$
(18)
A player’s mixed strategy consists in choosing randomly between many decisions, according to a causal probability law,
$$𝐝P(.,𝐝)𝒫(𝐝_1=𝐝|.),\underset{𝐝_1}{}P(.,𝐝_1)=1.$$
(19)
$`𝒫(X)`$ means the probability of the event $`X`$. The optimal mixed strategies are such that the support of the probability law (19) is a subset of $`S_u(.)`$ (among them are pure optimal strategies).
From the von Neumann-Morgenstern theorem,
$$u(.)=\underset{𝐝_{1/2}}{}p(.,𝐝_{1/2})u(.,𝐝_{1/2}).$$
(20)
where $`p`$ is a causal probability law, expressing the providential strategy and rules. Because of utility conditions, such as (14, 15), there are, in (20), products $`p\times u`$ of the undetermined form $`0\times \mathrm{}`$, which ought to be replaced by zero (or the summation ought to be properly restricted).
Combining (17, 20), or conversely,
$`u(.)`$ $`=`$ $`\underset{𝐝_1}{\mathrm{max}}{\displaystyle \underset{𝐝_{3/2}}{}}p(.,𝐝_{3/2})u(.,𝐝_1,𝐝_{3/2}),`$ (21)
$`u(.)`$ $`=`$ $`{\displaystyle \underset{𝐝_{1/2}}{}}p(.,𝐝_{1/2})\underset{𝐝_1}{\mathrm{max}}u(.,𝐝_{1/2},𝐝_1).`$ (22)
The composition of $`\mathrm{max}\mathrm{moy}`$ operations names the algorithm, which is the classical zero-sum game $`\mathrm{max}\mathrm{min}`$, where the rational opponent has been replaced by neutral providence. (21, 22) are consistent with (13): after the judgment, they simply repeat the utility forever, so that the $`\mathrm{max}\mathrm{moy}`$ operations can be chained ad infinitum, no matter the end of round. Thus, the judgment can be arbitrarily postponed, without affecting strategy. If judgment times have an upper bound (e. g. the number of fates is finite), then all judgments can be postponed until a (collective) last judgment a time $`J,J\mathrm{max}J_\phi `$, e. g. the maximum round duration.
Relaxing the rule (3), and taking $`𝐝_0`$ as a parameter, the problem of fate management, i. e. finding optimal strategies, is *self-similar* under time-shifts, except for the parameter “renormalization” (as in statistical physics)
$$(D_0,J,𝐝_0,(\mathrm{}𝐝_j),(p,u)(.))(D_j,Jj,𝐝_j,(),(p,u)(\mathrm{}𝐝_j,.)).$$
(23)
Let $`\chi _V`$ be the characteristic function of $`VB(D)`$. The round providential strategy is determined by (1) and
$$p(.,𝐝_0,𝐝_{1/2})p(𝐝_{1/2})\chi _{B^+(D_0)}(𝐝_{1/2}).$$
(24)
The expected utility is computed with (17, 20), from the last judgment backward in time:
$$u(\mathrm{}𝐝_{J1},𝐝_{J1/2}),u(\mathrm{}𝐝_{J1})\mathrm{}u(𝐝_0,𝐝_{1/2},𝐝_1),u(𝐝_0,𝐝_{1/2}),u(𝐝_0),$$
(25)
e. g., for $`J=3`$, and using (24),
$`u(\mathrm{}𝐝_{3/2})=\underset{𝐝_2}{\mathrm{max}}`$ $`{\displaystyle \underset{𝐝_{5/2}B^+(D_2)}{}}p(𝐝_{5/2})u(\mathrm{}𝐝_{3/2},𝐝_2,𝐝_{5/2}),`$
$`u(𝐝_0,𝐝_{1/2})=\underset{𝐝_1}{\mathrm{max}}`$ $`{\displaystyle \underset{𝐝_{3/2}B^+(D_1)}{}}p(𝐝_{3/2})u(𝐝_0,𝐝_{1/2},𝐝_1,𝐝_{3/2}).`$
$`u(𝐝_0)=`$ $`{\displaystyle \underset{𝐝_{1/2}B^+(D_0)}{}}p(𝐝_{1/2})u(𝐝_0,𝐝_{1/2}).`$
## 3 Fate as a stochastic process
For a given strategy, what is the presence density (of the die system in a subset of phase space)? What is the expectation of an arbitrary utility, for which the given strategy is not necessarily optimal?
### 3.1 The Kolmogorov equation on expected utility
Fate is a stochastic process, not only because it contains random events (the probability law $`p`$), but also random decisions, according to mixed strategies (the probability law $`P`$). For any causal process like (2), the sequence of histories
$$(𝐝_0),(𝐝_0,𝐝_{1/2}),(𝐝_0,𝐝_{1/2},𝐝_1)\mathrm{}$$
is a discrete Markov chain, for which classical results are available \[12, ch. 6\], \[13, ch. 15\], originating mostly from Brownian motion studies \[10, ch. 15\].
The fate stochastic evolution equation, the Langevin equation, is just a random sum, obeying (4, 5):
$`𝐝_{j+1}`$ $`=`$ $`𝐝_j+\widehat{𝐝}_{j+1},`$
$`𝒫(\widehat{𝐝}_{j+1}=𝐝|\mathrm{}𝐝_j,𝐝_{j+1/2})`$ $`=`$ $`P(\mathrm{}𝐝_j,𝐝_{j+1/2},𝐝_j+𝐝),`$
$`𝒫(𝐝_{j+1/2}=𝐝|\mathrm{}𝐝_j)`$ $`=`$ $`p(\mathrm{}𝐝_j,𝐝).`$
$`\widehat{𝐝}_{j+1}`$ is a random source term, conditioned by history, according to the mixed strategies $`P,p`$. $`𝐝_j`$ undergoes a strategy-driven Brownian motion as, for example, a charged Brownian particle driven by electrophoresis.
The Chapman-Kolmogorov equation yields the probability of transition, or jump, from one state to the other, in one time step:
$$\sigma (.,𝐝_0𝐝_1)𝒫(𝐝_1=𝐝|.,𝐝_0)=\underset{𝐝_{1/2}}{}p(.,𝐝_0,𝐝_{1/2})P(.,𝐝_0,𝐝_{1/2},𝐝).$$
(26)
Let $`P`$ be a player’s mixed strategy, possibly not optimal. From the von Neumann-Morgenstern theorem, twice applied,
$$u(.,𝐝_0)=\underset{𝐝_{1/2}}{}p(.,𝐝_0,𝐝_{1/2})\underset{𝐝_1}{}P(.,𝐝_0,𝐝_{1/2},𝐝_1)u(.,𝐝_0,𝐝_{1/2},𝐝_1).$$
(27)
Reversing the order of summation, using (26) and assuming that *utility does not depend on events, but only on states*, which is true in the 421 round, I obtain the Kolmogorov equation on the expected utility:
$$u(.,𝐝_0)=\underset{𝐝_1}{}\sigma (.,𝐝_0𝐝_1)u(.,𝐝_0,,𝐝_1).$$
(28)
(By hypothesis, $`u`$ does not depend on $``$.)
As opposed to the $`\mathrm{max}\mathrm{moy}`$ algorithm, (28) does not produce any decision, but, given the mixed strategies $`P,p`$ (effective through $`\sigma `$), determines the expectation of any utility, for which $`P`$ may not be optimal.
Nevertheless, if $`P`$ is optimal, from (17) and (19), there is an equality, between operators on $`u(.,𝐝_0,𝐝_{1/2},𝐝_1)`$:
$$\underset{𝐝_1}{\mathrm{max}}=\underset{𝐝_1}{}P(.,𝐝_0,𝐝_{1/2},𝐝_1).$$
(29)
Taking (29) into (27) returns (22).
### 3.2 The Fokker-Planck equation on presence density
I define the state fate $`\psi (𝐝_j)_{j=0,1\mathrm{}}`$ (fate with only states, not events). From (26),
$$𝒫(\psi =(.,𝐝_0,𝐝_1))=\sigma (.,𝐝_0𝐝_1)𝒫(\psi =(.,𝐝_0)),$$
(30)
so that the sequence of past states
$$(𝐝_0),(𝐝_0,𝐝_1),(𝐝_0,𝐝_1,𝐝_2)\mathrm{}$$
also is a Markov chain.
Summing (30) over all state fates converging to the same state $`𝐝`$ at time $`j+1`$ gives the presence density $`\rho _{j+1}(𝐝)`$:
$`\rho _0(𝐝)`$ $`=`$ $`\delta _{𝐝,𝐝_0},`$ (31)
$`j,\rho _{j+1}(𝐝)`$ $`=`$ $`{\displaystyle \underset{𝐝_0\mathrm{}𝐝_j}{}}𝒫(\psi =(𝐝_0\mathrm{}𝐝_j)).`$ (32)
In the round, from (11), $`\rho _j`$ is stationary, as soon as $`jJ`$.
I assume that *utility is a function of state and time only*, less general than causal (12):
$$u(\mathrm{}𝐝_{J_\phi })=u_{J_\phi }(𝐝_{J_\phi }).$$
(33)
The end-of-round utility is indeed of the kind (33), because end-of-set ranking (see the rules) only depends on round results, not on intermediary states and, for the first player, the effective round duration.
For all player’s optimal (or rational) mixed strategy $`P`$ derived from a utility of the kind (33),
$`P_j(𝐝_j,𝐝_{j+1/2},𝐝_{j+1})`$ $``$ $`P(.,𝐝_j,𝐝_{j+1/2},𝐝_{j+1}),`$ (34)
$`\sigma _j(𝐝_j𝐝_{j+1})`$ $``$ $`\sigma (.,𝐝_j𝐝_{j+1})`$ (35)
Taking (33, 35) into (28) allows to extend (33) to all time (for the expected utility), by induction:
$$j,u_j(𝐝_j)u(\mathrm{}𝐝_j).$$
(36)
The process $`(j,𝐝_j,𝐝_{j+1/2})`$ is Markovian.
The consequence (35) of (33), taken into (32), allows to express $`\rho _{j+1}(𝐝_{j+1})`$ as a functional on $`\rho _j`$:
$$\rho _{j+1}(𝐝_{j+1})=\underset{𝐝_j}{}\rho _j(𝐝_j)\sigma _j(𝐝_j𝐝_{j+1}),$$
(37)
the Fokker-Planck equation.
As opposed to (37), (28) does not need (33). Nevertheless, with (33), (28), becomes
$$u_j(𝐝_j)=\underset{𝐝_{j+1}}{}\sigma _j(𝐝_j𝐝_{j+1})u_{j+1}(𝐝_{j+1}),$$
(38)
adjoint to (37).
(37, 38) are the evolution equations, adjoint to each other, linear, unstationary, of presence density and expected utility. Their inputs are a player’s mixed strategy and utility.
### 3.3 Computing result probabilities by duality
Let $`(B^+(D),)`$ be the space of numerical functions on $`B^+(D)`$, with the scalar product
$$(f,g),f,g\underset{𝐝B^+(D)}{}f(𝐝)g(𝐝).$$
(39)
$`\sigma _j`$ is an operator, a linear endomorphism on $``$, fully determined by the Markovian matrix $`\sigma _j(𝐝_0𝐝_1)`$. Its transposed operator is $`\sigma _j^t`$, of matrix
$$\sigma _j^t(𝐝_1𝐝_0)\sigma _j(𝐝_0𝐝_1).$$
In operator notation, (37, 38) become
$$\rho _{j+1}=\sigma _j^t\rho _j,u_j=\sigma _ju_{j+1}.$$
As $`\sigma _j^t`$ et $`\sigma _j`$ are adjoint to each other, the expected utility follows a conservation law:
$`u_j,\rho _j=\sigma _ju_{j+1},\rho _j`$ $`=`$ $`u_{j+1},\sigma _j^t\rho _j=u_{j+1},\rho _{j+1},`$
$`u_j,\rho _j`$ $`=`$ $`u_0,\rho _0=u_0(𝐝_0).`$ (40)
The last equality is a consequence of (31). Given the player’s mixed strategy $`P`$, (40) holds for any utility.
The direct computation of $`u_j,\rho _j`$ consists in solving for $`\rho _j(𝐝_j)`$ the Fokker-Planck equation, which must be repeated, to complete the scalar product, at least for all $`𝐝_j`$ where $`u_j`$ does not vanish. More shrewdly, $`u_j,\rho _j`$ can be computed indirectly, as the r. h. s. of (40): the Kolmogorov equation is solved only once for the expected utility at the trunk of the fate tree, or the initial expected utility. The indirect computation is faster than the direct computation, by a factor which is the cardinal of the support of $`u_j`$. The indirect computation benefits from the unicity of the fate tree, and the diffusive growth of the support of $`u_j`$.
Moreover, to obtain the Kolmogorov algorithm from the $`\mathrm{max}\mathrm{moy}`$ algorithm, one merely has to replace, in (22), the operator $`\mathrm{max}`$ appearing at the l. h. s. of (29), by the operator $`\mathrm{strat}`$ appearing at the r. h. s. of (29). (These operators differ if $`P`$ is not optimal.) The Kolmogorov equation is thus solved by a $`\mathrm{strat}\mathrm{moy}`$ algorithm.
Here are examples of using the Kolmogorov equation and (40):
1. The probability of the result to be in $`VB^+(D)`$ (independently of time) is the initial expectation of the stationary utility $`u_j=\chi _V`$.
2. The probability of $`D_{j_0}`$ is the initial expectation of the utility $`u_j=\delta _{j,j_0}\chi _{B^+(DD_{j_0})}`$.
### 3.4 Analogy with linear transport theory
The round is a linear transport phenomenon, with respect to the face variable. Face, expected utility, presence density, transition probability correspond respectively, in transport theory , to phase (position, velocity), importance , flux and cross section. Harris shows that a monokinetic particle population grown by branching (e. g. neutrons produced by nuclear fission) follows a Galton-Watson process. Similarly, in appendix B, I discuss the Galton-Watson character of the first player’s live dice population $`D_j`$.
## 4 Simple optimal policies for one-goal utilities
Taking for goal a unique combination $`𝐝^{}B^+(D)`$, the utility is a binary Kronecker function $`\delta _{𝐝^{},.}`$ (modulo an affine transform), and optimal strategies are simply constructed.
### 4.1 The ratchet and Bernoulli policies
I examine two first player’s policies, with a one-goal utility:
1. The Bernoulli policy consists in accumulating no die, unless the goal has been attained (then, all dice are accumulated); the cast sequence is a stationary Bernoulli process (a sequence of independent trials terminated by success or failure).
2. The ratchet policy consists in putting aside as many dice as possible, contributing to the goal:
$`(j,j+1<J),𝐝_{j+1}`$ $`=`$ $`𝐝^{}(𝐝_j+𝐝_{j+1/2}),`$ (41)
$`P_j(𝐝_j,𝐝_{j+1/2},𝐝_{j+1})`$ $`=`$ $`\delta _{𝐝_{j+1},𝐝^{}(𝐝_j+𝐝_{j+1/2})}.`$ (42)
($``$ is the infix notation of the minimum in the partially ordered space $`^F`$, generalizing, in Lagrangian notation, the ensemble intersection $``$.)
The ratchet strategy towards $`𝐝^{}`$ is optimal, with respect to the $`𝐝^{}`$-goal utility, if and only if $`p`$ decreases in $`B^+(D)`$. This means that as many dice as possible should be accumulated, in order to maximize the success probability at any future time. For unloaded dice, from (1),
$$(𝐝B^+(F),𝐝+𝐞_1B^+(F)),\frac{p(𝐝+𝐞_1)}{p(𝐝)}=\frac{1}{F}\frac{|𝐝|+1}{d_1+1}1,$$
(43)
i. e. $`p`$ decreases on $`B^+(F)`$. The ratchet strategy is optimal if and only if $`DF`$, strictly if and only if $`D<F`$.
For example, with $`D=3<F=6,J>1,𝐝^{}=421,𝐝_{1/2}=651`$, the ratchet decision (to accumulate 1) is optimal, because $`p(421)<p(42)`$ (it will be easier to obtain 42 than 421). With $`D=3>F=2,𝐝^{}=211,𝐝_{1/2}=222`$, the Bernoulli decision (to replay all dice) is optimal, because $`p(11)=1/F^2=2/8<p(211)=3/F^3=3/8`$. With $`D=F=2,𝐝^{}=21,𝐝_{1/2}=11`$, both Bernoulli and ratchet decisions are optimal.
A next player’s maximum round duration is imposed. In case of a premature success, he is in a *dilemma*, having to decide between equally unpleasant ways of breaking the goal, obtained too early. For $`F>2`$, optimal decisions consist in replaying any one die; the number of pure optimal strategies is thus the number of distinct faces in the goal combination, at the power $`J1`$. If the goal is a brelan, then no dilemma exists.
### 4.2 Optimal one-goal strategy result probabilities
For any strategy, I consider the probability to obtain any result, e. g. 111 after three casts. According to section 3.3, this probability is the initial expected utility, determined by the Kolmogorov equation and the final condition of a Kronecker utility on the result. This probability depends on the player $`i=1,2`$ (first or next), the (renormalized) maximum round duration $`J_1`$, the player’s mixed strategy $`P`$, the delay $`j`$, and the result $`𝐝`$:
$$p_i(J_1,P,j,𝐝),0jJ_1J,𝐝B^+(D).$$
(44)
The set of result probabilities, for all possible pure strategies and $`(D,F,J)=(3,6,3)`$, is (much larger than the fate tree, itself very large and) too large to be extensively listed. Thus, I will work on a reduced strategy subset, for which a reasonable choice is the set of optimal one-goal strategies, for all possible goals. As far as the goal determines the optimal strategy, the variable $`P`$ in (44) is simply replaced by the goal $`𝐝^{}`$:
$$p_i(J_1,𝐝^{},j,𝐝),0jJ_1J,(𝐝,𝐝^{})B^+(D)^2$$
(45)
which looks like the Markovian matrix of section 3.3, except that $`𝐝^{}`$ is not actual, but contemplated. There are diagonal ($`𝐝=𝐝^{}`$) and non-diagonal result probabilities.
For the first player, the optimal one-goal strategy is unequivocally defined by the goal ($`D<F`$: the ratchet) and the function $`p_1`$ is defined everywhere. This in not true for $`p_2`$, because of dilemmas. However, next player diagonal probabilities are unaffected by dilemmas, so that $`p_2`$ is defined on the diagonal, $`𝐝=𝐝^{}`$; it is even defined for all $`(𝐝^{},𝐝)`$, if and only if $`𝐝^{}`$ is a brelan, since brelans do not produce dilemma, as noticed at end of section 4.1.
Here are a few properties of the functions $`p_i`$:
$`p_i(0,𝐝^{},0,𝐝)`$ $`=`$ $`\delta _{𝐝^{},𝐝},`$ (46)
$`p_i(J,\mathrm{𝟎},j,\mathrm{𝟎})`$ $`=`$ $`\delta _{j,0},`$ (47)
$`p_i(1,𝐝^{},1,𝐝)`$ $`=`$ $`p(𝐝),`$ (48)
$`p_i(J,𝐝^{},j,𝐝)`$ $`=`$ $`0,j<J,𝐝^{}𝐝,`$
$`p_1(J,𝐝,j,𝐝)`$ $`=`$ $`p_i(j,𝐝,j,𝐝),j<J,`$
$`p_2(J,𝐝^{},j,𝐝)`$ $`=`$ $`0,j<J,`$ (49)
$`{\displaystyle \underset{𝐝B^+(|𝐝^{}|)}{}}{\displaystyle \underset{j=0}{\overset{J}{}}}p_1(J,𝐝^{},j,𝐝)`$ $`=`$ $`1.`$
$`{\displaystyle \underset{𝐝B^+(|𝐝^{}|)}{}}p_2(j,D𝐞_f,j,𝐝)`$ $`=`$ $`1.`$
Let the cumulative diagonal probability be
$$s_i(J,𝐝)\underset{j=1}{\overset{J}{}}p_i(J,𝐝,j,𝐝).$$
(50)
Because of the next players’ round duration condition
$$J>1,s_1(J,𝐝)>s_2(J,𝐝)=p_2(J,𝐝,J,𝐝)>p_1(J,𝐝,J,𝐝).$$
To reduce the $`p_i`$ computational domain, I use invariance with respect to face permutations (for unloaded dice). Firstly, diagonal probabilities depend on only one combination. As in (1), two combinations are equivalent, modulo the functions $`𝐝p_i(J,𝐝,j,𝐝)`$, for all $`(i,J,j)`$, if and only if their occupation numbers (Lagrangian components) form the same combination, e. g. $`441655`$. With $`(D,F)=(3,6)`$, the quotient set contains three classes: that of brelans ($`111`$), that of sequences ($`123`$)<sup>3</sup><sup>3</sup>3I do not mean that all combination in the class of sequences is a sequence., that of pairs ($`112`$). Secondly, non-diagonal probabilities depend on a couple of combinations. Two *couples* of combinations are equivalent, modulo the functions $`(𝐝^{},𝐝)p_i(J,𝐝^{},j,𝐝)`$, for all $`(i,J,j)`$, if and only if their couples of occupation numbers form the same combination, e. g. $`(421,442)(321,211)`$. A face permutation transforms a next player’s optimal one-goal strategy into another, possibly different if the goal is not a brelan.
Taking into account (46) and face permutation invariance, the result probabilities (45) are computed, for $`(D,F,J)=(3,6,3)`$, by applying $`\mathrm{strat}\mathrm{moy}`$ on optimal $`𝐝^{}`$-goal strategies and $`𝐝`$-Kronecker utilities. As a consequence of self-similarity (23), the probabilities after the initial time ($`J_1<J`$), are obtained as intermediary results in the computation of a priori probabilities ($`J_1=J`$). The results are presented in the probability charts 6, 7, 8, 9, 10 (appendix C), which do not fill more than a few pages thanks to the extensive use of face permutation invariance and other properties (46…). There are 31 classes of three-die combination couples (including the three diagonal classes).
## 5 Goal identification programming
I will propose heuristic policies, based on the global maximization of expected utility, with respect to the subset of optimal one-goal strategies, for which result probabilities were obtained in the last section.
### 5.1 Motivation: bounded complexity
The $`\mathrm{max}\mathrm{moy}`$ backward induction algorithm is optimal, short, but the number of numerical operations per time step, already large for $`(D,F,J)=(3,6,3)`$, is unbounded as a function of the maximum round duration $`J`$. Information theory teaches that a message will be transmitted faster by a specialized code. $`\mathrm{max}\mathrm{moy}`$ backward induction is slow, for the general reasons that it is unspecialized (and optimal).
To speed-up policy, possibly at the expense of brevity and optimality, specialization is necessary. For example, consider the game of Nim \[7, § 1.3\]: besides $`\mathrm{max}\mathrm{moy}`$ backward induction, a stratagem is found, based on congruence, producing optimal strategies, with a bounded number of operations per time step. The ratchet ($`D<F`$) would be a stratagem of 421, if only the goal were known.
I propose to identify the goal, rigorously, by considering not only the utility, but also the result probabilities (45), obtained in section 4. I will obtain goal identification heuristic policies, that may be considered as quasi-Markovian, from the remark following (45). Roughly, they transfer the complexity of $`\mathrm{max}\mathrm{moy}`$ backward induction to the result probabilities, with the advantage that the latter can be compiled once for all (and the inconvenience that they must be remembered).
For a one-goal utility, goal identification is simple. For a constant utility, as well: any goal is optimal. Difficulties are thus with utilities somewhere between peaked and flat, “*fuzzy*”, e. g. with peaks of about the same height, playing the roles of attractors, that one has to choose between.<sup>4</sup><sup>4</sup>4Like Buridan’s donkey, starving from hesitating between bushels of oats and water.
### 5.2 Reduced horizon
I consider a time and state dependent utility, as in (33, 36), in a round of maximum duration $`J`$. $`𝐝_0`$ is the state at time $`j_0J`$. I define the “evaluation function”,
$$u_{j_0}^0(𝐝_0)\underset{𝐝^{}B^+(D_0)}{\mathrm{max}}\underset{j=0}{\overset{Jj_0}{}}p_i(Jj_0,𝐝^{},j,𝐝^{})u_{j_0+j}(𝐝_\mathrm{𝟎}+𝐝^{}),$$
(51)
where $`j`$ is the renormalized time and $`D_0=D|𝐝_\mathrm{𝟎}|`$. Evaluation functions are often used in stage game (chess, othello, checkers…) programming, but they are usually defined empirically, unlike (51), which is probabilistic.
To take into account serendipity – that a result other than the goal may be not so bad, after all – (51) is improved:
$$u_{j_0}^1(𝐝_0)\underset{𝐝^{}B^+(D_0)}{\mathrm{max}}\underset{j=0}{\overset{Jj_0}{}}\underset{𝐝B^+(D_0)}{}p_1(Jj_0,𝐝^{},j,𝐝)u_{j_0+j}(𝐝_\mathrm{𝟎}+𝐝),$$
(52)
which cannot be used for next players, because of dilemmas. For all $`𝐝_0B^+(D)`$, considering (47), the evaluation functions (51, 52) simply return the utility.
$`\mathrm{max}\mathrm{moy}`$ backward induction is particularly slow, because it needs to completely analyze the round even before its first decision. Hence the idea that short-sighted policies may be faster. At time $`j_0`$, a horizon $`h`$ may be chosen, such that $`j_1=j_0+hJ`$, and the round is virtually terminated at $`j_1`$, taking for ersatz utility the evaluation function $`u_{j_1}^s`$ given by (51) or (52), depending on the serendipity bit $`s\{0,1\}`$. With $`j_1=J1`$, considering (48), (52) reproduces the deepest $`\mathrm{max}\mathrm{moy}`$ iteration, so that an optimal strategy is generated.
I will further examine $`h=0,1`$. With $`h=0`$, the goal is found by maximizing $`u_{j_0}^s`$, independently of the first event. With $`h=1`$, as there is no interest in thinking before casting the dice, the decision $`𝐝_1`$ is rather taken after the first event $`𝐝_{1/2}`$, according to
$$\underset{𝐝_1}{\mathrm{max}}u_{j_0+1}^s(𝐝_1).$$
(53)
In case of many optimal decisions in (53), the corresponding states, written as increasing Lagrangian lists, e. g. 124, are discriminated according to the lexicographic order (only pure strategies are generated). In case of many optimal goals in (51) or (52), we need not discriminate between them, and the policy reproduces the human character of *duplicity*. Dilemma implies duplicity, but the converse is false.
### 5.3 Dynamic programming and goal revision
The strategy may be revised to take into account new events, which is an instance of dynamic programming or belief revision , realizing a feedback of fate on strategy. By self-similarity of the round, a policy may be applied at any time, with suitable parameter renormalization. Self-similar revision based on the $`\mathrm{max}\mathrm{moy}`$ backward induction policy would just confirm the optimal strategy, computed a priori: it is therefore useless. Only fallible policies are worth revising.
A heuristic policy of horizon $`h1`$ forecasts, at any given time, only the next $`h`$ decisions. Thus, it must be run with the period at least $`h`$. The revised serendipitous goal identification policy of horizon $`h`$ is optimal in its last $`h`$ decisions. The goal identification policy with $`h=0`$ does not require revision and is very simple (short and fast). It may be the only rational policy, simple enough for unaided human players in normal game conditions.
### 5.4 Policy benchmark and interpretation
For $`(D,F,J)=(3,6,3)`$, I consider a few increasingly fuzzy stationary utilities:
1. $`u=\delta _{123}`$, a one-goal utility,
2. $`u=\delta _{123}+\delta _{224}+\delta _{345}`$, a three-goal utility,
3. $`u=t`$, the transfer function defined by table 5 in appendix,
4. the sum of faces.
These utilities are unrealistic, in the sense that they may not be possible within a real 421 set (see section 2.3). I consider the policies: $`\mathrm{max}\mathrm{moy}`$ backward induction, and the four goal identification policies $`(h,s)\{0,1\}^2`$; the $`h=0`$ policies are without revision.
From the final utility, on the leaves of the fate tree, every policy yields a pure strategy, and its initial expected utility $`u_0`$, on the trunk, is obtained by solving the Kolmogorov equation exactly, with the $`\mathrm{strat}\mathrm{moy}`$ algorithm. Optimality is defined as the ratio of the expected utility, over the first player optimal expected utility $`u_{0r}`$. The numerical results (approximated by decimal numbers) are copied from into the tables 1, 2, 3, 4.
Table 1 confirms that for a one-goal utility, all goal identification policies are by definition optimal. Compared to the first player, next players are handicapped, but less with a fuzzier utility. The numerical results show a positive contribution of serendipity, much greater with the greater horizon and revision. The contribution of horizon and revision is positive with serendipity. Without serendipity, the contribution of horizon and revision is positive for peaked utilities, negative for fuzzy utilities (3, 4).
I take advantage of this effect to give a (less fuzzy) definition of fuzziness: a utility is fuzzy if and only if introducing horizon and revision without serendipity contributes negatively to its expectation. Thus, I have constructed fuzzy utilities, for which introducing horizon and revision decreases the expected utility, even though it is more complex. The response of expected utility with respect to complexity is non-increasing (this effect compares, in electricity, with a negative resistance).
## 6 Conclusions
The mathematics of fate in 421 leave as the only unsolved difficulty “bifurcations”, that maximizing the expected utility does not always determine a unique decision, as in next players’ dilemmas. Here is a toy example: a game with three players, P, A, B. If P says white, then A gives one euro to B; if P says black, then B gives one euro to A. P earns nothing anyway; A, B take no decision. Maximizing P’s expected utility does not determine its decision. Introducing a mixed strategy amounts to consider P as a random generator, with unknown probabilities. A classic postulate of statistical theories is to maximize the entropy or missing information , which here sets the probabilities of either outcome to $`1/2`$. Are the postulates of mixed strategy and maximum entropy so easily acceptable? We cannot exclude hidden determinism or bias in P. For example, P may always choose the first answer in the lexicographic order (black), or P may have a secret agreement with A to share his gain.
Bounded complexity, similar to *bounded rationality* in , motivates heuristic policies, where characters close to actual human behavior are found, in agreement with . These characters are fate, dilemma, goal identification and revision, restricted horizon, serendipity, duplicity and panic. When the policy belongs to an organization, we are in management. When an individual decides for himself, we are in psychology. For example, the same mathematical effect is behind counterproductive management or panic.
Goal identification consumes a bounded number of operations per time step, whatever the round duration, because it does not resolve all decisions in the fate tree, but only those which are compatible with the present state, and before the horizon. Goal identification is not generally optimal, as opposed to a common assertion in business courses. Only $`\mathrm{max}\mathrm{moy}`$ backward induction, which has no goal, just like random playing, is generally optimal. In the round, the ratchet stratagem allows the immediate translation from goal to decision. I used probability theory as the logic of goal identification, à la Jaynes . Complexity hides in the result probabilities, to be compiled before playing, as a kind of training.
Depending on complexity resources and utility, policies may be variably applicable or good. Starting from a given policy, one may increase optimality, by modifying its characters or the utility: this is the task of human resources management, when the policy is that of an individual taking decisions for a company, a manager. The short-sighted manager ($`h=0`$) gets hardly any help from serendipity. The unserendipitous manager should avoid fuzzy utilities and favor precise goal assignments. I obtain examples of counterproductive management: with a fuzzy utility and no serendipity, goal revision dramatically reduces the optimality. The role of serendipity was pointed out, on purely qualitative ground, by N. Wiener, about scientific and technical invention . The present work also pertains to Wiener’s cybernetics.
Rationality can be further reduced. At the extreme, the fool manager can be trusted only for a flat utility. The study of irrational or illogical but actual behavior is the task of sophistry . It may be quite useful in game practice, to produce best responses.
I thank researchers of the GREQAM in Marseilles, for fruitful discussions.
## Appendix A The (tentative) rules of 421
I define the game, from oral tradition and . The hardware consists of three dice and eleven tokens, initially in a pot. There are two or more players who can always see the positions of dice and tokens.
In the first part of the game, the charge, players get tokens from the pot. In the second part of the game, the discharge, players get tokens from each other. A player wins when he gets no token during the charge (many players may thus win), or when he first gets rid of his tokens during the discharge.
The charge or discharge is a sequence of sets. In every set, each player at his turn plays a round against the dice, while the others wait. The active player casts the dice up to three times; after every cast, he can put aside any number of dice, thus accumulating a combination. Next players must cast dice as many times as the first player.<sup>5</sup><sup>5</sup>5The order of players in the set matters, but I could not find definite rules for its determination. End-of-round accumulated combinations, obtained by all players in the set, are ranked in the hierarchic order
$$\begin{array}{c}421111611666511555411444311\hfill \\ \hfill 333211222654543432321665\mathrm{}221,\end{array}$$
(54)
where $``$ means ‘higher than’. The combinations, implicit in (54), are ordered as the numbers formed by their faces in a decreasing sequence: e. g. $`655654`$. The dominant combination 421 and the dominated combination 221, known as “nénette”, differ only by one die. $`fff`$ is the $`f`$-brelan, $`f11`$ is the $`f`$-pair ($`f1`$), $`654,543,432,321`$ are the sequences.
At end of set, the last<sup>6</sup><sup>6</sup>6The adjective ‘last’ is my own suggestion for automatic tie-breaking. player who has got the lowest combination gets the number of tokens determined by table 5, e. g. if the highest combination is 411, then the last player with the lowest combination (whatever it is) gets 4 tokens. During the charge, tokens are taken from the pot, if possible. When the pot is empty, the discharge begins, and tokens are now taken from the player who has got the highest combination.
## Appendix B A Galton-Watson process in the 421 round
Taking the genealogic point of view, each die is considered as an individual, dying after being cast, either without a child, in case of accumulation, or with a single child (itself indeed). The child number being lower than one, the number of live dice $`D_j`$ (section 3.3) decreases in time. Moreover, the population becomes extinct after $`J`$ casts (or sooner).
A Galton-Watson process is obtained when the offspring of each individual is independent of others’. With an optimal $`𝐝^{}`$-goal strategy, the dice dying without children have their faces in $`𝐝^{}`$, but the converse is not true. For example, with $`𝐝^{}=221,𝐝_1=211,J>1`$, the two dice 11 have correlated offspring: one has a child if and only if the other has none. Dice have independent offspring if and only if $`𝐝^{}`$ is a brelan and the player is first.
I apply the Galton-Watson theory \[12, §6.2\] to obtain the probability law of $`D_j`$, for an optimal $`𝐝^{}`$-goal strategy, where $`𝐝^{}=D𝐞_F`$ is the $`F`$-brelan. Dice are indexed by $`d=1\mathrm{}D_j`$ . Let $`Z_d\{0,1\}`$ the number of children of the die indexed by $`d`$.
$$D_j=\underset{d=1}{\overset{D_{j1}}{}}Z_d.$$
(55)
The $`Z_d`$ are random variables, with the same law $`q_i𝒫(Z_d=i)`$, of generating function
$$g(z)z^{Z_d}=q_0+q_1z,q_0=\frac{1}{F},q_1=1q_0.$$
The $`Z_d`$ are always independent if and only if $`𝐝^{}`$ is a brelan and the player is first. When this is true, from (55), the generating function of $`D_j`$, conditioned by $`D_{j1}`$, is
$$z^{D_j}|D_{j1}=d=g(z)^d.$$
The generating function of $`D_j`$ is thus determined by
$$g_0(z)=z^D,g_J(z)=1,$$
$$\begin{array}{c}(j,1j<J),g_j(z)z^{D_j}=\underset{d=0}{\overset{D}{}}z^{D_j}|D_{j1}=d𝒫(D_{j1}=d)\hfill \\ \hfill =\underset{d=0}{\overset{D}{}}𝒫(D_{j1}=d)g(z)^d=g_{j1}g(z).\end{array}$$
(56)
By induction,
$$g_j=g_0g^j.$$
The composition powers of the affine function $`g`$ are
$$g^j(z)=1q_1^j+q_1^jz.$$
Therefore
$`g_j(z)`$ $`=`$ $`(1q_1^j+q_1^jz)^D,`$
$`𝒫(D_j=d)`$ $`=`$ $`\left(\begin{array}{c}D\\ d\end{array}\right)(1q_1^j)^{Dd}q_1^{jd}.`$ (57)
$`D_j`$ follows a binomial law, directly obtained by considering that a die dies when accumulated, or stays alive, with the probability $`q_1`$ per time step, independently of others: a Bernoulli process is obtained, with the law (57). The interest of considering a Galton-Watson process is in the analogy with branching processes .
## Appendix C Realization with mathematica
The present article is supported by , an open source software and data base in the mathematica language , which, like LISP, is interpreted and allows functional and recursive treatments on arbitrary expressions, equivalent to trees. The mathematica frontend allows literate programming in the form of notebooks, gathering live code, outputs and comments, within a tree structure, that can be unfolded at will.
Combination manipulation differs slightly from list manipulation (since order does not matter in combinations) or ensemble manipulation (since repetitions are allowed in combinations). A tool box is developed. The numerical parameters $`(D,F,J)`$ are arbitrary, which realizes a scalable model, invaluable for development. Fate trees are created recursively. All fates converging to the same state at the same time are merged by indexing, so that the size grows only linearly with the depth $`J`$ and remains easily manageable for $`(D,F,J)=(3,6,3)`$. In exchange, the computing time is increased and the history is lost, which allows to treat only time and state dependent utilities (as required in the 421 set).
Starting from the leaves of the fate tree, where utility is grafted, optimal strategies and expected utilities are build recursively, according to the $`\mathrm{max}\mathrm{moy}`$ algorithm. A utility-strategy tree is finally obtained, from which the strategy can be extracted, then piped into the $`\mathrm{strat}\mathrm{moy}`$ algorithm, a variation on $`\mathrm{max}\mathrm{moy}`$, solving the Kolmogorov equation.
$`\mathrm{max}\mathrm{moy}`$ produces the expected optimal one-goal strategies, Bernoulli or ratchet, depending on $`D<F`$, and dilemmas. The result probabilities are computed, saved, and many properties are checked systematically. Some result probabilities are checked by Monte Carlo simulations, with success. The charts 6, 7, 8, 9, 10 are generated automatically. There is very little room for errors, and if there are any, they are traceable.
The goal identification heuristic policies are realized. Their wrong decisions are pointed out. They are exactly evaluated with $`\mathrm{strat}\mathrm{moy}`$, which is very slow, since it requires the computation of every heuristic decision in the fate tree, according to an algorithm actually longer and slower, for one decision, than the simple maximization in $`\mathrm{max}\mathrm{moy}`$. Obtaining the truth about heuristic policies is a lengthy task.
Probability charts player’s guide
p1, p2 mean first or next players. In every box of a diagonal probability chart stands a column of the probabilities, ordered from top to bottom by growing delay, to obtain the goal written at head of line.
In every box of a non-diagonal probability chart, stand two columns: at left, from top to bottom, the goal and the result; at right, the probabilities, ordered from top to bottom by growing delay, to obtain the result, with the goal in mind (and taking optimal decisions as determined by the ratchet). Moreover, for easy access, the couples (goal, result) are represented in a square array, where heads of lines and columns are the respective representatives of goal and result, modulo face permutations (section 4.2). The three-die representative 3X3 array is spread onto the three charts 8, 9, 10, one for each goal class.
Here is an example for using non-diagonal charts. Let the goal be $`641`$ and the result $`652`$. The representatives of $`641`$ and $`655`$ are, separately and respectively, $`123`$ and $`112`$. (Representatives are chosen so as to minimize the sum of their faces.) The representative of the *couple* $`(641,655)`$ is $`(123,144)`$. $`123`$ takes us to chart 10 (the third line of the representative square array), whence $`112`$ takes us to the second column, $`(123,144)`$ to the third row, where finally are the probabilities to obtain, with the goal 641, the result 652, after one, two or three casts.
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# Untitled Document
HOW HIGH THE TEMPERATURE OF A LIQUID
BE RAISED WITHOUT BOILING?
Mala Das, B. K. Chatterjee, B. Roy and S. C. Roy
Department of Physics, Bose Institute
93 / 1 A. P. C. Road, Calcutta 700009, India
Lines to be deleted from the manuscript for blindfold refereeing
## Abstract
How high a temperature of a liquid be raised beyond its boiling point without vaporizing (known as the limit of superheat) is an interesting subject of investigation. A new method of finding the limit of superheat of liquids is presented here. The superheated liquids are taken in the form of drops suspended in a dust free gel. The temperature of the superheated liquid is increased very slowly from room temperature to the temperature at which the liquid nucleates to boiling. The nucleation is detected acoustically by a sensitive piezo-electric transducer, coupled to a multi channel scaler (MCS) and the nucleation rate is observed as a function of time. The limit of superheat measured by the present method supersedes other measurements and theoretical predictions in reaching the temperature closest to the critical temperature of the liquids.
PACS numbers: 64.60.-i, 64, 64.70.Fx, 64.60.Qb, 64.60.My
Any fluid that exists in the liquid form above its boiling temperature is said to be superheated. These liquids are in a metastable state in the thermodynamic sense and can be nucleated to form vapor by homogeneous nucleation or by the presence of heterogeneous nucleation sites such as gas pockets, vapor bubbles, solid impurities etc. or by the radiation interactions caused by charged particles, neutrons etc. Vapor embryos of different sizes, which are responsible for homogeneous nucleation, are produced at thermal equilibrium in the superheated liquid. The superheated state owes its existence to an energy barrier which causes the vapor embryo to collapse, rather than lead to nucleation, if it is less than a critical size.
A liquid can not be superheated up to the critical temperature, there is a limit to the maximum attainable temperature for any given liquid without boiling. This limit is called the ‘limit of superheat of the liquid’ (T<sub>sl</sub>), where the height of the energy barrier which maintains the superheated state is of the order of kT and this temperature is a characteristic of any liquid. In addition to its importance in basic science, the knowledge of T<sub>sl</sub> is important in a number of industrial operations where a hot, nonvolatile liquid comes in contact with a cold volatile liquid. If the temperature of the hot liquid reaches to the limit of superheat of the cold liquid, explosive boiling would result. This explosive boiling is a potential hazard in damaging equipment and injure personnel in the vicinity of the blast . The study of T<sub>sl</sub> has another importance since the discovery of bubble chamber by Glaser and superheated drop detector. The operation of this detector depends on the degree of superheat of the liquid, more the liquid is superheated more sensitive is the detector to lower energy radiations. The minimum energy detectable by such detector is therefore limited by the limit of superheat of the detecting liquid. The limit of superheat of liquids can be estimated from the theory and can be measured experimentally. Theoretical calculations are performed either from the pure thermodynamic considerations or using the statistical mechanics. Very good and comprehensive reviews on homogeneous nucleation of lquid and on the limit of superheat are available in literature. One has to note that theoretical calculations are performed for ’pure’ homogeneous nucleation where the chance of heterogeneous nucleation arising out of various interfaces with different surface energies e.g. gas-liquid, liquid-liquid, solid-gas etc. is completely excluded.
Experimental results reported so far are far below the critical temperature of the liquids. One of the reasons being that observing ’pure’ homogeneous nucleation experimentally, without any chance of heterogeneous nucleation is difficult to achieve. Hence, the goal is to reduce the chance of heterogeneous nucleation as far as possible and to use an improved method of quantitative detection of nucleation to see how close one can reach experimentally to the predicted limit of superheat. The present experiment is designed to achieve this goal. Superheated sample used in this investigation is a homogeneous suspension of superheated drops of three liquids (R-12 : C$`Cl_2F_2`$, R-114 : $`C_2Cl_2F_4`$ and R-22 : CHCl$`F_2`$) in a dust free, visco-elastic, degassed gel medium. Suspending the superheated liquid in another liquid (gel) reduces the chance of heterogeneous nucleation. Nucleation is detected acoustically by a piezo-electric transducer and the pulses thus received are digitized and recorded as a function of time by a multichannel scaler. This improved method of determining T<sub>sl</sub> supersedes all other measured values in reaching closest to the critical temperatures. Reviews on previous experimental techniques of measuring the limit of superheat of liquid have been described in detail by Avedisian. As has been found from this literature, all previous experiments except one rely on the qualitative observation of the nucleation visually and therefore the present measurement constitutes the first quantitative measurement of T<sub>sl</sub> using digital electronics.
The limit of superheat can be estimated either from the thermodynamic stability theory or from the analysis of the dynamics of formation of the critical sized vapor embryos (statistical mechanical theory). The superheated state of a liquid is a metastable state and the limit of this metastable state is represented on the P-V diagram by the spinodal curves. For a pure liquid, the spinodal curve or the thermodynamic limit of superheat is defined by states for which
$$\left(\frac{dP}{dV}\right)_T=0$$
(1)
Temperley calculated the value of maximum superheat temperature using van der Waals’ equation of state. The maximum limit of superheat of a given liquid can be expressed as
$$t_m=\frac{27T_c}{32}$$
(2)
where t<sub>m</sub> is the limit of superheat of the liquid . For mathematical simplicity this has been calculated by considering the ambient pressure to be zero. At atmospheric pressure i.e. at P=1, t<sub>m</sub> will be slightly greater than the corresponding value at P=0. Other equations of state such as modified Bertholet equation and Redlich-Kwong equation have also been used to calculate the limit of superheat. As has been observed by Blander and Katz, experimental values of thermodynamic limit clearly exceeded the Van der Waals limit at least for five liquids.
For most of the organic liquids the thermodynamic limit of superheat can be represented empirically by
$$T_{sl}=T_c[0.11(P/P_c)+0.89]$$
(3)
where T<sub>c</sub> is the critical temperature, P<sub>c</sub> is the critical pressure and P is the ambient pressure.
Another method of estimating T<sub>sl</sub> using statistical mechanics involves considerations of the rate processes of nucleation to form vapor embryos in a superheated liquid. This method does not yield an absolute value of T<sub>sl</sub> but it allows one to estimate the probable rate of formation of critical-sized vapor embryos in a superheated liquid at a given temperature. If the rate is very low within the time scale of the experiment, one considers no nucleation would occur, while if the rate is very high, then one assumes that T<sub>sl</sub> has been reached. The rate of homogeneous nucleation (J) as given approximately by the Volumer-Doring formula is given by
$$J=Nf\mathrm{exp}\left(\frac{B}{kT}\right)$$
(4)
where J is the expected rate of formation of critical sized vapor embryos per unit volume, f is a frequency factor which in general is of the order of 10<sup>11</sup> sec<sup>-1</sup>, N is the number density of molecules in the superheated liquid and B is the minimum amount of energy needed to form a vapor bubble of critical size as given by Gibbs from reversible thermodynamics is
$$B=16\pi \gamma {}_{}{}^{3}\left(T\right)/3(p_vp_o)^2$$
(5)
where $`\gamma (T)`$ is the liquid-vapor interfacial tension, $`P_v`$ is vapor pressure of the superheated liquid and $`P_o`$ is the ambient pressure. It is to be noted in this connection that which value of J is proper to calculate T<sub>sl</sub> is not defined and therefore one has to make some ’judicious choice’ of a rate which would correspond to T<sub>sl</sub>. A J value of 10<sup>6</sup> nucleation/cm<sup>3</sup>.sec is often used to define the limit of superheat temperature.
It is to be noted that all the above discussions are related with the classical theory of nucleation. Effect of other factors like diffusion, viscosity and other hydrodynamical constraints are discussed by Blander and Katz . As has been pointed out by them, contributions arising out of these effects in calculating T<sub>sl</sub> of pure liquids are not very significant.
The experiment is carried out with superheated liquids of R12 (b.p. -29.79 $`{}_{}{}^{o}C`$), R114 (b.p. 3.6$`{}_{}{}^{o}C`$) and R22 (b.p. -40.5$`{}_{}{}^{o}C`$). The superheated drops are suspended in dust free, de-gassed visco-elastic gel. The gel is a mixture of ’aquasonic’ gel available commercially and glycerine. A glass vial containing the superheated drops homogeneously suspended in gel is placed on the top of a thin layer of degassed gel taken in a beaker. The gel in the beaker improves the acoustic coupling between the superheated drops in the vial and the transducer. The beaker is placed on a piezoelectric transducer with a coupling gel. Some pure gel is placed on the top of the sample and a thermometer was inserted in the pure gel so as to avoid any contact with the superheated liquid sample, thus reducing the chance of heterogeneous nucleation from the liquid-glass interface. The nucleation in superheated drops is detected by the transducer, the output of the transducer is digitized and recorded by a multi channel scaler. The vial was wrapped with a heating coil covering the gel and sample. The temperature of the sample is increased slowly from room temperature and the count rate (dN/dt) is recorded in MCS. As nucleation proceeds, the number of superheated drops are depleted and hence the nucleation rate is normalized with respect to the number of drops present. What we expect ideally is $`(\frac{1}{N})`$$`\frac{dN}{dt}`$ is zero till the temperature reaches the limit of superheat where there will be a sudden increase in $`(\frac{1}{N})`$$`\frac{dN}{dt}`$ (entire liquid nucleates) and will be no nucleation beyond this temperature. Considering the experimental uncertainty, one may observe the similar behavior as presented in Fig. 1. The comparison of observed limit of superheat with other experimental results is presented in the table below. The reduced limit of superheat defined as $`T_{sl}`$/$`T_c`$ (taken in <sup>o</sup>K) for these liquids is also presented in the table along with theoretically predicted values and other experimental results.
| | | observed T<sub>sl</sub> $`{}_{}{}^{o}C`$ | | Reduced limit of superheat \[T<sub>sl</sub>(K)/T<sub>c</sub>(K)\] | | | | |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| Liquid | T<sub>c</sub> | | | Predicted values from | | | Experiment | |
| | $`{}_{}{}^{o}K`$ | Present | Others | (eqn.2) | (eqn.3) | (eqn.4) | Present | Others. |
| R12 | 384.5 | 80.0 | 72.0 | 0.84 | 0.89 | 0.90 | 0.92 | 0.90 |
| R114 | 418.7 | 120.5 | 102.0 | 0.84 | 0.89 | 0.91 | 0.94 | 0.90 |
| R22 | 369.0 | 57.5 | 54.0 | 0.84 | 0.89 | 0.89 | 0.89 | 0.89 |
| | | | | | | | | 0.89 |
As could be seen from the table, the measured limit of superheat exceeded the predicted limit of superheat and other experimental values. It is to be noted in this connection that all theretical predictions are approximate as discussed before. Therefore the present experimental measurements indicate the need of improved calculation of limit of superheat. That the Van der Waals’ limit is exceeded was reported before by Blander and Katz. This table also gives an useful insight about the nucleation process. As can be seen from the table, for liquids with lower boiling points it is harder to reach closer to the critical temperature. This is quite expected as the chances of heterogeneous nucleation increases in case of liquids with lower boiling points. Whether complete elimination of heterogeneous nucleation in experimental measurement is possible or not is an open question. No other measurement have been able to reach so close to the critical temperature. It is to be noted in this connection that the limit of superheat of only 14 liquids out of 56 liquids studied by Blander and Katz hardly exceeded 90% of the critical temperature.
Therefore, by reducing the chances of heterogeneous nucleation by suspending the superheated sample in another ‘pure’ liquid and using precise electronic measurement we have been able to reach closer to the critical temperature hitherto unattainable. Inspite of the fact that theoretical calculations are performed for ‘pure’ homogeneous nucleation, they fall below the experimental values indicate the inadequacy of the present method of calculation discussed here and warrants improved calculations.
\[Present address : Chemistry Dept. Univ. of Utah. Salt Lake City. USA.\]
\[Author for correspondence : Fax no. 91 33 350-6790,
email : scroy@bosemain.boseinst.ernet.in\]
Lines to be deleted from the manuscript for blindfold refereeing
1. R. C. Reid Advances in Chemical Engineering 12, 199 (1983).
2. D. A. Glaser Phys. Rev 8, 665 (1952).
3. R. E. Apfel US Patent 4,143,274 (1979).
4. R.E. Apfel, S. C. Roy and Y.C. Lo Phys. Rev. A. 31, 3194 (1985).
5. Blander M. and Katz J. L. AIChE 21, 833 (1975).
6. Avedisian C. T. J. Phys. Chem. Data 14, 695, (1985).
7. Basu D. K. and Sinha D. B. Ind. J. Phys. 42, 198, (1968).
8. Apfel R. E. and Roy S. C. Rev. Sci. Inst. 54, 1397, (1983).
9. Temperley H. N. V. Proc. Phys. Soc. 59, 199, (1947).
10. Gibbs J. W. Translations of the Connecticut Academy III, p.108 (1875).
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# First Results from the X–ray and Optical1footnote 11footnote 1Based on observations performed at the European Southern Observatory, Paranal, Chile Survey of the Chandra Deep Field South
## 1 INTRODUCTION
In the rocket flight which discovered X-ray stars (Giacconi et al 1962), the presence of a diffused X-ray background (XRB) was also observed. Reviews of the findings up to the middle 70’s are given in Gursky and Schwartz (1977), Fabian and Barcons (1992) and references therein. The measurement by Uhuru of the high degree of isotropy led to the conclusion that the XRB had to be extragalactic. The logN–logS relation for sources at high galactic latitude ($`b>20`$) showed a slope of 1.5 which would account for the entire XRB already at fluxes larger than 10<sup>-14</sup> erg cm<sup>-2</sup> s<sup>-1</sup> (Matilsky et al 1973). The required number of discrete sources had to be large ($`N>10^6`$ sr<sup>-1</sup>) to be consistent with the value of the fluctuations of $`<3`$% over 2 degrees (Schwartz et al 1976).
The spectrum of the XRB had been measured over the range 3 to 100 keV and could be fitted with two power laws of index $`\mathrm{\Gamma }`$ = 1.4 for $`E<25`$ keV and $`\mathrm{\Gamma }=2.4`$ for $`E>25`$ keV (Gursky and Schwartz 1977). However, the results of HEAO–1 gave an excellent fit over the $`330`$ keV range of the XRB spectrum to a thermal bremstrahlung from a hot plasma at 40 keV, leading the authors (Marshall et al 1980) to suggest that the flux could be due to a hot intergalactic medium as had been considered by Field and Perrenod (1977). The finding that the spectrum of AGNs had been measured to be steeper than that of the XRB (Mushotzky 1984) further strengthened the view that the XRB could not be due to the sum of known sources.
On the other hand, the Einstein survey showed that about 25% of the soft XRB (1–3 keV) was resolved into discrete sources at a flux of order $`3\times 10^{14}`$ erg s<sup>-1</sup> cm<sup>-2</sup> (Giacconi 1980, Griffiths et al. 1983). A large fraction of these sources were AGNs. A reasonable extrapolation of the X-ray properties and optical counts of extragalactic sources led to the conclusion that it was unlikely that discrete sources contributed less than 50% of the XRB (e.g., Schmidt & Green 1986; Setti 1985). If so, and their spectrum continued to be steep (energy index $`\alpha 0.71.2`$), then one could show that a diffuse bremsstrahlung hypothesis for the residual was untenable (Giacconi & Zamorani 1987). The COBE result entirely ruled out the possibility that the XRB could be due to thermal bremsstrahlung from a hot plasma (Mather et al. 1990).
The ROSAT survey established that 70–80% of the XRB is resolved into discrete sources in the 1-2 keV range at a flux level of $`1\times 10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup> (Hasinger et al. 1998). At higher energies (2–10 keV) the ASCA and BeppoSAX satellites have achieved detection limits of $`35\times 10^{14}`$ erg s<sup>-1</sup> cm<sup>-2</sup> and have resolved $`25\%`$ of the 2-10 keV XRB into discrete sources (Ogasaka 1998; Cagnoni et al. 1998; Della Ceca et al. 1999a; Giommi, Fiore & Perri 1998). The population of sources appears to be composed of obscured and unobscured AGNs. The majority of hard sources appear to have counterparts in the soft X–ray band and are believed to be obscured AGNs. In this case, the soft emission may be due to electron scattered radiation. In addition to the scattered radiation, there are in principle several possibilities why obscured AGN could show up in the soft band: a very high redshift, so that the absorption cutoff is shifted to the soft energy band; a starburst component; photoionized gas emission (as observed by XMM in NGC 1068, Paerels 2000).
The very high sensitivity, broad energy range and high angular resolution of the Chandra observatory (Weisskopf, O’dell & van Speybroeck 1996) will permit the final resolution of the origin of the XRB at least in the 0.5 to 10 keV range. However, the most important aspect of the deep surveys (exposure time $`1`$ Msec) to be conducted with Chandra will be the study of the individual classes of objects rather than their integrated properties. Based on the limiting sensitivity, exposure map and source surface density found in the first Chandra data, we expect a limiting flux of $`f_{[0.52keV]}=2\times 10^{17}`$erg s<sup>-1</sup> cm<sup>-2</sup> in the center of the field for 1.5 Msec of exposure. Such a depth will permit the study of: (1) the formation and evolution of AGNs at large redshifts ($`z5`$); (2) clusters of galaxies at redshifts greater than $`2`$ and (3) study of star–forming regions in galaxies at moderate redshifts ($`z1`$).
We selected a field for our observations which has the following properties: (1) very low Galactic neutral hydrogen column ($`N_H8\times 10^{19}\mathrm{cm}^2`$) – a Southern sky equivalent to the Lockman Hole; (2) no stars brighter than $`m_v=14`$ and (3) well suited to observations with 8 meter class telescopes such as the VLT and Gemini. Our observations are centered at RA = 3:32:28.0 DEC = -27:48:30 (J2000). The results in this paper pertain to the initial 130 ksec of observations out of the currently planned 500 ksec (another exposure of 500 ksec has been approved, for a total of 1 Msec, for the end of the year 2000). The plan of the paper is as follows. In §2 we describe the data reduction. In §3 we will present the results on the X–ray data including spectral properties of the sources, number counts and correlation function. In §4 we will describe the optical data. In §5 we will discuss our results. Finally, we will summarize our conclusions in §6.
## 2 X–RAY DATA: OBSERVATIONS AND DATA REDUCTION
The Chandra Deep Field South (CDFS) data are obtained from the combination of two exposures taken when the Advanced CCD Imaging Spectrometer–Imaging (ACIS-I) detector was at a temperature of -110 K. The first observation (Obs ID 1431\_0) was taken on 1999-10-15/16, in the very faint mode, for a total 34 ks exposure. The second observation (Obs ID 1431\_1) was taken on 1999-11-23/24, in the faint mode, for a total 100 ks exposure.
The data were reduced and analyzed using the CIAO software (release V1.3, see http://asc.harvard.edu/cda). The Level 1 data were processed using the latest calibration files and the latest aspect solution. We included the new quantum efficiency uniformity files, that correct the effective area for loss due to charge transfer inefficiency at a temperature of -110 K. This correction compensates for the loss of events far from the readout, especially at high energies, and it is particularly relevant when fitting the total spectrum in the broadest energy range.
The data were filtered to include only the standard event grades 0,2,3,4 and 6. All hot pixels and columns were removed as were also the columns close to the border of each node, since the grade filtering is not efficient in these columns. We looked for flickering pixels, defined as pixels with more than 4 events contiguous in time. The removal of columns and pixels reduces slightly the effective area of the detector and this effect has been included when computing the total exposure maps. Time intervals with background rates larger than $`3`$ sigma over the quiescent value ($`0.31`$ counts s<sup>-1</sup> per chip in the 0.3–10 keV band) were removed. This procedure gave 25 ksec of effective exposure out of the first observation, and 93 ksec out of the second, for a total of 118 ks. The two observations had different roll angles, so that the exposure in the combined image of the field ranges between 25 and 118 ksec over a total coverage of 0.096 deg<sup>2</sup>. Since the two observations have the same nominal aimpoint, the exposure time is 118 ksec over the majority of the field of view. The chip S–2 was active during both observations. We did not include data from this chip in the following analysis, since the point spread function (PSF) is very broadened ($`15`$ arcsec) and it adds a very small effective area especially at low fluxes.
We selected a soft band from 0.5–2 keV and a hard band from 2–7 keV in which to detect sources. The hard band was cut at 7 keV since above this energy the effective area of Chandra is decreasing, while the instrumental background is rising, giving a very inefficient detection of sky photons. A wavelet detection algorithm run on the 7–10 keV image, gives no sources over our detection threshold. Note also that the counts expected on the ACIS-I detector in the 7–10 keV band for the hardest sources, are always no more than 3% of the total hard photons. Thus, including the 7–10 keV band would always decrease the signal for the hard sources. However, we will always quote the fluxes in the canonical 2–10 keV band, as extrapolated from the flux measured from the 2–7 keV count rate, in order to have a direct comparison with the previous results.
Figure 1 shows a sky map of the field. We used a wavelet detection algorithm (Rosati et al. 1995) to find X–ray sources. In order to match the PSF variation as a function of the off-axis angle, the wavelet analysis was carried out at five scales, $`a_i=(\sqrt{2})^i`$ pixels, for $`i=0,1,2,3,4`$ and pixel size 0.984 arcsec. Simulations have shown that simple aperture photometry is very accurate across the ACIS–I detector and was therefore preferred to time–consuming wavelet photometry (see Rosati et al. 1995). The area of extraction of each source is defined as a circle of radius $`R_s=2.4\times FWHM`$ (with a minimum of 5 pixels $`(5\mathrm{}`$). The FWHM (in arcsec) is modeled as a function of the off-axis angle $`\theta `$ (in arcmin) as $`FWHM(arcsec)=_{i=0,3}a_i\theta ^i`$, with $`a_i=\{0.678,0.0405,0.0535\}`$ (see the Observatory Guide, http://asc.harvard.edu/udocs/docs/docs.html). The background was calculated locally for each source in an annulus with outer radius of $`R_s+8^{\prime \prime }`$ and inner radius of $`R_s+2^{\prime \prime }`$, after masking out other sources. With this choice, the average number of counts in the background regions is $``$ 11–22 in the soft band and $``$ 17–34 in the hard band. Thus the local estimate of the background has a poissonian fluctuation always less than 30%. We define the $`S/N`$ ratio as $`S/\sqrt{S+2B}`$ where S are the net counts in the extraction region of radius $`R_s`$, and B are the background counts found in the annulus defined above and rescaled to the extraction region. Since the lowest $`S/N`$ ratio of our final catalog is 2 across the field, the flux threshold across the field in the soft band, including the effect of vignetting and of the point spread function, is $`2\times 10^{16}`$ erg s<sup>-1</sup> cm<sup>-2</sup> within the central 6 arcmin, and grows to $`5\times 10^{16}`$ erg s<sup>-1</sup> cm<sup>-2</sup> at 10 arcmin off–axis. In the hard band, the flux limit is $`2\times 10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup> within the central 6 arcmin, and $`4\times 10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup> 10 arcmin off–axis. In the following we describe the analysis of point–like sources only, leaving the analysis of diffuse sources to a subsequent paper.
We performed extensive simulations to determine the accuracy of our aperture photometry and our completeness limit. The aperure photometry with the aforementioned parameters is accurate within 3%. From the catalog of candidate sources selected by the wavelet detector, we remove all sources with a signal-to-noise ratio $`S/N<2`$, corresponding to 7 net counts in the center. With these criteria, our detections have a less than 5% probability of including a fake source on the total ACIS-I field. This is a robust limit since the background is very low and we are always signal limited. More importantly, the simulations provide a test for the sky coverage model, which is defined as the area of the sky where a point–like source of a given flux can be detected by the wavelet algorithm and has a $`S/N>2`$ in the extraction region of radius $`R_S`$. The effective solid angle of the Chandra observations is equal to the geometrical solid angle (0.096 deg<sup>2</sup>) only at fluxes $`>2\times 10^{15},2\times 10^{14}`$ erg s<sup>-1</sup> cm<sup>-2</sup> in the soft and hard bands respectively, whereas it drops at 50% of these values at fluxes $`<3.2\times 10^{16},2.3\times 10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup>. By comparing the input LogN–LogS with the output LogN–LogS in the simulations we verified that our model for the sky coverage is accurate within 5%. A further check comes from the preliminary analysis of a total exposure of 300 ksec of the CDFS. We found the presence of 5% of spurious sources, to be added to six sources associated with flickering pixels that were not removed from the exposure.
We used two separate conversion factors to derive the energy flux from the observed count rate for the soft and hard bands. The conversion factors were $`(4.8\pm 0.3)\times 10^{12}`$ erg s<sup>-1</sup> cm<sup>-2</sup> per count s<sup>-1</sup> in the soft band, and $`(2.7\pm 0.3)\times 10^{11}`$ erg s<sup>-1</sup> cm<sup>-2</sup> per count $`\mathrm{s}^1`$ in the hard band assuming an absorbing column of $`8\times 10^{19}`$ cm<sup>-2</sup> (Galactic value) and a photon index, $`\mathrm{\Gamma }=1.7`$, consistent with the average spectrum of the bright sample. The uncertainties in the conversion factors reflect the range of possible values for the photon index, $`\mathrm{\Gamma }=1.42.0`$. As suggested by the spectral analysis of the stacked spectra, these values are representative of our sample (see §3.1). The conversion factors were computed at the aimpoint (which is the same for both exposures). Before being applied to the net count rate of a given sources, the conversion factors are corrected for vignetting. The correction is given by the ratio of the value of the exposure map at the aimpoint to the value of the exposure map at the source position (averaged over the extraction region). This is done separately for the soft and the hard band, using the exposure maps computed for energies of 1.5 keV (soft) and 4.5 keV (hard). Sources with more than 100 total counts in the two bands were analyzed also with XSPEC (Arnaud 1996), fitting an absorbed power law to the data binned with a minimum of 20 photons per bin. The soft and hard fluxes of the fitted spectra are always within 10% of the values obtained using the conversion factors quoted above.
## 3 X–RAY DATA: RESULTS
### 3.1 X–ray Spectra and Hardness Ratio
The majority of the sources are faint. Thus we measure the average stacked spectrum of both the faint sources and bright sources. To this purpose, we use only the longest exposure (93 ksec), in order to have a uniform exposure for the stacked spectra. From our sample of 159 sources, 5 were excluded because they were present only in the field of view of the first exposure. We divided the remaining sample into two groups of 29 and 125 with a dividing count rate of $`1.0\times 10^3`$ cts/sec in the total band 0.5–7 keV. The background spectrum was obtained using the event file of the total field of the same exposure, after the removal of the detected sources. The background is scaled by the ratio of the total exposure maps of the sources and of the background. Such a procedure guarantees a correct background subtraction despite the non-uniformities of the instrumental background across ACIS-I. The ancillary response matrix for the effective area is obtained from the counts–weighted average of the matrices of the single sources. The response matrix is assumed to be the one computed in the aimpoint. We recall that we use the quantum efficiency uniformity file appropriate for a temperature of -110 K. This file corrects the quantum efficiency with respect to the pre–flight values; especially, it takes into account the correction due to the charge transfer inefficiency.
We used XSPEC to compute the slope of a power law spectrum with $`N_H`$ absorption at low energy, in the energy range 0.5–10 keV and in the 2–10 keV band only. The value of $`N_H`$ is fixed at the galactic value when the fit is done using the hard energy range, while, if the total energy range is used, $`N_H`$ is left free to vary. We exclude bins below $`0.5`$ keV because the calibration is still uncertain below this energy.
For the total sample, including all the 154 sources on the largest exposure, we performed a power-law fit over the energy range 0.5–10 keV, and obtained a photon index $`\mathrm{\Gamma }`$ of $`1.53\pm 0.07`$ and a column density $`N_H=(6.0\pm 3.0)\times 10^{20}`$ cm<sup>-2</sup>. Errors refer to the 90% confidence level. Then we perform the fit using only the 2–10 keV energy range, with Galactic $`N_H=8\times 10^{19}`$ cm<sup>-2</sup>. We obtained $`\mathrm{\Gamma }=1.61\pm 0.11`$. The contribution of the total sample to the XRB is $`(4.0\pm 0.3)\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup> deg<sup>-2</sup> in the soft band and $`(1.07\pm 0.15)\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> deg<sup>-2</sup> in the hard band.
For the bright sample we obtained a photon index $`\mathrm{\Gamma }`$ of $`1.71\pm 0.07`$ and $`N_H=(7.0\pm 2.0)\times 10^{20}`$ cm<sup>-2</sup> from the fit in the 0.5–10 keV energy range. Using the 2–10 keV energy range with galactic $`N_H`$ we obtained $`\mathrm{\Gamma }`$ of $`1.70\pm 0.12`$. The contribution of the bright sample to the XRB is $`(6.3\pm 0.9)\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup> deg<sup>-2</sup> in the hard band.
We note that our bright sample shows a spectrum that is softer than that measured at about the same fluxes by ASCA in the energy range 2–10 keV (Della Ceca et al. 1999a find $`\mathrm{\Gamma }=1.36\pm 0.16`$ at $`S6\times 10^{14}`$ cgs, Ueda et al. 1999 find $`\mathrm{\Gamma }=1.6`$). This may be due to cosmic variance, since we poorly sample the bright end of the number counts.
For the faint sample we obtained a photon index $`\mathrm{\Gamma }`$ of $`1.26\pm 0.10`$ and $`N_H=(8.0\pm 4.0)\times 10^{20}`$ cm<sup>-2</sup> from the fit using the 0.5–10 keV energy range, and $`\mathrm{\Gamma }`$ of $`1.35\pm 0.20`$ with galactic $`N_H`$ using the 2–10 keV energy range. The contribution of the faint sample in the hard band is $`(4.3\pm 1.0)\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup> deg<sup>-2</sup>. In this case the spectral shape of the faint sample is consistent with the average spectrum of the XRB, which is around 1.4. We conclude that the average spectrum of the detected sources is approaching the average shape of the hard background (see also Mushotzky et al. 2000).
The results of the spectral fits are shown in Table 1 and 2. The new quantum efficiency file provided with the CXC software was used which includes the effect of a reduced quantum efficiency at higher energies caused by the radiation damage.
To test the accuracy of our background subtraction, we perform the same spectral analysis in two ways. First, we build a synthetic spectrum using the software of M. Markevitch (see http://hea-www.harvard.edu/ maxim/axaf/acisbg/) and we extracted the background from the same regions used for the extraction of the source spectra. In this way we eliminate the uncertainties due to the use of different extraction regions for the sources and the background. The results are the same within a few percent. Second, we build a background file summing all the background events found in the annuli around each source, and scaling the resulting file by the ratio of the total area af the extraction radius to the total area of the annuli. In this case too, the spectral fits change by a few percent.
A more detailed view of the spectral properties as a function of the fluxes comes from the hardness ratio $`HR=(HS)/(H+S)`$ where H and S are the net counts in the hard (2–7 keV) and the soft band (0.5–2 keV), respectively. The distribution of the hardness ratios as a function of the count rate in the soft band is shown in Figure 2. There are 15 sources which are observed only in the hard band plotted at $`HR=1`$ (corresponding to 9% of the total combined sample) and 68 which are observed only in the soft band, and are plotted at $`HR=1`$. When the hardness ratio of the stacked spectrum of the only hard and only soft sources are computed, we find respectively $`HR=0.52\pm 0.07`$ and $`HR=0.65\pm 0.05`$. These hardness ratios are plotted as big asterisks. It is clear that faint sources are harder than bright ones, as already shown by the fits of the faint and bright stacked spectra.
We note that the stacked spectrum of the 15 sources detected only in the hard band, is consistent with the typical hardness ratio of the sources detected in both bands at the count rate of $`10^4`$ cts/sec (see asterisks in the upper left of figure 2). In fact, we detected soft emission at 3 sigma levels in the stacked spectrum of these 15 sources. They have been missed in the soft band only because their emission were just below the detection threshold, while they would have been detected in a longer exposure. With this result, we have no evidence for sources in which the soft emission lower is lower than a 5–10 % of the energy emitted in the total (0.5–10 keV) band. This is expected in synthesis models of the XRB, where a soft component is considered in obscured AGNs.
### 3.2 LogN–LogS and Total Flux from Discrete Sources
We compute the number counts in the soft and hard bands. We show in Figure 3 a comparison of the soft 0.5–2 keV band to the ROSAT results. We find the Chandra results in excellent agreement with ROSAT in the region of overlap $`S_{min}>10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup>. The Chandra data extend the results to $`2\times 10^{16}`$ erg s<sup>-1</sup> cm<sup>-2</sup>. In Figure 4 we show the LogN–LogS distribution for sources in the hard band, extending to a flux limit of $`2\times 10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup>, which represents a substantial improvement over previous missions.
We performed a maximum likelihood fit with a power law in the range $`2\times 10^{16}<S<3\times 10^{14}\mathrm{cgs}`$ for the soft band and $`2\times 10^{15}<S<5\times 10^{14}\mathrm{cgs}`$ in the hard band. The resulting best fit to the soft LogN–LogS is:
$$N(>S)=370\left(\frac{S}{2\times 10^{15}}\right)^{0.85\pm 0.15}\mathrm{𝑠𝑜𝑢𝑟𝑐𝑒𝑠}deg^2$$
(1)
and to the hard LogN–LogS:
$$N(>S)=1200\left(\frac{S}{2\times 10^{15}}\right)^{1.0\pm 0.20}\mathrm{𝑠𝑜𝑢𝑟𝑐𝑒𝑠}deg^2$$
(2)
The errors on the slope correspond to 2 sigma. The normalization of the soft counts is consistent within 1 sigma with the estimate of Mushotzky et al. (2000). On the other hand, we find that the previous Chandra results for the hard counts reported by Mushotzky et al. (2000), shown as crosses, are higher by 40% (see §5). From the maximum likelihood analysis we find that the normalization of the hard counts is, in fact, more than three sigma lower than the value found by Mushotzky et al. (2000). This is partially due to a different $`\mathrm{\Gamma }`$ used in deriving the converson factors. If we use an average $`\mathrm{\Gamma }=1.4`$, which is appropriate for the faint end of the number counts, the discrepance in normalization between our hard counts and that of Mushotzky et al. (2000) is reduced to 3 sigma. Our hard counts are consistent with those derived from other deep ACIS-I pointings, namely the Lynx Field (Stern et al. 2000, in preparation) and the HDF–N field (Garmire 2000, priv. communication), and in a completely independent fashion, from the first XMM deep surveys of the Lockman Hole (Hasinger et al. 2000, A&A, in press). In order to further investigate the difference with Mushotzky et al. results, a quantitative understanding of the clustering of X–ray sources on scales $`5`$ arcmin would be needed, together with a knowledge of the sky coverage of that survey.
The integrated contribution of all sources within the flux range $`10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup> to $`2\times 10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup> in the 2-10 keV band is $`(1.05\pm 0.2)\times 10^{11}`$ erg s<sup>-1</sup> cm<sup>-2</sup> deg<sup>-2</sup>. When this value is added to the contribution for fluxes $`>10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup> from ASCA (Della Ceca et al. 1999b), we have a total contribution of $`(1.3\pm 0.2)\times 10^{11}`$ erg s<sup>-1</sup> cm<sup>-2</sup> deg<sup>-2</sup>, which is close to the value $`1.6\times 10^{11}`$ erg s<sup>-1</sup> cm<sup>-2</sup> deg<sup>-2</sup> from UHURU and HEAO-1 (Marshall et al. 1980). More recent values of the 2–10 keV integrated flux from the BeppoSAX and ASCA surveys (e.g., Vecchi et al. 1999; Ishisaki et al. 1999 and Gendreau et al. 1995) appear higher by 20–40%. The integrated contribution of all sources in the field plus the bright sources seen by ASCA (data kindly provided by R. Della Ceca) is shown in Figure 5, together with the best estimates of the total background. The inclusion of the ASCA data at bright fluxes minimizes the effect of cosmic variance. We conclude that, given the uncertainty on the value of the total background, a fraction between $`20`$% and $`40`$% is still unresolved.
As for the total contribution to the soft X–ray background, we refer to the 1–2 keV band, following Hasinger et al. (1993, 1998) and Mushotzky et al. (2000). We find a contribution of $`5\times 10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup> deg<sup>-2</sup> from discrete sources for fluxes lower than $`10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup>, corresponding to $`11`$% of the unresolved flux ($`4.38\times 10^{12}`$ erg s<sup>-1</sup> cm<sup>-2</sup> deg<sup>-2</sup>). If this value is summed to the contribution at higher fluxes (see Hasinger et al. 1998), we end up with a total contribution of $`3.5\times 10^{12}`$ erg s<sup>-1</sup> cm<sup>-2</sup> deg<sup>-2</sup> for fluxes larger than $`2\times 10^{16}`$ erg s<sup>-1</sup> cm<sup>-2</sup>, corresponding to 80% of the unresolved value.
### 3.3 Angular Correlation Function
We have performed an analysis of the clustering of the X–ray sources in the CDFS using the sources which were identified in both the soft and hard bands. Since the limited statistics of our sample is expected to provide a weak clustering signal, we consider the average value of the two-point angular correlation function, $`\overline{\omega }(\vartheta )`$, defined as the excess of source pairs at separation $`\vartheta `$ with respect to a random distribution (e.g., Peebles 1980). The random control sample was generated by distributing 50,000 points at random within an area of the sky with the same geometry as our field and a distribution modulated by the exposure map. Statistical errors in our value of $`\omega (\vartheta )`$ are estimated using the standard bootstrap resampling technique. The errors estimated with this technique are those due to statistical noise and do not include cosmic variance, which can only be assessed by repeating the analysis over several independent fields.
In Figure 6 we report the results for both the whole source population and for the subsample of 103 sources with fluxes $`>10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup>. For reference, we also plot the power–law shape, $`\overline{\omega }(\vartheta )=(\vartheta _c/\vartheta )^{1\gamma }`$, where $`\gamma =2`$ is the slope of the spatial correlation function and $`\vartheta _c=10`$ arcsec. Our analysis shows a statistically significant correlation signal out to scales of about 100 arcsec, while the noise dominates at larger separations. There is a marginal indication for the brighter sources to be more clustered at smaller separations. If confirmed, this would indicate that close pairs tend to be formed preferentially by bright objects, as is apparent even by visual inspection of the CDFS.
Vikhlinin & Forman (1995) analysed the correlation function for the the angular distribution of sources identified within ROSAT PSPC deep pointings. Quite remarkably, and despite the different nature of the sample they used, their correlation function turns out to be consistent with our Chandra data. Although their much larger statistics and the wider PSPC field–of–view allowed them to extend their analysis out to scales $`\vartheta 800`$ arcsec, the exquisite angular resolution achievable with Chandra provides here a detection of a correlation signal down to scales $`10`$ arcsec, about four times smaller than those sampled by Vikhlinin & Forman (1995).
We defer to a forthcoming paper a detailed correlation analysis for the distribution of the Chandra discrete sources and its implication for models aimed at explaining the nature of X–ray emission from AGNs.
## 4 OPTICAL AND NEAR–IR DATA
### 4.1 Imaging Survey Data
For identification of sources we have carried out an imaging survey in 3 bands, VRI with the FORS–1 camera at the ANTU telescope (UT–1 of VLT) (Programme Id 64.O–0621). We obtained exposures between 4000 to 7000 seconds in 4 adjacent fields of 6.8 x 6.8 arcmin<sup>2</sup> covering a large fraction of the CDFS 16x16 arcmin<sup>2</sup> field. Some shallower R-band fields of about 1200 seconds each in were obtained with the FORS imager in 4 additional positions in order to cover the entire area in the sky swept by different orientations of the Chandra field of view. For a few objects we were able to obtain spectroscopic information using the FORS–1 multislit capabilities in March 2000, just before the CDFSbecame inaccessible.
Our identification process was complemented with data from the ESO Imaging Survey (Rengelink et al. 1998) in the J and K bands obtained with SOFI at the NTT and in the UV bands obtained at the ESO NTT + SUSI–2. We also used a $`30\times 30`$ arcmin image in B band obtained during the commissioning of the WFI camera at the ESO/MPI 2.2 meter telescope. The near-IR imaging covers, at present, the central $`9.4\times 9.3`$ arcmin of the CDFS.
In the source identification process, we have found a positional offset of order of 0.3 arcsec and then found the correspondence of some of the brightest X–ray sources to likely candidates. We find an positional offset (-0.2, +1.4) arcsec between the optical and X–ray data. Using this correction we find X–ray/optical position deviations highly concentrated in a radius of rms 0.67 arcsec. A 2 arcesc correlation radius has been used for further analysis.
### 4.2 Optical counterparts of X–ray sources
Spatial coincidence allows identification of a large fraction of the sources. Approximately 10% of the X–ray sources with more than 10 detected counts have no immediately identifiable optical counterpart at $`m_R<26`$. This could be due to a flux ratio $`S_X/S_{opt}`$ larger than average but still close to the values observed for most of the sources (see Figure 7). We find approximately $`1/3`$ of the sources are extended, which we associate with galaxies, and $`2/3`$ of the sources appear are point like. If we compare the diagram of X–ray flux vs R magnitude for ROSAT (Hasinger et al. 1999) and the Chandra Deep Field South, we find that a large part of the Chandra sources appear to be consistent with an extension to fainter fluxes of the ROSAT sources, with an average $`S_X/S_{opt}=1`$. However a significant fraction of the sample ($`10\%`$) appear to have a very low $`S_X/S_{opt}1/10`$ and is detected only in the soft band at fluxes $`<10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup>. Most of these objects are identified as galaxies. Another subsample of objects that have $`S_X/S_{opt}10`$, appear below fluxes $`3\times 10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup> (soft band). For these sources (about 20) with no optical counterpart at $`R26`$, it could be that the corresponding objects are fainter, thereby increasing the $`S_X/S_{opt}`$ to more than $`10`$, or that they correspond to clusters.
We also plot R-K vs R magnitudes for the Chandra sources (see Figure 8). The color magnitude plot of the sources in the Chandra Deep Field South essentially overlaps the same plot with the ROSAT sources and extends to fainter magnitudes. The faintest Chandra sources appear to have very red colors, R–K $`5`$, with two of them at R–K $`6`$, possibly indicating high redshift obscured AGNs.
### 4.3 Spectroscopic Information on Selected Sources
We have obtained spectroscopic data for a dozen optical counterparts of the X–ray sources using the multislit capability of FORS–1. The results are reported in Table 3, where we give the redshift, the counts in the soft and hard bands, the hardness ratio, the 0.5–2 keV X–ray flux in erg s<sup>-1</sup> cm<sup>-2</sup>, the 0.5–2 keV intrinsic luminosity in units of $`10^{43}`$ erg s<sup>-1</sup> cm<sup>-2</sup>, and the preliminary optical classification based on the width of the emission lines for Seyferts and QSOs. The objects described in Table 3 are not meant to be a well defined sample. They are just a random sample of bright CDFS sources for which we obtained the first spectra. The optical follow up started October 26 2000. At present we identify an elliptical galaxy where we observe an early type continuum with strong H and K absorption and weak, or no emission lines.
In total we identified 4 QSOs, 5 Seyfert 2 galaxies, 1 elliptical galaxy and 2 galaxies in an interacting pair. While we make no claim for completeness, it is interesting to note that the type of objects we observe are not extremely different from those we have observed in the ROSAT surveys. This can also be shown by plotting the computed values of $`L_X`$ against z for the Chandra identifications (Figure 9). The many objects shown in the diagram are from previous surveys, as specified by different symbols. We note that the luminosity of the objects we have classified as QSOs is of order $`17\times 10^{43}`$ erg s<sup>-1</sup>, while the luminosity of those we classify as Seyfert 2 range between $`10^{41}`$ and $`5\times 10^{43}`$ erg s<sup>-1</sup> cm<sup>-2</sup>. The elliptical galaxy has a luminosity of $`7\times 10^{40}`$ erg s<sup>-1</sup> and the nearby interacting pair galaxies have luminosities of 6 and $`9\times 10^{39}`$ erg s<sup>-1</sup>. None of these objects appear to exhibit unusual properties.
## 5 DISCUSSION
From the sources detected in the Chandra Deep Field South we resolve 60–80% of the hard X–ray background. Their X–ray and optical properties are consistent with AGNs being the dominant population.
Absorbed AGNs, which have hard X–ray spectra, are missed in shallow surveys because of absorption dimming. They should be detected efficiently at low fluxes. The average spectrum of sources detected in X–ray surveys is expected to become harder with decreasing flux. This trend, which has already been observed in surveys performed with other satellites, is shown in Fig. 2.
Although the average spectrum of the sources gets harder with decreasing flux, only a small fraction ($`9\%`$) of the sources are detected in the hard band and not in the soft band. The hardness ratio of their stacked spectrum is consistent with the hardness ratios of the hard sources already detected both in the soft and hard bands, showing that they may lie just below the soft detection limit. This might imply that there is not a sizeable population of sources visible only in the hard band. Indeed, even in highly obscured AGNs, soft X–ray emission could be produced by partial covering of the nuclear emission, scattered components, or circumnuclear starbursts associated with AGN. Evidence for soft components in obscured AGNs is commonly observed in good quality X–ray spectra (e.g. Turner et al. 1997). Furthermore, the analysis of the X–ray color–color plot for recent ASCA and BeppoSAX surveys (Della Ceca et al. 1999b, Fiore et al. 2000, Giommi, Perri & Fiore, 2000) is consistent with this scenario, with the soft components being 1–10% of the nuclear emission.
For about 30% of the X–ray sources a host galaxy is resolved in the optical images. Most of the objects clustered on the lower–left corner of the soft $`S_X`$ vs R diagram of Figure 7 are associated with bright galaxies. These galaxies are overluminous in the R band with respect to the other X–ray detected sources. Their X-ray emission generally comes from the center and is soft. Many of them are undetected in the hard band. The upper limits of the hardness ratios indicate that for the four brightest galaxies the average slope is $`\mathrm{\Gamma }1.5`$. For three of these galaxies we have measured redshifts of z = 0.075, 0.075 and 0.215 with $`L_X`$ of $`9\times 10^{39}`$, $`6\times 10^{39}`$ and $`7\times 10^{40}`$ erg s<sup>-1</sup> respectively as given in Table 3. The remainder of the galaxies are optically fainter and presumably more distant with luminosity greater than $`10^{41}`$ erg s<sup>-1</sup>. All our soft galaxies fall below the detection threshold of the deepest ROSAT survey, $`10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup> (Hasinger et al. 1998), and thus Chandra has opened up a new parameter space for the study of moderate redshift galaxies in the X–ray band.
Two sources have been detected in the K but not in the R band (see Figure 8). They have R-K $`6`$, and can be classified as Extremely Red Objects (EROs, Elston et al. 1988). Very red colors would be produced both by a dust reddened starburst or by an old stellar population in high redshift galaxies ($`z1`$, c.f. Crawford et al. 2000). In the first case the X–ray emission could arise from a surrounding cluster. Alternatively these objects could be high redshift obscured AGNs. The last hypothesis seems more likely for two objects with a similar color found by ROSAT in the Lockman Hole, based on their hard X-ray spectrum (Lehmann et al. 2000). However, the two EROs in our field are detected only in the soft band, suggesting that the X–ray emission is likely due to a starburst component or to a cluster. Another interesting possibility is that we are observing the soft scattered component in highly obscured Compton-thick AGNs at high redshift.
The cumulative number counts in the CDFS are fully consistent with those of Mushotzky et al. (2000) in the soft band, while in the hard band we found a surface density lower by $`3040`$% depending on the average spectral slope $`\mathrm{\Gamma }`$ assumed for the X–ray source population. In the soft band the LogN-LogS is within the region allowed by ROSAT fluctuation analysis (Hasinger et al. 1993), while in the hard band the BeppoSAX fluctuation analysis (Perri & Giommi 2000) does not extend to sufficiently low fluxes to allow a full comparison. The source counts are in agreement within the errors with the predictions of AGN synthesis models (Comastri et al. 1995, Gilli et al. 1999). Furthermore, the source density observed at low fluxes seems to favor models where the cosmological evolution of absorbed AGNs is faster than that of unabsorbed ones Gilli, Salvati & Hasinger (2000).
## 6 CONCLUSIONS
In the Chandra Deep Field South we have a similarly large survey area and limiting flux to the Chandra Deep Field North carried out at Penn State. Our results are in general agreement with those quoted by this group who, in their two recent papers (Hornschemeier et al. 2000, Brandt et al. 2000), concentrated on a limited number ($`10`$) of sources where they had optical identifications.
The limiting flux of the Mushotzky et al. (2000) results is similar but the area surveyed is about one third of the two Deep Fields. While we agree on the soft counts, we have a $``$ 40% lower normalization of the hard counts, which implies that an important fraction of the XRB, at least 20%, has yet to be resolved.
Our data confirm that the XRB is due to the summed contribution of individual sources. In the 0.5–2 keV band this result had been foreshadowed by the Einstein measure of an individual source contribution of 25% (Giacconi 1980) and by the ROSAT determination of a 70% contribution (Hasinger et al. 1998). At higher energies (2–10 keV) the contribution to the background from sources was of 25% (Cagnoni et al. 1998; Della Ceca et al. 1999a; Ueda et al. 1999; Giommi, Fiore & Perri 1998). Our results bring the individual source contribution to the XRB in the 2–10 keV range to the level of 60–80%. Since there is still a non-negligible fraction of order 20-40% of the XRB to be resolved, it is important to push the detection limit at least down to $`5\times 10^{16}`$ erg s<sup>-1</sup> cm<sup>-2</sup> in the hard band, or even at lower fluxes if the number counts flatten below $`10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup>.
With our increased number of sources we determined that the hardness ratio increases with fainter limiting flux and the spectrum of the fainter sources matches the background. Much of this was expected from earlier work of Setti and Woltjer (1989), Schmidt and Green (1986), Madau, Ghisellini & Fabian (1993), Comastri et al. (1995), and the ASCA results by Della Ceca et al. (1999a).
In addition, the angular correlation function of the X–ray sources in our field exhibits significant power on scales of order of the survey size. The slope and correlation length is consistent with that observed for galaxies. This indicates the presence of large–scale clustering at high redshift, whose spatial extent requires knowing the redshift distribution of the identified sources, by either spectroscopic or photometric methods. Detailed analyses of the correlation function in Chandra deep pointings will provide important complementary information on the nature of the discrete X-ray sources (e.g., Haiman & Hui 2000, and references therein).
Only 9% of the sources are detected in the hard band and not in the soft band. Most of these sources show some flux in the soft band below our detection threshold. As discussed in §3.1, the average hardness ratio of these faint sources is $`HR=0.52\pm 0.07`$ consistent with that found for the faint end of the sample. We conclude that we find no evidence for sources which are completely obscured in the X–ray soft band.
Our colour-magnitude and flux–magnitude diagrams follow the trends well known from the ROSAT data. A small fraction of the X–ray sources ($`<10\%`$) with more than $`10`$ detected counts are not immediately identified at R $`<`$ 26.
A substantial number of nearby galaxies appear in the Chandra survey at very low $`S_x/S_{opt}`$ and below the detection threshold of ROSAT. The X–ray emission from these galaxies is characterized by a soft spectrum and thus does not contribute significantly to the X–ray background at high energies. However, their study is of great intrinsic interest for the understanding of the properties and evolution of normal and starburst galaxies.
Chandra sensitivity is higher than or equal to that anticipated. Longer exposures to the confusion limit will still be in a signal–limited regime. For a $`2\times 10^6`$ s exposure we should reach $`S_{min}10^{17}`$ erg s<sup>-1</sup> cm<sup>-2</sup> in the soft band, which would allow detection of $`L_X=5\times 10^{40}`$ erg s<sup>-1</sup> cm<sup>-2</sup> at $`z0.5`$. This will enable a statistically significant study of normal and star–forming galaxies to significant look back times. Studies of bright QSOs can reach very large redshifts (5–10) and even observations of low surface brightness features (e.g., clusters) can be extended to $`z3`$.
A deep observation of 385 ksec in the same field is planned during the guaranteed time of one of us with XMM. The larger area of XMM will contribute significantly to the study of the spectra of point sources and the study of diffuse emission from galaxy clusters.
We thank the entire Chandra Team for the high degree of support we have received in carrying out our observing program. In particular, we wish to thank Antonella Fruscione for her constant help in the use of the CXC software. We thank also Maxim Markevitch and Pasquale Mazzotta for discussions. We thank Roberto Della Ceca for discussions and for providing the ASCA number counts. We thank Ingo Lehmann for his contribution to the spectroscopic redshifts. G. Hasinger acknowledges support under DLR grant 50 OR 9908 0. R. Giacconi and C. Norman gratefully acknowledge support under NASA grant NAG-8-1527 and NAG-8-1133.
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# Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 2311footnote 11footnote 1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.
## 1 Introduction
There has long been considerable circumstantial evidence that at least some luminous, low-redshift QSOs are the result of strong interactions or mergers of galaxies Toomre & Toomre, 1972; Gunn, 1979; Stockton, 1982; Hutchings & Neff, 1992; McLeod & Rieke, 1994, e.g.,; see Stockton, 1999 for a review. However, a concrete suggestion for an evolutionary scenario for such objects was lacking until Sanders et al. (1988a) showed that ultraluminous infrared galaxies (ULIGs), virtually all of which are compelling examples of ongoing mergers, had bolometric luminosities and space densities similar to those of QSOs. These similarities suggested the possibility that ULIGs are dust-enshrouded QSOs which, after blowing away the dust, become classical QSOs.
If this hypothesis is correct, one should be able to observe examples of objects that are at intermediate stages of this evolutionary sequence. We are conducting a study of a sample of low-redshift objects that may be in such a transitionary state. These objects are recognized as bona-fide QSOs and are found at an intermediate position in a far infrared (FIR) color-color diagram between the regions occupied by typical QSOs and ULIGs (see Fig. 1). FIR color–color diagrams have been used as tools to detect and discriminate different types of activity in the nuclear and circumnuclear regions of galaxies. Different kinds of objects such as QSO/Seyfert, starbursts, and powerful IR galaxies, occupy fairly well defined regions in the diagram (see, e.g., Neugebauer et al., 1985; de Grijp et al., 1987; Taniguchi et al., 1988; Lípari, 1994). With deep imaging and spectroscopic observations of the host galaxies, we are attempting to construct interaction histories for each of these “transition” objects.
If strong interactions triggered the QSO activity and induced starbursts, one might expect both events to occur roughly simultaneously, since both are plausibly dependent on gas flows to the inner regions. Thus, we are placing these objects on an age sequence by measuring the time elapsed since the last major starburst event. This age sequence along with interaction histories can help us answer the question of whether the intermediate position of these objects is indicative of evolution from the ULIG to the classical QSO population, or whether it simply indicates a range of characteristics in QSOs.
Our sample is drawn from the Neugebauer et al. (1986), Low et al. (1988), and Clements (1996) samples of Infrared Astronomical Satellite<sup>2</sup><sup>2</sup>footnotemark: 2 (IRAS) objects, and it consists of those objects which have: (1) a luminosity above the cutoff defined for quasars by Schmidt & Green (1983), i.e., $`M_\mathrm{B}=23`$ for $`H_0=50`$ km s<sup>-1</sup> Mpc<sup>-1</sup> (or $`M_\mathrm{B}=22.1`$ for $`H_0=75`$ km s<sup>-1</sup> Mpc<sup>-1</sup>), (2) a redshift $`z0.4`$, (3) a declination $`\delta 30\mathrm{°}`$, (4) firm IRAS detections at $`25\mu `$m, $`60\mu `$m, and $`100\mu `$m, and (5) a position in the FIR color–color diagram which is intermediate between the ULIG and QSO loci (Fig. 1). Although Mrk 231 just misses the luminosity threshold given above, its active nucleus is known to suffer heavy extinction (see §5), apart from which it would clearly be a member of the sample. We know of no other objects satisfying the other criteria for which this is true. We have therefore chosen to include it for the present, although it may be approriate to exclude it from some of the analyses of the whole sample, which will be presented in a subsequent paper.
<sup>1</sup><sup>1</sup>footnotetext: The Infrared Astronomical Satellite was developed and operated by the US National Aeronautics and Space Administration (NASA), the Netherlands Agency for Aerospace Programs (NIVR), and the UK Science and Engineering Research Council (SERC).
So far, we have presented results for two of the nine objects in the sample: 3C 48 (Canalizo & Stockton, 2000, hereafter CS2000), an ongoing merger near the peak of starburst activity; and PG 1700+518 Canalizo & Stockton, 1997; Stockton et al., 1998, hereafter CS97 and SCC98; see also Hines et al., 1999, where a tidally disturbed companion with a dominant 85 Myr old post-starburst population may be in the process of merging with the host galaxy. In this paper we present the results for three additional objects: Mrk 1014, IRAS 07598+651, and Mrk 231. We assume $`H_0=75`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and $`q_0=0.5`$ throughout this paper, so that the projected physical length subtended by 1″ is 2.43 kpc for Mrk 1014, 2.26 kpc for IRAS 07598+651, and 0.77 kpc for Mrk 231.
## 2 Observations and Data Reduction
Spectroscopic observations for the three objects were carried out using the Low-Resolution Imaging Spectrometer (LRIS; Oke et al., 1995) on the Keck II telescope. For IRAS 07598+6508, we used a 600 groove mm<sup>-1</sup> grating blazed at 5000 Å yielding a dispersion of 1.28 Å pixel<sup>-1</sup>. For Mrk 231 and Mrk 1014, we used a 300 groove mm<sup>-1</sup> grating blazed at 5000 Å with a dispersion of 2.44 Å pixel<sup>-1</sup>. The slit was 1″ wide, projecting to $``$5 pixels on the Tektronix 2048$`\times `$2048 CCD. We obtained two or three exposures for each slit position, dithering along the slit between exposures. Table 1 shows a complete journal of observations, with specification of the slit positions, and total integration times.
The spectra were reduced with IRAF, using standard reduction procedures. After subtracting bias, dividing by a normalized halogen lamp flat-field frame and removing sky lines, we rectified the two-dimensional spectra and placed them on a wavelength scale using the least-mean-squares fit of cubic spline segments to identified lines in a Hg-Kr-Ne lamp. We calibrated the spectra using spectrophotometric standards from Massey et al. (1988) observed with the slit at the parallactic angle. The distortions in the spatial coordinate were removed with the IRAF apextract routines. For each slit position, we had two or three individual frames; we averaged the spatially corrected spectra using the IRAF task scombine. We then corrected the spectra for Galactic extinction, using the values given by Schlegel et al. (1998).
Since we were aiming to observe the youngest populations in the host galaxies of these objects, we chose the slit positions based on previously obtained color maps of the host galaxies. We obtained imaging data for the three objects with the University of Hawaii 2.2 m telescope as specified in Table 2. Observations with the $`U^{}`$ filter (centered at 3410 Å with a bandpass of 320 Å) sample the spectral energy distribution (SED) of galaxies on the short side of both the 4000 Å break and the Balmer limit. This region is very sensitive to the age of the stellar population, and will be brightest in regions of very recent star formation (ages $`100`$ Myr), as well as regions with scattered QSO continuum. Observations in $`B`$band sample the region redwards of the Balmer limit, where $``$A-type stellar populations are expected to peak, while longer wavelength optical and near infrared images will map the distribution of late-type stellar populations. The $`U^{}B`$ color maps will then highlight the regions of most recent star formation while $`BR^{}`$ will point to the somewhat older ($`200`$ Myr) populations. Thus our slit positions generally cover the regions brightest in these color maps. In general, the spectra from each slit position were subdivided into regions corresponding to some of the main features observed in the optical ground-based images, and one-dimensional spectra were extracted by summing pixels corresponding to the regions of interest.
The spectra of regions close to the QSO nucleus were contaminated by scattered QSO light. The scattered light was removed by subtracting from each region a version of the quasar nuclear spectrum, scaled to match the broad-line flux. In the case of Mrk 1014, none of our slit positions went through the QSO, so we obtained a separate 200 s exposure of the nucleus.
Spectra from those slit positions which actually went through the QSO nucleus suffered from strong light scattering within the spectrograph, particularly in the spectral region around H$`\alpha `$. Therefore we were unable to obtain spectra of regions closer than $`2`$″ from the QSO nucleus for these slit positions.
HST WFPC2 images of the three objects were obtained from the HST data archive. We used three 600 s WFC2 images of Mrk 1014 in the F675W filter, one 400 s and two 600 s PC1 images of IRAS 07598+6508 in the F702W filter, and two 1100 s PC1 images of Mrk 231 in the F439W filter. Most cosmic rays were removed by subtracting a median image from each of the individual frames, then thresholding the difference at a 3$`\sigma `$ level, setting points above this threshold to the median of the difference image. Pixels near the position of the peak of the QSO were excluded from this process. The corrected difference image was then added back to the median image, giving a corrected version of the original image. The few cosmic rays within the relevant region that escaped this process were removed manually with the IRAF task imedit. In the case of Mrk 231, where only two images were available, we first subtracted one from the other, and proceeded as above. The procedure was repeated, interchanging the images. All the corrected images for each object were then averaged.
### 2.1 Modeling Spectra
We use Bruzual & Charlot (1996) isochrone synthesis models to fit the spectra of the host galaxies. We will see in the following sections that the three objects which are the subject of this paper show strong evidence of having undergone some major tidal interaction. As we have described in our previous work (CS97, CS2000), spectra of the host galaxies of such objects show features from both young (e.g., strong Balmer lines) and old (e.g., Mg I b absorption) stellar populations, and can usually be fitted satisfactorily by a two-component model. This model includes an old underlying stellar population, presumably the stellar component present prior to interaction, and a younger instantaneous burst model, presumably produced as a result of the interaction. A population with no age dispersion will be a reasonable approximation of the actual starburst as long as the period during which the star formation rate was greatly enhanced is short compared to the age of the population itself.
We have also noted previously (SCC98) that the age of the superposed starburst is remarkably robust with respect to the different assumptions about the nature of the older stellar component. Thus, we select a reasonable old underlying population (with certain variations as described in each section below) and assume that the same underlying population is present everywhere in the host galaxy. To this population we add instantaneous burst, Scalo (1986) initial mass function, solar metallicity models of various ages. We then perform a $`\chi ^2`$ fit to the data to determine the scaling of each component and the age of the most recent starburst. The errors in the starburst ages that we quote are estimated by noting the youngest and oldest best fits for which $`\chi ^2`$ changes by 15% with respect to the minimum value (see CS2000 for details).
In some cases, stellar absorption features (most often the Balmer lines) were contaminated by emission coming from the extended narrow emission line region around the QSO. In some cases, we subtracted a scaled synthetic spectrum of the recombination lines assuming Case B. However, in calculating $`\chi ^2`$ for the model fitting, we generally excluded those lines that suffered most from contamination. All spectra are displayed as observed (i.e., without line subtraction), unless otherwise specified.
As we discussed in CS2000, because of our limited spatial resolution (generally $`1\mathrm{}\times 1\mathrm{}`$) and projection along the line of sight, we are likely observing the integrated spectrum of several starbursts of different ages, and the age we determine will be somewhat older than the youngest starbursts. Therefore the ages we report should be regarded as upper limits to the most recent episodes of star formation along the line of sight.
In objects with recent starbursts, the effect of reddening by dust is an obvious concern. However, studies of low-redshift AGNs and ULIGs at millimeter and submillimeter wavelengths indicate that dust is generally heavily concentrated within $`1`$ kpc of the nucleus (Andreani et al., 1999; Bryant & Scoville, 1996). In addition, have found (SCC98) that even in a case where the optical—IR spectral index is strongly affected by dust, the stellar ages from spectral features in the rest-frame 3200–5200 Å region remain fairly robust. This relative insensitivity to dust can be attributed to the fact that we are largely dealing with dust that is intermixed with the stars and that we preferentially observe regions with low extinction and, thus, low reddening. Even in regions where dust along the line of sight is significant, the reddening in our observed bandpass is likely to be largely compensated by blue light scattered into our line of sight.
## 3 Mrk 1014
Mrk 1014 ($`z=0.163`$) is a luminous (M$`{}_{B}{}^{}=23.9`$), infrared loud (e.g., Sanders et al., 1988b) radio-quiet QSO which shows a luminous host galaxy. The host galaxy has two large “spiral-like arms” (MacKenty & Stockton 1984; see Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. and Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.). Its spectrum indicates a mixture of old and young stars (MacKenty & Stockton, 1984; Heckman et al., 1984; Hutchings & Crampton, 1990).
Recently, Nolan et al. (2000), in a spectroscopic survey of 26 RLQs, RQQs, and radio galaxies, observed the host galaxy of Mrk 1014 with the Mayall 4 m telescope at Kitt Peak National Observatory and with the 4.2 m William Herschel telescope at La Palma. They modeled the spectra and determined an age of 12 Gyr for the host galaxy of Mrk 1014. Their approach is exactly opposite to ours (see §2.1): they fix the age of a possible young single starburst population to 0.1 Gyr and let that of a second, older single starburst population (what we would call the “underlying population”) vary. Their reason for including the 0.1 Gyr population is to account for “the spectral shape of the blue light” which they attribute to either a recent burst of star formation or contamination of the slit by scattered light from the QSO nucleus. Their results will be discussed further in §3.1.
### 3.1 Stellar Populations
The host galaxy of Mrk 1014 shows stellar absorption features with redshifts remarkably close to those of the QSO broad and narrow emission lines ($`z_{\mathrm{QSO}}=0.1634`$, as measured from our spectrum).
We have obtained and modeled spectra of different regions in the host galaxy of Mrk 1014, and we shall refer to them according to their label in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.. We have chosen a 10 Gyr old population with an exponentially declining star formation rate with an e-folding time of 3 Gyr as an old underlying population. This model fits reasonably well the spectrum of a galaxy 60″ west-southwest of the QSO at the same redshift and only slightly smaller than the “bulge” (see below) of Mrk 1014. (This galaxy was in our slit while obtaining the 200 s exposure of the QSO). The assumption that the host galaxy (or parent galaxies) of Mrk 1014 had a similar star formation history to this galaxy need not be accurate as the precise model we use as the pre-existing population makes little difference in the age determination of the starburst population (SCC98). For comparison, we have tried using a generic elliptical galaxy spectrum, and a model with a longer e-folding time (i.e., 5 Gyr) as underlying populations in the modeling of the spectra in Mrk 1014. We obtain the same starburst ages, though slightly different flux contributions from the old population, regardless of the model used.
Figure Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. includes a $`B^{}R^{}`$ color map of Mrk 1014. This image emphasizes those regions with a steeper blue continuum spectrum, peaking just redwards of the Balmer limit. These regions are concentrated mainly along the north edge of the tail (regions $`a`$ and $`b`$ in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.), in a clump on the east end of the tail ($`c`$), directly east of the nucleus (not covered by our slits), and southwest of the QSO nucleus ($`e`$).
Spectroscopy of these regions confirms that they are indeed the youngest stellar populations we find in the host, with ages ranging from 180 Myr in region $`e`$ to 290 Myr in region $`d`$. Figure 2 shows the spectrum of region $`a`$, with the best $`\chi ^2`$ fit of the model to the data superposed, and the relative contributions of the 200 Myr starburst and the old underlying population. The error in these ages is typically $`\pm 50`$ Myr. Figure Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. shows compact knots at the positions of regions $`a`$ through $`d`$, and Planetary Camera HST images (Surace et al., 1998) show a larger and very bright blue knot at the position of region $`e`$. Even though there are knots at the positions of $`a`$ and $`b`$, it is evident from colors and spectroscopy that there is recent star formation all along the north edge of the tail (i.e., between $`a`$ and $`b`$).
Other regions of the host galaxy appear redder in the $`B^{}R^{}`$ color map in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.. These regions, sampled by $`f`$ and $`g`$, appear to be dominated by an older, $`1`$ Gyr population, and to have very little, if any, contribution from the old underlying population (Fig. 3, top panel). It is not entirely clear whether these are truly intermediate age populations, or if they are simply, as in the case of region $`h`$, dominated by an old underlying population with a very small contribution from a younger population like those found in regions $`a`$ through $`e`$. We attempted to fit models with the latter characteristics to the observed spectra. The resulting fits are reasonable, with minimum values of $`\chi ^2`$ 10% and 25% larger than that obtained for a dominant intermediate-age population for regions $`f`$ and $`g`$ respectively. The bottom panel of Fig. 3 shows the spectrum of region $`f`$ with the best fit to the data of the sum of an old underlying population and a 250 Myr instantaneous burst population. While a reasonable fit, it does show significant discrepancies in the region around the 4000 Å break and the Ca II $`K`$ line.
The potential presence of this intermediate age population suggests that a better fit for the younger populations might be achieved by adding a third component that accounts for this intermediate age component. However the flux of the young starburst population in those regions is so dominant (typically contributing 80% of the total flux at rest-frame 5000 Å), that a third component makes a negligible difference to the fit.
Region $`i`$ shows very strong emission, which almost certainly comes from gas ionized by the QSO rather than from star-forming regions. The equivalent widths of the emission lines here are greater than those of the emission lines elsewhere in the galaxy by at least a factor of 5, and at least twice those of the QSO nucleus for \[O III\]. The emission line ratios indicate a power-law ionizing continuum (Veilleux & Osterbrock, 1987). This region is clearly seen as a large, discrete knot in Fig. 2$`d`$ of Stockton & MacKenty (1987). The gas has an approaching velocity of $`180\pm 50`$ km s<sup>-1</sup> with respect to the stellar absorptions in that region. The underlying spectrum is similar to that of region $`h`$.
Region $`h`$ also shows emission lines with approaching velocities of $`200`$ km s<sup>-1</sup> with respect to the stellar features, but this region shows an additional, weaker emission component clearly visible in \[O II\], \[O III\], and H$`\beta `$, blueshifted by $`1150\pm 50`$ km s<sup>-1</sup> with respect to the stronger component. The stellar population here seems to be dominated by an old population with a small fraction of the flux coming from younger stars.
We have added spectra from two different slit positions at region $`j`$ to improve the signal to noise in this very faint region. This is part of the long extension on the west side clearly seen in the high contrast image in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.. We find a continuum with a red SED and a clear 4000 Å break at a redshift close to, but slightly larger than that of the main body of the galaxy.
Region $`k`$, which is very bright in the $`R^{}`$ image in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555., would appear to be an extension of the east tail. However, its spectrum shows that it is a background galaxy at $`z=0.5843`$. Likewise, region $`l`$ on the west tail is a galaxy at $`z=0.5857`$. A third galaxy SE of the QSO ($`\mathrm{\Delta }\alpha =+11\stackrel{}{\mathrm{.}}9`$, $`\mathrm{\Delta }\delta =40`$″) has a similar redshift as well ($`z=0.5838`$), so there seems to be a group or cluster of galaxies at this redshift.
The 200 s exposure spectrum of the elongated object 9$`\stackrel{}{\mathrm{.}}`$2 west-southwest of the QSO shows narrow emission lines at the same redshift as the QSO superposed on a red continuum. Fig. 2$`d`$ of Stockton & MacKenty (1987) shows that there is strong emission just north of this object coming from extended gas ionized by the QSO and not necessarily associated with the object. The spectrum is too noisy to distinguish stellar features, so it is unclear whether this object is interacting with the system or if it is a chance projection. The near edge-on galaxy $`23`$″ west of the QSO, on the other hand, shows \[O III\] emission confined to the galaxy at $`z=0.162`$ in a University of Hawaii 2.2 m telescope spectrum (Canalizo 2000, unpublished), as suggested by Stockton & MacKenty (1987) from their \[O III\] imaging.
How do our results compare to those of Nolan et al. (2000)? As mentioned above, Nolan et al. use a fixed 0.1 Gyr starburst population. They determine that this population makes up 1.1% of the total luminous mass along the line of sight, and that the rest of the flux is well characterized by a 12 Gyr instantaneous burst. In contrast, we find several regions of recent ($`0.2`$ Gyr) star formation where the star forming mass typically amounts to 12%, and sometimes up to 30%, of the total luminous mass along the line of sight. One of our slit positions is very similar to one used by Nolan et al. (see Fig. 2 in Hughes et al. 2000), but our slit is narrower and slightly closer to the nucleus. As a way to compare, we added the flux along the slit as Nolan et al. seem to have done, including the background galaxy $`k`$, which is also in their slit. Even if we add up all the flux along this slit (subtracting the QSO light, which amounted to 5% of the total flux), we still find that the star forming regions make up 8% of the total luminous mass along the line of sight. Obviously, the difference in ages will lead to a smaller percentage for their choice of parameters. So, we tried fixing our parameters to match theirs (i.e., 12 Gyr + 0.1 Gyr populations), and this yields a 5% by mass for the young population, but the fit is much inferior to the ones discussed in this section. Their slit position may have fortuitously missed the major star forming regions, thus leading to this smaller percentage. We have previously cautioned (CS2000) that different slit positions can lead to different age determinations, and we emphasize the importance of carefully selecting slit positions if one wishes (as we do) to find the major starburst regions. It is also far more difficult to obtain a reliable age for an older population in the presence of a contaminating younger population than the reverse.
### 3.2 Interaction History
Images of Mrk 1014 show a very prominent tail extending to the northeast, reminiscent of the tidal tail of 3C 48 (CS2000). Like the tidal tail in 3C 48, this tail has a number of small clumps (see HST image in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.) of star formation, which are commonly found in merging systems. The bulk of the bright portion of the tail appears to be dominated by an intermediate age population ($``$1 Gyr) with a very small contribution, if any, from an old underlying population; alternatively, it could be dominated by an older population, with some flux scattered from the bright star-forming regions ($`a`$ through $`d`$) or fainter small regions distributed along the host. However, the north edge of the tail appears as a very sharp feature in the short wavelength images, and contains stellar populations which are as young as those of the clumps. As we find no redshift variations along this edge, it would appear that we are observing the tail nearly face-on, and that this is the leading edge where the material has been compressed, thus producing star formation. This, again, is similar to the blue leading edge observed in 3C 48 (Fig. 2$`e`$ in CS2000), presumably sharper in Mrk 1014 because of the lower inclination angle.
Both the ground-based and the HST images of Mrk 1014 (Figs. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. and Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.) show a long, low-surface-brightness extension of the bright tail on the east of the nucleus and arching towards the south (more evident after the removal of the background galaxy, $`k`$) as well a very extended faint secondary tail, rotationally symmetric to the bright (primary) tail. Each tail extends for as much as $`40`$″ or $``$100 kpc (note the inset in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. showing the well-known local interacting system M 51/NGC 5195 at the same scale). Spectra of the secondary tail west of the nucleus indicates that the tail is made up by older stars, and this is consistent with this tail being visible in our $`H`$ (not shown) and $`K^{}`$ images. No bright clumps of star formation are evident along the secondary tail; this absence is not unusual as tidal dwarf formation appears not to be a ubiquitous process in mergers (Hibbard & Yun, 1999). The HST image of the nucleus (see inset in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.) shows a small extension on the south side which follows the direction of the secondary tidal tail. A similar extension is seen in HST NICMOS images by Scoville et al. (2000). There appears to be a “bulge” or enhanced brightness area elongated roughly along the axis connecting the beginning of both tails (i.e., NW–SE), with a half radius of $`4`$ kpc. Assuming a projected velocity of 300 km s<sup>-1</sup> on the plane of the sky, the dynamical age for the tails is $`330`$ Myr. The tails are then older than every major post-starburst knot they contain, consistent with the idea that the latter were formed after the tails were first launched. At the same time, the tails are dynamically younger than the bulk of the stars that form them, though perhaps only slightly so if the stellar populations observed in $`f`$ and $`g`$ are truly $`1`$ Gyr old, instead of being $`10`$ Gyr with a small admixture of younger stars.
Mrk 1014 shows some striking similarities to 3C 48 (CS2000): the morphology of the (primary) tidal tail, the clumps of star formation along the tail as well as its blue leading edge, the relation of the starburst ages to the dynamical age of the tails (both ages $`200`$ Myr younger in 3C 48), and the clumpy extended emission line region (Stockton & MacKenty, 1987) with high velocity ($`>1000`$ km s<sup>-1</sup>) components. As with 3C 48, the data strongly suggest that Mrk 1014 is the result of a merger of two galaxies of comparable size, both of which were disks in this case. The starbursts in the main body of the host of Mrk 1014, however, appear to be less intense and less widespread than those of 3C 48.
If the intermediate age (1 Gyr) population we see in regions $`f`$ and $`g`$ is real, it may be the relic of a starburst ignited at an initial passage of the two interacting galaxies. Indeed, we know of another system (UN J1025$``$0040; Canalizo et al., 2000) where the interacting galaxies have a difference in starburst ages of this order, possibly coincident with their orbital period. If this were the case for Mrk 1014, one might expect most of the star formation at the present to be very strongly concentrated towards the nucleus, as most of the gas would have been driven towards the center starting $`1`$ Gyr ago.
## 4 IRAS 07598+6508
The luminous, $`z=0.148`$, radio-quiet QSO IRAS 07598+6508 was first detected by IRAS, identified as an AGN candidate by de Grijp et al. (1987), and spectroscopically identified as a QSO by Low et al. (1988). The optical spectrum of the QSO is dominated by extremely strong Fe II emission (Lawrence, 1988; Lípari, 1994), and the UV spectrum shows low- and high-ionization broad absorption lines (BAL) extending to blueshifts of 5200 to 22000 km s$`1`$ (Hines & Wills, 1995; Boroson & Meyers, 1992).
Figures Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. and Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. show, respectively, ground-based and HST images of IRAS 07598+6508 (see Boyce, 1996, for a PSF-subtracted version of the HST image). These images show two clumpy regions $`7`$″ south and southeast of the QSO, bright in the PSF subtracted $`U^{}`$ image, but barely visible in $`H`$-band images (not shown). The great number of knots in these regions, presumably OB associations, already argues for recent star formation.
Our slit positions cover these two regions, labeled $`a`$ and $`c`$ in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555., as well as some of the fainter emission surrounding the nucleus. Spectra of these regions are shown in Fig. 4. Regions $`a`$ and $`c`$ are fit by single starburst models of ages $`70\pm 15`$ and $`32\pm 7`$ Myr respectively. The light from these starbursts dominates the spectra as we do not find any significant contribution from an older component. Although single burst models fit the data better than two-component models, there are still some discrepancies in the fit. The observed spectra show an excess with respect to the model in the region between 3900 Å and 4100 Å , apparently because the observed continuum is steeper in this region. To test whether the discrepancy could be an indication of scattered QSO light, we subtracted QSO spectra with several different scalings from the stellar spectra, but we were unable to obtain better fits.
In contrast to $`a`$ and $`c`$, the spectra of the regions closer to the nucleus (labeled $`d`$ and $`e`$ in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.) show an older stellar population (Fig. 4, bottom panel). The SEDs of these spectra are much redder and there is a clear, though not very prominent, 4000 Å break. We detect this old population from $`2\mathrm{}`$ to $`8\mathrm{}`$ north of the QSO and from $`2\mathrm{}`$ to $`4\stackrel{}{\mathrm{.}}5`$ south of the QSO.
We were unsuccessful in attempting to subtract the QSO scattered light closer to the nucleus because of the strong light scattering along the slit in these regions. Therefore we are unable to determine whether the populations closest to the nucleus are as old as those of $`d`$ and $`e`$ or if there might be a younger starburst concentrated around the nucleus.
Regions $`a`$ and $`c`$ have a slightly higher redshift than the fainter emission around the QSO: $`z_a=0.1490\pm 0.0001`$ and $`z_c=0.1488\pm 0.0001`$, compared to $`z_{d,e}=0.1485\pm 0.0002`$ for the regions closer to the nucleus, but the difference is barely significant. Measuring a precise redshift for the QSO is difficult because of the strong Fe II emission contaminating the broad emission lines, and the absence of narrow emission lines. Values in the literature include $`z_{\mathrm{QSO}}=0.1488`$ as measured from CO (Solomon et al., 1997) and $`z_{\mathrm{QSO}}=0.1483`$ as measured from broad Balmer emission lines (Lawrence, 1988); both are within 180 km s<sup>-1</sup> (and within the errors) of the values we measure for the host galaxy.
IRAS 07598+6508 has some limited extended narrow emission as seen in the \[O III\] image in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.. The ionized gas seems to be correlated with the stellar component: the velocity difference between emission and absorption lines in regions where both are present is never larger than 50 km $`s^1`$ (notice the profile of H$`\beta `$ in Fig. 4; emission is visible slightly redshifted with respect to the absorption).
Region $`b`$ shows the strongest emission and appears brightest in the \[O III\] image. The emission line spectrum has an underlying blue continuum with stellar features. The 2-dimensional spectrum shows that this region is broken up into two discrete clumps a little over 1″ each, having some small scale velocity structure. The line flux ratios indicate that the south clump has lower ionization, and it may be an H II region rather than gas ionized by the QSO. A second, fainter component with a velocity gradient from 0 to $`500`$ km $`s^1`$ apparently originating from the north clump extends 1″ towards the south.
About 11″ north of the QSO, we find another emission region (labeled $`g`$) which is visible in the \[O III\] image, but not in any of the broad band images. No stellar continuum is evident from the spectrum, either. This region is also broken into two discrete, somewhat larger clumps, but in this case the clump closer to the nucleus has lower ionization. Both clumps are at a lower redshift than those of region $`b`$, i.e., $`z_{\mathrm{em}}=0.1483\pm 0.0001`$ vs. $`z_{\mathrm{em}}=0.1491\pm 0.0001`$ in $`b`$.
Our deep $`R^{}`$ image (Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.) shows a tidal tail extending from north to the east and arching towards the south of the nucleus for $`22`$″ or $`50`$ kpc. The dynamical age of this feature (assuming a projected velocity of 300 km $`s^1`$) is $`160`$ Myr, again older than the starbursts ages found for this object. We have a very faint spectrum of this tail at region $`f`$. The SED of the spectrum is similar to that of $`d`$ and $`e`$, but the spectrum is too noisy to show stellar features. An \[O II\] $`\lambda 3727`$ emission line is confined to this region (as seen in the 2-d spectrum) at a redshift of $`z_{\mathrm{em}}=0.1487`$, which is probably close to that of the stellar component as in other regions.
This tidal tail is strongly suggestive of a merger. The fact that only one tail is evident may indicate that we are seeing the result of a merger of a spiral with an elliptical galaxy. This configuration, however, could also result from the merger of two spiral galaxies, one of which is counter rotating to the relative orbit. In mergers with such geometry, the gas and stars in the counter-rotating disk are only slightly perturbed and are not pulled into tidal bridges and tails (Toomre & Toomre, 1972; Hibbard & Yun, 1999). It is also possible that a second tail may not be evident because of projection effects.
Regions $`a`$ and $`c`$ may be part of the host galaxy with enhanced surface brightness due to the recent star formation, or they may be remnants of companion galaxies which have strongly interacted with the host galaxy. In either case it is clear that these regions are tidally disturbed and they will likely be completely mixed with the host within a few crossing times.
The two galaxies $`14`$″ south of the QSO are likely to be companion galaxies since they are bright in our \[O III\] image (Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.), but we do not have a spectrum to confirm this. The colors for the SE galaxy indicate that it may have recent star formation, if at the same redshift as IRAS 07598+6508. We find two additional objects with the same redshift as IRAS 07598+6508 that happened to fall in our slits: a faint emission line object $`20`$″ south-southwest of the QSO, and a bright galaxy with an absorption and emission line spectrum $`60`$″ southwest of the QSO.
## 5 Mrk 231
Mrk 231 ($`z=0.042`$), often classified as a Seyfert 1 galaxy, is slightly below the luminosity cutoff for QSOs defined by Schmidt & Green (1983). However, the nucleus is heavily reddened, with an estimated $`A_V2`$ of foreground extinction Boksenberg et al., 1977; Lípari, 1994; see also Goodrich & Miller, 1994, so it would be well above this threshold if it were unobscured. The central plateau of the host galaxy is off center with respect to the nucleus, and the presence of tails to the north and south, as well as a low-surface-brightness extension to the east clearly indicate a recent merger (Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.). Closer to the nucleus, HST imaging shows a large number of stellar associations indicating recent star formation (Surace et al. 1998; see Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.).
As in the case of IRAS 07598+651, the spectrum of the QSO is dominated by strong Fe II emission (Lípari et al., 1994) and shows a peculiar low-ionization BAL system with velocities up to $`7800`$ km s<sup>-1</sup> (Smith et al., 1995; Foster et al., 1995). Previous spectroscopy of the host galaxy indicates the presence of a young stellar population (Hamilton & Keel, 1987).
We have obtained and modeled spectra of several regions in the host galaxy as labeled in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.. In regions $`k`$ and $`i`$, we find no evidence for a significant young starburst population. A 10 Gyr-old population with an exponentially decreasing SFR and an e-folding time of 5 Gyr fits this spectrum quite well, as shown in the bottom panel of Fig. 5. Therefore we use this model as an underlying population to fit the rest of the spectra in the host galaxy.
Unfortunately, because of the poor sensitivity of LRIS shortwards of 4000 Å, and the fact that Mrk 231 has a lower redshift ($`z=0.042`$) than the other two objects, we were unable to obtain the near-UV portion of the spectrum, which is the most helpful region in determining the ages of post-starburst populations. We are left then with the region on the long side of 3800 Å, where more than one combination of models can give similar fits to the spectra.
Indeed, we found degenerate fits in some regions of the host galaxy, particularly in those regions with the youngest populations and those where the absorption lines were heavily contaminated by emission so that we could not use Balmer lines to discriminate between models. In such regions we obtained two $`\chi ^2`$ minima, generally with one corresponding to a model with a small contribution (i.e., small percentage of total mass along the line of sight) from a very young starburst, and the other to a model with a large contribution from a somewhat older starburst. We illustrate the problem in Fig. 6, where we have plotted two models resulting from a 4 Myr and a 42 Myr old populations contributing, respectively, 1% and 12% of the total luminous mass. The models are nearly identical in the spectral region redwards of 3800 Å, and both are good fits to the data (see inset in Fig. 6) with $`\chi ^2`$ values for these models differing only by 10%. The models, however, diverge quickly at shorter wavelengths.
In order to discriminate between the two “best fits” to the data, we have measured photometry of the different regions ($`a`$ through $`k`$ as indicated in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.) from our ground-based optical images. We have been careful to measure fluxes only within the areas limited by the slits, at similar resolutions, and subtracting scattered light from the QSO where necessary. The solid circles in Fig. 6 indicate the photometry of region $`a`$. Clearly, the older model (red line) is a better fit to the data. It is important to note, however, that if there is considerable reddening by dust along the lines of sight to these regions in the host galaxy, the flux values of $`U^{}`$ will be depressed; therefore the photometric points can only help us to determine upper limits to the ages of the stellar populations. For consistency, the age errors quoted in this section are as defined in §2.1, and do not take into account additional constraints placed by the photometry.
Region $`a`$ corresponds to the west side of the arc-shaped structure (the “horseshoe”; Hamilton & Keel, 1987), $`4`$″ south of the QSO nucleus. The HST F702W image in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. shows that region $`a`$ is formed by multiple knots, much like those of regions $`a`$ and $`c`$ in IRAS 07598+6508. However this structure, unlike those of IRAS 07598+6508, has a remarkably similar morphology from the near UV to the near IR as shown in the top right panel of Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.. As discussed above, we determine an age of 42 (+22, -17) Myr. This age is younger than the 225 Myr estimated from $`UB`$ colors of this region by Surace & Sanders (2000).
While several authors have noted the blue color of the “horseshoe” (e.g., Kodaira et al., 1979; Hutchings & Neff, 1987), we find a region that is relatively much brighter at shorter wavelengths 16″ south of the QSO, labeled $`c`$ in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.. This region appears as a bright blue extended blob dominating the $`U^{}B`$ color map in Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.. The spectrum has a very steep blue continuum, and very strong emission lines at the same redshift as the absorption lines, likely from an H II region (Hamilton & Keel, 1987). The top panel of Fig. 5 shows the spectrum of region $`c`$ with the Balmer emission lines subtracted as described in §2.1. We find an age for this region of 5 Myr. As such a young starburst age is comparable to the expected duration of a typical individual starburst, and we are likely observing a collection of several such starbursts, it is virtually impossible to distinguish between continuous star formation and instantaneous bursts. The age we find is, therefore, simply indicative of ongoing star formation. As in region $`a`$, we find degeneracy in the modeling of this spectrum, with a second $`\chi ^2`$ minimum at 100 Myr. The near-UV continuum of the older model, however, falls $`2\sigma `$ below the $`U^{}`$ photometric point. Neither the deep optical ground-based images nor the HST image show evidence for stellar associations in region $`c`$. This is somewhat surprising since, although the starburst population contributes only 1% to the luminous mass along the line of sight, its flux amounts to 42% of the total flux at 5000 Å(rest wavelength), and even more at shorter wavelengths.
Regions $`b`$ and $`e`$ appear as single clumps in the images and have ages of 140 (+80, $``$70) Myr and 180 (+60, $``$80) Myr respectively. The starburst populations dominate the spectra in these regions, contributing $`75`$% of the total flux at rest frame 5000 Å and up to 36% of the total luminous mass along the line of sight. Regions $`d`$, $`f`$ (Fig. 5), $`g`$, and $`h`$ all show very similar spectra, as expected from the color maps (Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.), with ages between 300 and 360 Myr. The starbursts in these regions typically contribute only $``$50% to the total flux and $`15`$% of the total luminous mass, except for region $`d`$, for which these values are very similar to those of regions $`b`$ and $`e`$.
We detect very weak continuum from the north tail (region $`l`$) at the $`2.5\sigma `$ level. The continuum is slightly blue, and there is a hint of Mg I$`b`$ at the redshift of the host galaxy. This region corresponds to the blue region in the $`BV`$ color map, just before the condensation at the end of the north tail. It is possible that this region contains a knot of star formation like those found in the tail of Mrk 1014.
Lípari (1994) reports the presence of an extended Na I D BAL in the west side of the host galaxy. The only place where we see evidence for this feature is in region $`j`$. The spectrum of $`j`$, uncorrected for QSO scattered light, shows a wide absorption feature ($`1900`$ km s<sup>-1</sup> FWHM compared to $`1050`$ km s<sup>-1</sup> for the QSO Na I D BAL) as well as narrow absorption and emission (see below) at the redshift of the stellar component. After correcting for QSO contamination, the feature appears narrower, but still at $``$6400 km s<sup>-1</sup> with respect to the emission line. The hypothetical BAL feature is, however, much weaker than that of the QSO, and it could be an artifact of an imperfect subtraction of the QSO light.
While the evidence for extended BAL is not strong in our data, Hamilton & Keel (1987) find an “excess flux” extending to the blue side of \[N II\] $`\lambda `$6548, which they interpret as a broad blueshifted component of H$`\alpha `$ or H$`\alpha `$ \+ \[N II\] indicative of an outflow with velocities up to 1500 km s<sup>-1</sup>. We observe this excess most clearly in regions $`h`$ and $`g`$ extending to even larger velocities ($`2500`$ km s<sup>-1</sup>). Boroson & Meyers (1992) find some level of excess light in the blue wing of the H$`\alpha `$ line in every low-ionization BAL QSO in their sample and suggest that this excess flux could be from the BAL material itself.
The narrow Na I D emission line shows a peculiar behavior. We observe narrow Na I D emission in regions $`h`$ and $`j`$ on the red side of the absorption line, forming what looks like a P-cygni profile. The only other emission lines evident in region $`h`$ are \[N II\] $`\lambda `$6583, and weak \[S II\] $`\lambda \lambda `$6717,6731, but these lines are slightly blueshifted with respect to the absorption lines. In region $`g`$ the Na I D emission line disappears, but the absorption line becomes very weak as well, so it is possible that the emission line is blueshifted into the absorption line; \[N II\], however, is at a higher redshift than in region $`h`$. Thus, the very low ionization gas appears to be decoupled from the moderately low ionization gas in these regions.
Hutchings & Neff (1987) find a “green” (i.e., visible only in their $`GB`$ image), jet-like feature extending from near the nucleus to the northeast. They suggest it “may be line emission in \[O III\], perhaps tracing ionizing radiation originating in the nucleus”. Our slit goes through the western side of where this feature would be (see Fig. 2 in Hutchings & Neff, 1987), and region $`j`$ should sample the brightest part of the ridge leading to the bright knot. However, we find only weak \[O III\] in the spectrum of this region, certainly much weaker than in other regions we sample. Therefore it is unlikely that this “green” feature, if real, is due to ionized gas.
Images of Mrk 231 show a greatly disturbed host galaxy both in large and small scales (Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.). At large scales, two symmetric tails extend on the east side of the nucleus for $`44`$″ or $`35`$ kpc each, and there is some very extended low surface brightness material east of the north tail. The nucleus is off-center on a bright “plateau” (Hamilton & Keel, 1987) 20 kpc across that shows complex morphology, including linear jet-like structures (region $`g`$), and curved tail-like structures (region $`d`$). Closer to the nucleus (Fig. Stellar Populations in the Host Galaxies of Mrk 1014, IRAS 07598+6508, and Mrk 231<sup>1</sup><sup>1</sup>1Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555.) we find the “horseshoe” 3 kpc to the south with numerous clumps of star formation, and arm-like features spiraling around the nucleus which connect to the outer tidal tails (Surace et al., 1998) This morphology is indicative of a merger between two disk galaxies of similar mass.
While we are unable to obtain stellar spectra within the central kpc ($`<2\mathrm{}`$) of Mrk 231, Smith et al. (1995) have found from UV spectropolarimetry of the nuclear region that high polarization in the optical falls off quite rapidly shortward of $`3000`$ Å; this effect is most easily interpreted as a dilution of the polarized component by O and B stars from an ongoing starburst (Smith et al., 1995). CO and radio observations also show evidence for a centrally concentrated starburst (Downes & Solomon, 1998; Taylor et al., 1999). It is possible that the knots along the arm-like features around the nucleus have ages similar to those of region $`a`$, if not younger.
We find a relatively wide range of starburst ages around the host galaxy. If, once again, we assume a projected velocity of 300 km s<sup>-1</sup>, the tidal tails have a dynamical age of 110 Myr. Mrk 231 is unusual in that there seems to be a significant post-starburst population somewhat older (300–360 Myr) than the tails, which may indicate that wide spread star formation was ignited prior to the stages of final merger. Numerical simulations of mergers (Mihos & Hernquist, 1996) predict such a scenario when the merging galaxies lack a significant bulge to stabilize them against an early dissipation. In CS2000 we noted that the apparent correlation between bulge mass and black-hole mass (Magorrian et al., 1998) suggests a possible correlation between QSO luminosity and the delay in star formation activity. Thus, the early starburst activity in Mrk 231 may indicate that the host galaxy did not have a substantial bulge, which would in turn imply a less massive black-hole that may result into a less luminous AGN. Although we do in fact observe a relatively less luminous AGN in Mrk 231, these connections are highly speculative, and they are based on relations (Magorrian et al., 1998) and models (Mihos & Hernquist, 1996) that still have many uncertainties.
## 6 Discussion
Mrk 1014, IRAS 07598+6508, and Mrk 231 show many similarities. We already knew that these three objects, while being optically or IR selected QSOs, were also part of the ULIG family. In this study we have found additional properties which they share and which may be related to their intermediate position in the FIR diagram.
While Mrk 1014 and Mrk 231 have long been known to have tails, we have shown that IRAS 07598+6508 also has at least one tidal tail. We have found evidence that all three objects have undergone a strong interaction and are now in the final stages of mergers. The morphology of the hosts (highly perturbed galaxies with tidal tails and destroyed disks) as well as the extent of the starbursts indicate that these are all major mergers between galaxies of comparable mass rather than accretion events of low mass dwarf companions (Mihos & Herquist, 1994; Mihos & Hernquist, 1996).
All three objects show spectra typical of E+A galaxies (Dressler & Gunn, 1983), that is, spectra characterized by the simultaneous presence of strong Balmer absorption lines indicative of a young stellar population and features from an older population such as Mg Ib, and the absence of strong emission lines typical of star forming galaxies (we reiterate that the emission lines seen in our spectra generally come from extended gas ionized by the QSO rather than from star forming regions, with a few exceptions). E+A galaxies are frequently linked to galaxy-galaxy mergers and interactions (e.g., Zabludoff, 1996). In the case of objects with E+A spectra undergoing tidal interactions, the young superposed population is clearly related to the interaction. The three objects discussed in this paper, along with 3C 48 (CS2000) and PG 1700+518 (CS97), all have interaction-induced starbursts.
In addition to the similarities noted in the host galaxies of the three objects, there are also some common characteristics in the spectra of the active nuclei themselves. All three objects show strong Fe II emission, with two of these objects (IRAS 07598+6508 and Mrk 231) being “extreme” Fe II emitters (i.e., showing ratios of $`I`$(Fe II $`\lambda 4570`$)/$`I`$(H$`\beta `$) $`>`$ 2; Lípari 1994). The latter two objects, like PG 1700+518 (CS97; SCC98), are also low-ionization BAL QSOs and have weak or absent narrow emission lines.
How do the ages of these starbursts relate to the merger/interaction stage? The dynamical ages for tidal features in these objects are uncertain because of projection effects. However, our rough estimates indicate that in two cases, Mrk 1014 and IRAS 07598+6508, the peak of the starburst occurred after the tidal tails were launched. This requires some mechanism to stabilize the gas contents of the galaxies against bursting in star formation until the later stages of the merger; numerical simulations indicate that a significant bulge in the host can provide this mechanism (Mihos & Hernquist, 1996). Mrk 231, on the other hand, shows starburst ages which indicate that much of the star formation activity may have started before the last stages of the merger; hence, one or both of the merging galaxies may have lacked a significant bulge.
Whenever we have found a significant range of starburst ages in the host galaxies of QSOs, we have also found evidence that the youngest major starburst regions are preferentially concentrated towards the center of the galaxy. In 3C 48 (CS2000) we found the clearest example of this, with starburst ages becoming progressively younger and more dominating as we approached the galaxy/QSO nucleus. Mrk 231 shows older starbursts in the “plateau” extending $`12`$ kpc around the nucleus, with some of the youngest populations only 3 kpc from the nucleus, and possibly even younger populations in the central kpc of the galaxy. Mrk 1014 shows stronger relatively recent star formation activity along the tail than any of the other objects. However, the starburst region we observe within 2 kpc of the nucleus is far more massive, luminous, and larger (Surace et al., 1998) than the starburst regions along the tail, and has the youngest age found in the host.
The interaction histories of the five objects discussed so far clearly favor a strong connection between interactions and vigorous bursts of star formation. Since the gas flows towards the inner regions can not only trigger star formation but also serve as fuel to the QSO, one might expect the age of the QSO activity to be closely related to the age of the initial starbursts in the central regions of the galaxy. However, observations of central starbursts are in every case hampered by the presence of the QSO; even if this region could be observed, the spectrum would likely be dominated by continuing recent starburst activity and not by a starburst that was coincident with the onset of the QSO activity. Furthermore, both starburst and QSO activity may be episodic. All of this is to say that there is an unavoidable intrinsic uncertainty in using starburst ages to place QSOs in an evolutionary sequence. We defer the detailed discussion of a possible age sequence to a subsequent paper where we will present the results of the four remaining objects in our sample (i.e., IRAS 00275$``$2859, IRAS 04505$``$2958, I Zw 1, and PG 1543+489).
We emphasize, though, that many of our general conclusions from the objects discussed here, as well as from our observations of 3C 48 (CS2000), have an interest that is quite independent of any attempt to use some sort of starburst age as a proxy for a QSO age. (1) The confirmation that these are all starburst or post-starburst objects and that they all show obvious tidal tails validates their close connection with other ULIGs, virtually all of which are mergers or strongly interacting pairs. While there has long been strong circumstantial evidence that a large fraction of the QSO population has resulted from triggering of the QSO activity by interactions and mergers, we now have much more direct evidence for this mechanism for at least one subclass of QSOs. (2) Both the spatial distribution and the time history of star formation in a QSO host galaxy give clues to nature of the galaxies that have participated in the merger. While more sophisticated models of star formation during interactions will be necessary to exploit these data fully, we already have some hints in terms of the enhanced star formation along the leading edge of the tails in 3C 48 and Mrk 1014, and in the relative ages of the tail structures and the star-forming regions contained within them. (3) The youth of the stellar populations in these objects reinforces previous suggestions connecting strong Fe II emission and low-ionization BAL features with the relatively recent triggering of QSO activity. We will discuss these connections in detail in the paper presenting the observations of the remaining four objects in our sample.
We thank Gerbs Bauer, Scott Dahm, and Susan Ridgway for assisting in some of the observations, and Bill Vacca for helpful discussions about IMFs. We also thank the referee, Dean Hines, for his very prompt review of the paper and his suggestions, which helped us improve both its content and its presentation. This paper was partly written while both authors were visitors at the Research School of Astronomy and Astrophysics of the Australian National University, and we thank both the Director, Jeremy Mould, and the staff there for their hospitality. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. This research was partially supported by NSF under grant AST95-29078.
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# On factorisation at small 𝒙Work supported by E.U. QCDNET contract FMRX-CT98-0194 and by the Fondazione A. della Riccia.
## 1 Introduction
The standard approach to calculating any quantity in Deep Inelastic Scattering (DIS) or Drell-Yan type processes is to assume factorisation, based on the operator product expansion. This allows us to consider parton distributions which evolve with scale as
$$\frac{g_\omega (Q^2)}{\mathrm{ln}Q^2}=\gamma _\omega (Q^2)g_\omega (Q^2)+𝒪\left(\frac{1}{Q^2}\right),$$
(1)
where $`\omega `$ is the Mellin moment of the parton distribution. Cross sections are then expressed in terms of the parton distribution at scale $`Q^2`$, multiplied by a coefficient function, which is itself a function of only $`\omega `$ and $`\alpha _\mathrm{s}(Q^2)`$, plus corrections of order $`1/Q^{2n}`$.
A few years ago however, it was suggested in that at small $`x`$ (small $`\omega `$) diffusion might place a limit on factorisation: namely that beyond a certain value of $`x`$ factorisation would break down. The argument was the following: at moderate values of $`x`$ parton evolution in $`Q^2`$ can be associated with a chain of emissions ordered in $`Q^2`$. Corrections to factorisation can arise from chains which evolve from some starting scale $`Q_{\text{start}}^2`$ up to a scale $`Q^2`$, evolve down to a non-perturbative scale $`\mathrm{\Lambda }^2`$ and then evolve back up to scale $`Q^2`$. At moderate $`x`$, such contributions are suppressed by a power of $`\mathrm{\Lambda }^2/Q^2`$. At small $`x`$ however, this process of going up and down in scale becomes an essential element of the evolution, and is no longer suppressed by powers of $`Q^2`$: so it would seemed that factorisation might break down at the values of $`x`$ such that this diffusion enters into the non-perturbative region. More recently this has been used as an argument for parameterising and fitting the small-$`x`$ behaviour of splitting functions, rather than trying to calculate them from first principles .
In plausibility arguments were presented suggesting that there would be no such breakdown. Physically speaking, perhaps the simplest explanation is to say that evolution that, while going down in $`x`$ drops below one’s starting scale $`Q_{\text{start}}^2`$, contributes to the small-$`x`$ part of the gluon distribution at the starting scale. Therefore it should not be included in the anomalous dimension in order to avoid counting it twice, so the anomalous dimension contains only evolution which remains above the starting scale, ensuring that it is perturbative.
However this explanation, and the arguments in , are not entirely satisfactory because they are not able to estimate the size or functional dependence on $`x`$ of any violations of factorisation.
In , a model was studied which was qualitatively similar to BFKL in that it had diffusion, and which correctly reproduced collinear evolution at all orders. It was possible to show that this model leads to *exact factorisation*, namely that there were no $`1/Q^2`$ corrections at all. Thus diffusion on its own is not incompatible with factorisation.
This collinear model had two particularities: the branching kernel, when converted to $`\gamma `$ space (the Mellin transform variable conjugate to transverse momentum) contained poles only at $`\gamma =0`$ and $`\gamma =1`$, i.e. the leading collinear and anti-collinear poles. The BFKL equation instead has poles at all integer values of $`\gamma `$. The second particularity was related to the scale of $`\alpha _\mathrm{s}`$: in the collinear model the scale was chosen to be $`\mathrm{max}(k^2,k_{}^{}{}_{}{}^{2})`$, where $`k`$ and $`k^{}`$ are the exchanged transverse momenta before and after the branching. In the BFKL equation, the next-to-leading corrections seem to indicate that the correct scale is rather $`q=|\stackrel{}{k}\stackrel{}{k}^{}|`$ which has the property that it can go to zero even when both $`k`$ and $`k^{}`$ are large.
Thus there is a need for a study of factorisation within the BFKL equation. The complications mentioned above make this quite difficult to do analytically. We therefore adopt a two-pronged approach. In section 3 we consider an extension of the collinear model with additional poles at $`\gamma =1`$ and $`\gamma =2`$. Some of its gross features can be deduced analytically and can be expected to carry over to the full BFKL equation.
In section 4 we then present a simple numerical method which for the first time allows the extraction of the exact (effective — see below) anomalous dimensions in the full leading logarithmic BFKL equation, with all-order running coupling corrections, allowing a detailed study of factorisation not limited by the simplifying assumptions of any given particular model.
## 2 Understanding factorisation
We shall start off by discussing factorisation for the non-integrated gluon distribution, or equivalently for the gluon Green’s function. Conceptually it it somewhat simpler — the generalisation to the integrated gluon distribution is then relatively straightforward.
We shall discuss factorisation by considering a number of models. We study a gluon Green’s function $`G_\omega (t,t_0)`$, $`t=\mathrm{ln}Q^2`$, in analogy with the unintegrated gluon distribution in the proton. and consider as the non-perturbative aspects of the problem the value of $`t_0`$ (the lower hard scale of the process) and the regularisation of the running coupling in the infra-red. The latter will be represented by some scale $`\overline{t}`$ at which the coupling is cut off (we could equally have chosen to freeze the coupling at that scale). The small-$`x`$ properties of the Green’s function will depend on both $`t_0`$ and $`\overline{t}`$. More specifically the position of the leading pole in the $`\omega `$-plane will depend on the value of $`\overline{t}`$ and its normalisation on both $`\overline{t}`$ and $`t_0`$.
If factorisation holds then the (non-integrated) effective anomalous dimension defined by
$$\stackrel{~}{\gamma }_\omega (t)=\frac{1}{G_\omega (t,t_0)}\frac{G_\omega (t,t_0)}{t},$$
(2)
should be independent of both $`\overline{t}`$ and $`t_0`$ at least to within higher-twist terms suppressed at least as $`e^t`$.
This is not quite a sufficient condition: indeed if the anomalous dimension contains higher-twist pieces which grow as a sufficiently large power of $`x`$ then these could dominate the scaling violations making them impossible to predict. Since the Pomeron singularity will affect $`G_\omega `$ anyway, yielding a large power behaviour of this kind, the problem arises of understanding whether approximate factorisation is still preserved in the small-$`x`$ region.
Since we are interested in pieces enhanced at small-$`x`$ we need to understand the singularity structure of the anomalous dimension. There are two possible origins for singularities of $`\stackrel{~}{\gamma }_\omega `$. One possible origin is that at some $`\omega `$, $`G_\omega /G_\omega `$ contains a non-factorisable singularity, the other is for $`G_\omega `$ to be zero while $`G_\omega `$ is finite.
Now, the simplest collinear model of provided us with a mechanism by which the Pomeron singularity in $`G_\omega `$ can be consistent with exact factorisation: the leading singularity of $`\stackrel{~}{\gamma }_\omega `$ was found to come from a $`t`$-dependent *zero* of $`G_\omega `$, leaving the Pomeron factorised away. We shall see in the following that this basic mechanism is still at work in collinear models with higher twist terms and in the BFKL equation itself, thus suppressing the non-factorisable singularities.
### 2.1 Recalling the structure of the collinear model
The collinear model of ref. gave a powerful handle for the study of anomalous dimensions. It was described by an equation of the form
$$\omega G_\omega (t,t_0)=\delta (tt_0)+𝑑t^{}K_2(t,t^{})G_\omega (t^{},t_0),$$
(3)
whose kernel is
$$K_2(t,t^{})=\overline{\alpha }_\mathrm{s}(t)\mathrm{\Theta }(tt^{})+\overline{\alpha }_\mathrm{s}(t^{})\mathrm{\Theta }(t^{}t)e^{(t^{}t)}.$$
(4)
In $`\gamma `$-Mellin transform space with respect to $`Q^2=\mathrm{exp}(t)`$, the leading order part of the kernel gives the following characteristic function:
$$\chi ^{2\mathrm{pole}}(\gamma )=\frac{1}{\gamma }+\frac{1}{1\gamma },$$
(5)
for which reason it is also referred to as a $`2`$-pole model. It has the convenient property that it can be expressed in terms of a second order differential equation, whose solutions have the following factorised form:
$$G_\omega (t,t_0)=F_\omega ^L(t_0)F_\omega ^R(t),t>t_0,$$
(6)
where $`F_\omega ^L`$ and $`F_\omega ^R`$ are the linearly independent solutions of the homogeneous equation (eq. (3) without the $`\delta `$-function term), that are regular to the left (negative $`t_0`$) and to the right (positive $`t`$) respectively.
The fact that $`F_\omega ^R(t)`$ is regular at large $`t`$ means that it is independent of $`\overline{t}`$ for $`t>\overline{t}`$. It is free of singularities in $`\omega `$-space, as illustrated in figure 1. All non-perturbative dependence, in particular the poles governing the high-energy behaviour of $`G`$, is contained in $`F_\omega ^L`$, also illustrated in figure 1.
One sees from these plots that zeroes of $`G_\omega `$ can arise from both $`F_\omega ^L`$ and $`F_\omega ^R`$. <sup>1</sup><sup>1</sup>1Note however that $`F_\omega ^L`$ does not always have zeroes: their presence depends on a variety of factors such as the relative sign of successive divergences, which are determined by the non-perturbative parameters. Only in the situation of $`t_0\overline{t}`$ is at least one zero guaranteed, just to the left of the leading divergence. However since the $`t`$-dependence lies entirely in $`F_\omega ^R`$, the $`t_0`$-dependent, non-factorisable zeroes of $`F_\omega ^L`$ lead to zeroes of both $`G_\omega `$ and of its $`t`$-derivative. Therefore they do not lead to divergences of the anomalous dimension. Only the $`t`$-dependent zeroes of $`F_\omega ^R`$ lead to poles in $`\gamma (\omega )`$, $`_tF_\omega ^R(t)`$ not usually being zero when $`F_\omega ^R(t)`$ is zero.
This fact has the consequence mentioned before: the Pomeron singularity, present in $`F_\omega ^L(t_0)`$, does not occur in the anomalous dimension, which is singular at $`\omega _c(t)`$, the leading zero of $`F_\omega ^R`$ (while we refer to the leading zero of $`F_\omega ^L(t_0)`$ as $`\omega _0(t_0)`$).
## 3 The 4-pole collinear model
One of the main differences between the BFKL equation and the collinear model just-discussed lies in the presence in the BFKL equation of poles of $`\chi (\gamma )`$ at all integer values of $`\gamma `$. The poles beyond $`\gamma =0`$ and $`\gamma =1`$ can give rise to higher-twist effects, in which we are particularly interested in this article.
A useful stepping stone to the full BFKL equation is an equation containing the first subleading higher-twist parts, namely poles at $`\gamma =1`$ and $`\gamma =2`$. The equation for the Green’s function is
$$\omega G_\omega (t,t_0)=\delta (tt_0)+𝑑t^{}K_4(t,t^{})G_\omega (t^{},t_0),$$
(7)
where the kernel is
$$\begin{array}{c}K_4(t,t^{})=\overline{\alpha }_\mathrm{s}(t)\mathrm{\Theta }(tt^{})\left(1+e^{(tt^{})}\right)+\overline{\alpha }_\mathrm{s}(t^{})\mathrm{\Theta }(t^{}t)\left(e^{(t^{}t)}+e^{2(t^{}t)}\right)\hfill \\ \hfill \frac{4}{3}\overline{\alpha }_\mathrm{s}(t)\delta (tt^{}).\end{array}$$
(8)
The $`\delta `$-function term is included so that the value of $`\chi _0(1/2)`$ be the same as for the 2-pole collinear model:
$$\chi ^{4\mathrm{pole}}(\gamma )=\frac{1}{1+\gamma }+\frac{1}{\gamma }+\frac{1}{1\gamma }+\frac{1}{2\gamma }\frac{4}{3}.$$
(9)
The physics content of this formulation is quite similar to the leading and first-subleading collinear and anti-collinear parts of the BFKL equation. A slight difference exists in that the higher-order higher-twist pieces of the BFKL equation also contain a non-local dependence on $`\alpha _\mathrm{s}`$. Their inclusion would considerably increase the difficulty of the analytic treatment.
This model as it stands can be expressed as a fourth-order differential equation (or two coupled second order differential equations, etc.), using the same approach as in , as discussed in detail in the Appendix. It is thus similar to the diffusion models with higher twist terms . Without having to write down the full equations in detail, we describe here some important properties of its solutions.
There are four linearly-independent solutions of the homogeneous equation. We denote them by $`F^{R,0}`$, $`F^{R,1}`$, the leading and sub-leading twist right-regular solutions, which for large $`t`$ go as a constant and $`e^t`$ respectively; and by $`F^{L,a}`$, $`F^{L,b}`$ for the two left-regular solutions. For $`t>t_0`$ the Green’s function has the expression
$$G_\omega (t,t_0)=F_\omega ^{L,a}(t_0)F_\omega ^{R,0}(t)+F_\omega ^{L,b}(t_0)F_\omega ^{R,1}(t).$$
(10)
as expected from , and proved in the Appendix. There we also show how to determine the $`F^L`$’s so as to satisfy the boundary conditions at $`t=t_0`$. The solution for $`t<t_0`$ can then also be obtained, by a simple exchange of variables. As in the 2-pole model, the right-regular solutions are independent of the non-perturbative parameters of the problem. The Pomeron singularity and the non-perturbative dependence all enter into the left-regular solutions.
It is convenient to divide the solution into its leading and higher-twist parts, $`G_\omega ^0`$ and $`G_\omega ^1`$, corresponding to the first and second terms of (10) respectively. For large $`tt_0,\overline{t}`$, we have that $`G_\omega ^1`$ is strongly suppressed (by a relative amount $`e^t`$) compared to $`G_\omega ^0`$; $`G_\omega ^0`$ should be qualitatively similar to the solution of the 2-pole model, having zeroes of perturbative and non-perturbative origin. We know therefore, from the start, that the violation of factorisation is uniformly of higher twist, just because of the additive nature of (10).
On the other hand, one can look at the same problem from the standpoint of the singularities of the anomalous dimension. Let us concentrate on the case with a discrete spectrum (running coupling with cutoff), in which we have to look for zeroes of $`G_\omega `$. The fact that $`G_\omega ^1`$ is strongly suppressed means that the distribution of zeroes of $`G_\omega `$ is determined essentially by the positions of the zeroes of $`G_\omega ^0`$. But in general where $`G_\omega ^0`$ is zero, $`G_\omega ^1`$ will be non-zero, causing the zeroes of $`G_\omega `$ to be slightly shifted compared to the zeroes of $`G_\omega ^0`$. The size of this shift will be related to the relative sizes of $`G_\omega ^1`$ and $`G_\omega ^0`$, i.e. it will be of the order of of $`e^t`$.
The effect of the shift will be different according to whether the zero is of perturbative or non-perturbative origin. In the case of the leading zero of perturbative origin, then the leading perturbative pole of the anomalous dimension is shifted, and its normalisation changes — both the effects are of order $`e^t`$ (the relation between relative changes to the normalisation and the position of the divergence depends on the nature of the $`t`$ and $`\omega `$-dependence of $`G^1`$, and so is difficult to predict).
The consequences of the rightmost (NP) zero $`\omega _0(t_0)`$ being shifted are somewhat more interesting, and possibly dangerous for factorisation. We recall that in the two-pole model, the non-perturbative zeroes did not lead to divergences of the anomalous dimension because the zero was in $`F^L`$ and so the derivative $`_tG(t,t_0)=F^L(t_0)_tF^R(t)`$ was also zero when $`G`$ was zero. This is no longer true when higher twist terms are present, and a singularity around $`\omega _0`$ is expected. The exact value of $`\omega _0`$, and indeed even the existence of this zero, depend rather subtly on the values of $`\overline{t}`$ and $`t_0`$. However, for $`t_0`$ significantly larger than $`\overline{t}`$ one expects that the zero exists, and is driven by the Pomeron term in $`F^L(t_0)`$, so that $`\omega _{}\omega _0e^{t_0}`$ is rather small and $`\omega _0`$ may be leading.
To see in detail what happens in the 4-pole model if $`\omega _0`$ is leading, we approximate $`G_\omega ^0(t,t_0)`$ around its zero by $`𝒩_0(\omega \omega _0)F_{\omega _0}^{R,0}(t)`$, where $`𝒩_0=_\omega F_\omega ^{L,a}(t_0)|_{\omega =\omega _0}`$, and $`G^1(t,t_0)`$ by $`𝒩_1F_{\omega _0}^{R,1}(t)`$. It is useful to define partial anomalous dimensions separately for $`F_{\omega _0}^{R,0}(t)`$ and $`F_{\omega _0}^{R,1}(t)`$:
$`\gamma _\omega ^0(t)`$ $`={\displaystyle \frac{_tF_\omega ^{R,0}(t)}{F_\omega ^{R,0}(t)}}={\displaystyle \frac{\overline{\alpha }_\mathrm{s}}{\omega }}+\mathrm{},`$ (11)
$`\gamma _\omega ^1(t)`$ $`={\displaystyle \frac{_tF_\omega ^{R,1}(t)}{F_\omega ^{R,1}(t)}}=1+{\displaystyle \frac{\overline{\alpha }_\mathrm{s}}{\omega }}+\mathrm{},`$ (12)
where we have neglected terms of order $`(\overline{\alpha }_\mathrm{s}/\omega )^2`$ in the expansion of the anomalous dimensions. Thus $`F_{\omega _0}^{R,1}(t)`$ is suppressed compared to $`F_{\omega _0}^{R,0}(t)`$ by an amount
$$\frac{F_{\omega _0}^{R,1}(t)}{F_{\omega _0}^{R,0}(t)}\mathrm{exp}\left[^t𝑑t^{}\left(\gamma _{\omega _0}^1(t)\gamma _{\omega _0}^0(t)\right)\right]e^{t+𝒪\left(\alpha _\mathrm{s}(t)/\omega _0^2\right)}.$$
(13)
The rightmost zero of $`G_\omega (t,t_0)`$, $`\omega _0^{}`$, is therefore shifted compared to that of $`F_\omega ^{R,1}(t)`$ by an amount
$$\omega _0^{}\omega _0\frac{𝒩_1F_{\omega _0}^{R,1}(t)}{𝒩_0F_{\omega _0}^{R,0}(t)}e^{t+𝒪\left(\alpha _\mathrm{s}/\omega _0^2\right)}.$$
(14)
At $`\omega _0^{}`$ the $`t`$ derivative of $`G_\omega `$ is non-zero, leading to a pole of the anomalous dimension, whose residue $`R(\omega _0^{})`$ is
$$R(\omega _0^{})=(\gamma _{\omega _0}^1(t)\gamma _{\omega _0}^0(t))\frac{𝒩_1F_{\omega _0}^{R,1}(t)}{𝒩_0F_{\omega _0}^{R,0}(t)}=\left(1+𝒪\left(\alpha _\mathrm{s}^2/\omega ^2\right)\right)(\omega _0^{}\omega _0)$$
(15)
So we have that the residue of the leading NP pole is the same as the shift of the zero and *both are higher twist*. In the case of the BFKL equation (section 4) one can test to see if this remains true, in order to verify that the same mechanisms are at work there as in the 4-pole model.<sup>2</sup><sup>2</sup>2We note that for the BFKL equation including next-to-leading corrections, the relation between $`\omega _0^{}\omega _0`$ and the residue of the pole can be modified by pieces of relative order $`\alpha _\mathrm{s}`$.
Since $`\omega _0`$, which may approach $`\omega _{}`$, is likely to be to the right of the perturbative pole $`\omega _c`$, at small $`x`$ this higher-twist non-perturbative contribution will dominate the anomalous dimension. At first sight this might seem to have worrying implications for the prediction of small-$`x`$ scaling violations. But in the present model we know that this is not the case, because of the validity of (10). From the point of view of $`t`$-evolution we are saved by the fact that at small $`x`$, in the convolution of the splitting function with the non-perturbative input distribution (which grows as $`\omega _{}>\omega _0`$), the higher-twist part of the effective splitting function gives a contribution of order
$$\frac{e^t}{\omega _{}\omega _0}$$
while the perturbative contribution is of order
$$\frac{\alpha _\mathrm{s}(t)}{\omega _{}},$$
both growing as $`x^\omega _{}`$. We thus recover in the latter contribution the Pomeron part of the leading twist term in (10), that we know to be factorised. Thus there is always a value of $`Q^2`$ such that the higher-twist corrections can be ignored at all $`x`$.
A small point worth bearing in mind is that our analysis so far has always been for the anomalous dimensions related to unintegrated gluon distributions. In practice one is more interested in the anomalous dimensions of the integrated distributions. It turns out that their properties are very similar: this is because the $`n`$-pole models can be expressed in terms of $`n`$ coupled linear differential equations, and the unintegrated and integrated gluon distributions are simply different linear combinations of the components of the equations.
## 4 The BFKL equation
We shall study the leading-order BFKL equation including a running coupling,
$$\begin{array}{c}𝒢(x,k,Q_0)=\delta (k^2/Q_0^21)+_x^1\frac{\mathrm{d}z}{z}\frac{\mathrm{d}^2\stackrel{}{q}}{\pi q^2}\overline{\alpha }_\mathrm{s}(q^2)\hfill \\ \hfill \times \left[\frac{k^2}{|\stackrel{}{k}\stackrel{}{q}|^2}𝒢(x,|\stackrel{}{k}\stackrel{}{q}|,Q_0)\mathrm{\Theta }(kq)𝒢(x,k,Q_0)\right].\end{array}$$
(16)
For notational convenience we have switched to using transverse momenta $`k,Q_0`$, rather than the logs of their squares $`t,t_0`$. The choice of the emitted transverse momentum as the scale for the running coupling, $`\alpha _\mathrm{s}(q^2)`$, is suggested by the form of the NLO corrections to the kernel . A normal DGLAP gluon distribution is expressed in terms of $`𝒢(x,k,Q_0)`$ via $`k`$-factorisation:
$$xg(x,Q^2)=^{Q^2}\frac{\mathrm{d}^2k}{\pi k^2}𝒢(x,k,Q_0).$$
(17)
In practice rather than solving the integral equation (16), it is easier to solve the related differential equation
$$\begin{array}{c}\frac{𝒢(x,k,Q_0)}{\mathrm{ln}1/x}=\frac{\mathrm{d}^2\stackrel{}{q}}{\pi q^2}\overline{\alpha }_\mathrm{s}(q^2)\left[\frac{k^2}{|\stackrel{}{k}\stackrel{}{q}|^2}𝒢(x,|\stackrel{}{k}\stackrel{}{q}|,Q_0)\mathrm{\Theta }(kq)𝒢(x,k,Q_0)\right],\hfill \end{array}$$
(18)
with initial condition $`𝒢(1,k,Q_0)=\delta (k^2/Q_0^21)`$.
### 4.1 The extraction of anomalous dimensions
Naively to obtain the effective splitting function, one would determine $`𝒢`$ and then solve for the function $`P_{gg,\mathrm{eff}}(z)`$ such that
$$x_{\mathrm{ln}Q^2}g(x,Q^2)=x_x^1\frac{dz}{z}P_{gg,\mathrm{eff}}(z)g(x/z,Q^2).$$
(19)
However such an approach turns out to be subject to considerable numerical instabilities. The reason is that any method of solution for $`𝒢`$ introduces small errors (typically of the relative order of $`10^2`$$`10^3`$). When carrying out the deconvolution it generally turns out that $`P(z)`$ for small $`z`$ contributes only a small amount to the scaling violations ($`xg(x,Q^2)`$ grows as $`x^\omega _{}`$ whereas perturbatively, $`P(z)`$ grows as $`x^{\omega _c}`$ and $`\omega _c\omega _{}`$). When $`x`$ is such that the small error on $`𝒢(x)`$ is of the same order as the contribution to the scaling violations from $`P(x)`$, then we no longer have a handle on the splitting function.
A solution is to choose an inhomogeneous term such that $`g(x,Q^2)`$ is independent of $`x`$. Then, for a given $`x`$, the convolution (19) is dominated by small $`z`$’s and small errors on $`g(x)`$ are no longer amplified when translated to $`P(x)`$. We introduce $`(x,k)`$ as being the unintegrated gluon distribution which, integrated, gives $`xg(x,Q^2)=1`$. It satisfies the equation
$$\begin{array}{c}\frac{(x,k)}{\mathrm{ln}1/x}=f(x,k)\delta (k^2/Q_0^21)+\frac{\mathrm{d}^2\stackrel{}{q}}{\pi q^2}\overline{\alpha }_\mathrm{s}(q^2)\hfill \\ \hfill \times \left[\frac{k^2}{|\stackrel{}{k}\stackrel{}{q}|^2}(x,|\stackrel{}{k}\stackrel{}{q}|)\mathrm{\Theta }(kq)(x,k)\right],\end{array}$$
(20)
where $`(1,k)=\delta (k^2/Q_0^21)`$ and $`f(x)`$ is given implicitly by
$$f(x)=^{Q^2}\frac{\mathrm{d}^2k}{\pi k^2}\frac{\mathrm{d}^2\stackrel{}{q}}{\pi q^2}\overline{\alpha }_\mathrm{s}(q^2)\left[\frac{k^2}{|\stackrel{}{k}\stackrel{}{q}|^2}(x,|\stackrel{}{k}\stackrel{}{q}|)\mathrm{\Theta }(kq)(x,k)\right].$$
(21)
It is trivial to verify that this leads to $`xg(x,Q^2)=1`$. The form of $`f(x)`$ depends on the $`Q^2`$ value at which we intend to consider the splitting function and on the initial scale $`Q_0`$. Using (19) it is now simple to obtain the effective splitting function:
$$xP_{gg,\mathrm{eff}}(x)=\frac{(x,Q)}{\mathrm{ln}1/x}.$$
(22)
This turns out to be numerically stable, at least until $`xP_{gg,\mathrm{eff}}`$ becomes comparable to the inverse of the machine precision.
Equation (20) is solved by discretising $`(x,k)`$ uniformally in $`\mathrm{ln}k`$ space, and then applying standard Runge-Kutta techniques for solving the resulting matrix differential equation. This method has the advantage over other potentially faster methods, such as a representation with a basis of Chebyshev polynomials , that it quite easily accommodates the large variations that arise in the value of $`(x,k)`$ (and errors are at worst of the relative order of the discretisation interval).
For $`\overline{\alpha }_\mathrm{s}`$ we take the asymptotic formula, i.e.
$$\overline{\alpha }_\mathrm{s}(q^2)=\frac{\mathrm{\Theta }(q^2\overline{Q}^2)}{b\mathrm{ln}q^2/\mathrm{\Lambda }^2},$$
(23)
where $`b=11/12`$ (we work with zero flavours). The cutoff at small momenta corresponds to the regularisation prescription used in the previous section for the collinear model.
### 4.2 Results
This section has two aims. Firstly to demonstrate that for the BFKL equation the splitting function is a truly perturbative quantity, and that any non-perturbative dependence is higher-twist. And secondly, to show that our understanding of higher-twist effects as obtained from the 4-pole model, carries over to the BFKL equation.
Let us start by examining some concrete examples of effective splitting functions. Figure 2 illustrates the effective splitting function as a function of $`x`$ in three situations, all with the same value of $`t`$, but different sets (a, b and c) of non-perturbative parameters, $`\overline{t}`$ and $`t_0`$. Going down in $`x`$ from $`x=1`$, one sees that initially the three splitting functions are almost identical (the inset with the larger scale reveals small differences between them). For moderately small $`x`$ the splitting function actually decreases (this phenomenon was recently observed also by ), and then starts to grow as $`x^{\omega _c}`$. The late onset of the power growth is related to the fact that in $`\omega `$-space the PT pole of the anomalous dimension has a residue of order $`\overline{\alpha }_\mathrm{s}^2`$.
At a certain point, two of the curves (b and c) change sign (since we use a logarithmic scale and plot the absolute value of the splitting function, the change of sign appears as a downward cusp) and start to grow with a much larger power ($`\omega _0^{}`$). This is the non-perturbative higher-twist component of the splitting function discussed in the section 3. As can be seen the exact value of $`\omega _0^{}`$ depends on the non-perturbative parameters. Furthermore there are situations (curve a) in which there is no NP power growth at all, corresponding to the absence of NP zeroes in the Green’s function.
Figure 2 is not sufficient to demonstrate that the NP corrections are truly higher twist. First we consider the ‘PT part’ of the splitting function. Figure 3 shows the ratio of two effective splitting functions, obtained with different non-perturbative parameter sets (a and b), chosen such that there is no component with the large NP power growth (which would complicate the interpretation of the ratio). We see that the NP parameters affect both the normalisation of the PT splitting function, and the exact value of the power growth, since the ratio grows as a power. This is as predicted in the 4-pole model, being due to the shift of the PT zeroes of the Green’s function. We observe that the effect on the power is relatively small (note the $`x`$ scale), and that from a practical point of view it will mostly be the effect on the normalisation that will be of interest. The curve at the higher value of $`t`$ shows a significantly decreased dependence on the NP parameters, confirming that their effect is higher-twist.<sup>3</sup><sup>3</sup>3The attentive reader will have noticed that the modification of the PT exponent changes sign at the higher $`t`$ value — indeed it turns out that the dependence on the NP parameters, while decreasing at least as fast as $`e^t`$, has a non-trivial $`t`$ dependence for moderate $`t`$ values.
We also need to demonstrate that the NP power component is higher twist. Figure 4 shows the splitting function for three different $`t`$ values, but with the same NP parameters. The initial parts of the splitting function now differ since the PT scales are different; the NP power growth has (approximately) the same power in the three cases, but its normalisation decreases rapidly with increasing $`t`$, roughly as $`e^t`$, confirming that it too is higher twist.
One of the non-trivial features of the 4-pole model was the (approximate) equality between the normalisation of the NP power component, and the quantity $`\omega _0^{}\omega _0`$, the difference between the exponent of the power growth and the position of the NP zero (obtained numerically from the exponent of the NP power growth in the limit $`t\mathrm{}`$). Both quantities are shown in figure 5 over a range of $`t`$ values, illustrating very clearly their closeness (as well as their higher-twist nature). A detailed study reveals that the relative difference between them is consistent with a term of $`𝒪\left(\alpha _\mathrm{s}^2\right)`$, as predicted in section 3. We note also that $`\omega _0`$ is found to be below the NP exponent $`\omega _{}`$ characteristic of the Green’s function, consistent with it being due to a zero to the left of the leading singularity.
## 5 Conclusions
For some time now there has been some debate as to whether diffusion might destroy small-$`x`$ factorisation or make it impossible to perturbatively predict the small-$`x`$ splitting functions . In a model was presented which contained diffusion and the correct collinear limits, but also displayed the property of exact factorisation. It failed though to make any statement about (the functional dependence of) the magnitude of any higher-twist corrections, leaving open the possibility that they could come to dominate at small $`x`$, and thus still destroy factorisation.
This paper has presented a much more complete study of the problem, both by an extension of the model so that it includes leading higher-twist components, and by the development of numerical techniques for studying the effective splitting function in the full (LL with running coupling) BFKL equation.
The basic conclusion of these studies is that higher-twist effects, while present (and dominant at small $`x`$ for the effective splitting function) are truly small *for scaling violations*. An explanation of this fact comes from the analysis of the 4-pole collinear model, whose Green’s function is a sum of two terms, each of which with a NP Pomeron part factorised from the $`t`$-dependence, the second being uniformly of higher twist.
Although oversimplified, the collinear example is expected to exhibit the mechanism at work in a realistic BFKL equation also. In fact, it can be generalised to the case of $`2n`$-poles, for arbitrary $`n`$, the Green’s function being a sum of $`n`$ terms of higher and higher twist . The basic point remains the fact that the Pomeron singularity, although non-perturbative, is stably factorised in front of the $`t`$-dependence in the leading twist term.
## Acknowledgements
We wish to thank Martina Taiuti for several discussions in the early stages of this work. One of us (GPS) would also like to thank Guido Altarelli, Yuri Dokshitzer and Al Mueller for interesting conversations.
## Appendix A Green’s function of the 4-pole model
Starting from eq. (9) there are various ways of deriving the differential equation of the model which differ by the treatment of the higher twist terms.
The simplest way is to insert the expression
$$b\omega t\chi (\gamma )$$
(24)
into the $`\gamma `$-representation of the (regular) solution
$$_\omega (t)=e^{\frac{1}{2}t}F_\omega (t)=\frac{d\gamma }{2\pi i}\mathrm{exp}\left[\left(\gamma \frac{1}{2}\right)t\frac{1}{b\omega }^\gamma \chi (\gamma ^{})𝑑\gamma ^{}\right]$$
(25)
and to notice that the result vanishes by partial integrations. The corresponding differential equation is thus obtained by the replacement $`\gamma \frac{1}{2}_t`$. By thus using the identity
$$\chi (\frac{1}{2}+_t)+\frac{4}{3}=\frac{4_t^23}{_t^4\frac{5}{2}_t^2+\frac{9}{16}}=\frac{𝒩_t}{𝒟_t}$$
(26)
and shifting the $`t`$ variable to incorporate the constant term $`4/3`$, we obtain the Green’s function equation
$`\left[\omega 𝒟_t+𝒩_t\overline{\alpha }_\mathrm{s}(t)\right]g(t,t_0)`$ $`=\left[\omega 𝒟_{t_0}+𝒩_{t_0}\overline{\alpha }_\mathrm{s}(t_0)\right]g(t,t_0)`$
$`={\displaystyle \frac{1}{\omega }}𝒩_t\delta (tt_0)`$ (27)
where $`G_\omega (t,t_0)=\omega ^1\delta (tt_0)+e^{\frac{1}{2}(tt_0)}g_\omega (t,t_0)\overline{\alpha }_\mathrm{s}(t_0)`$ is the gluon Green’s function discussed in the main text.
Note first that $`g(t,t_0)`$ satisfies a fourth order differential equation in both the $`t`$ and $`t_0`$ variables, for any low-$`t`$ regularisation of $`\overline{\alpha }_\mathrm{s}(t)`$, starting from the $`1/bt`$ expression of asymptotic freedom (the regularisation depends on the $`\overline{t}`$ parameter of the main text). Some care is needed in order to treat the boundary conditions that $`g`$ has to satisfy due to the peculiar distribution occurring in the right-hand side of eq. (A). This distribution can be taken into account by assuming $`g(t,t_0)`$ to be a linear combination of left (right) regular solutions of the homogeneous equation, as in (10), satisfying at $`t=t_0`$ the discontinuity requirements
$$\mathrm{\Delta }g=\mathrm{\Delta }g^{\prime \prime }=0,\mathrm{\Delta }g^{}=\frac{4}{\omega ^2},\mathrm{\Delta }g^{\prime \prime \prime }=\frac{7}{\omega ^2}+\frac{16\overline{\alpha }_\mathrm{s}(t_0)}{\omega ^3},(t=t_0).$$
(28)
Given an arbitrary basis $`_L^0(t_0),_L^1(t_0)`$ of left-regular solutions, the Green’s function takes the form
$$\begin{array}{c}g_\omega (t,t_0)=(_R^0(t)_L^a(t_0)+_R^1(t)_L^b(t_0)))\mathrm{\Theta }(tt_0)\hfill \\ \hfill +(_L^0(t)_R^a(t_0)+_L^1(t)_R^b(t_0)))\mathrm{\Theta }(t_0t)\end{array}$$
(29)
and should also be symmetrical under $`t,t_0`$ interchange. Because of (29), the discontinuity conditions (28) at $`t=t_0`$ can be viewed as a system of four linear equations in the four unknowns $`_L^a(t_0)`$, $`_L^b(t_0)`$, $`_R^a(t_0)`$, $`_R^b(t_0)`$, which can be solved by standard methods.
As a consequence, the left-regular solutions $`_L^a`$, $`_L^b`$ occurring in (10) and satisfying (28) are given by the following expressions
$`_L^a(t)`$ $`={\displaystyle \frac{detW(\mathrm{\Delta }(t),_R^1(t),_L^0(t),_L^1(t))}{detW(_R^0,_R^1,_L^0,_L^1)}},`$ (30a)
$`_L^b(t)`$ $`={\displaystyle \frac{detW(_R^0(t),\mathrm{\Delta }(t),_L^0(t),_L^1(t))}{detW(_R^0,_R^1,_L^0,_L^1)}}.`$ (30b)
Here the $`W`$’s are the Wronskian matrices of the corresponding functions, where in the numerators the column vector of derivatives is replaced, in the proper place, by the discontinuity vector ($`i=0,1,2,3`$)
$$\mathrm{\Delta }(t)=(\mathrm{\Delta }g^{(i)})=\frac{1}{\omega ^2}(0,4,0,7\frac{16\overline{\alpha }_\mathrm{s}(t)}{\omega }).$$
(31)
While the Wronskian in the denominator is constant, one can check that the numerators are indeed solutions of the basic homogeneous equation in (A). Furthermore the expressions (29) turn out to be independent of the choice of the left-regular basis, by the linearity properties of Wronskian matrices. Finally the Green’s function is determined by symmetry for $`t<t_0`$.
The explicit determination of $`_L^a(t_0)`$ and $`_L^b(t_0)`$, and of the corresponding Pomeron singularity requires the solution of a matching problem for scattering in a fourth-order framework, depending on the regularisation of $`\overline{\alpha }_\mathrm{s}(t)`$ to the left. The outcome of this procedure is the separation of the left-regular solutions into irregular and regular ones on the right, as follows
$`_L^0(t_0)`$ $`=_R^0(t_0)+\sigma _{00}(\omega )_R^0(t_0)+\mathrm{},`$ (32a)
$`_L^1(t_0)`$ $`=_R^1(t_0)+\sigma _{10}(\omega )_R^0(t_0)+\mathrm{},`$ (32b)
where $`\sigma _{00}(\omega )`$, $`\sigma _{10}(\omega )`$, …are scattering coefficients which carry the (non-perturbative) Pomeron singularity and $`_R^0(_R^0)^1`$, $`_R^1(_R^1)^1`$ are the irregular solutions on the right. The relation of the non-perturbative to perturbative contributions in each of the terms of eq. (10) is thus very similar to that found in the two-pole model .
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# Shell model description of “mixed-symmetry” states in 94Mo
## I Introduction
Recently performed photon scattering experiments and $`\gamma \gamma `$–coincidence studies of the nucleus <sup>94</sup>Mo indicate the existence of low-lying valence shell excitations with proton-neutron symmetry different to that of the ground state in that nucleus. The measurements of absolute $`E2`$ and $`M1`$ transition strengths have been interpreted in terms of $`J^\pi =1^+,2^+,3^+`$ mixed-symmetry states in the framework of the proton-neutron version of the Interacting Boson Model (IBM-2). The proton-neutron symmetry of an IBM-2 wavefunction is quantified by the $`F`$-spin quantum number , which is the isospin for basic proton and neutron bosons. The IBM-2 predicts the lowest-lying collective states to be dominantly isoscalar excitations of the almost proton-neutron symmetric ground state with maximum $`F`$-spin quantum number $`F=F_{\mathrm{max}}`$. This is put in evidence by the existence of $`F`$-spin multiplets with rather constant energies . The IBM-2 predicts also valence shell excitations with wavefunctions, which are not symmetric with respect to the proton-neutron degree of freedom . Such states have $`F`$-spin quantum numbers $`F<F_{\mathrm{max}}`$ and are called mixed-symmetry (MS) states. The observations for the $`1_1^+,2_3^+,3_2^+`$ states of <sup>94</sup>Mo agree with MS assignments . The key-signatures used for the assignments of MS character to these states of <sup>94</sup>Mo were the measured relatively strong $`M1`$ transitions and weakly-collective $`E2`$ transitions to low-lying symmetric states.
Separate proton and neutron quadrupole surface vibrations, which lead to eigenstates with different symmetries with respect to proton-neutron permutations, have been considered in a geometrical model already in the sixties by Faessler . At that time these states were predicted to exist above the particle threshold, at a much higher energy than recently observed. In the eighties the isovector vibrational model was improved giving quantitatively better description of the $`2^+`$ isovector vibrational states. The existence of collective orbital isovector $`M1`$ transitions in the geometrical Two Rotor Model for deformed nuclei was realized by Lo Iudice and Palumbo . A more general type of enhanced magnetic dipole transitions in the valence shell of all open–shell nuclei, not only deformed, has been predicted in the IBM-2 as the decays of MS states. Following these predictions the dominantly isovector $`J^\pi =1^+`$ state, generally called scissors mode, was discovered by Richter et al. in inelastic electron scattering experiments in Darmstadt. The typically fragmented $`1_{\mathrm{sc}}^+`$ scissors mode was further investigated mainly by electron scattering , photon scattering and neutron scattering experiments.
Many theoretical studies were published to explain the structure of this mode, e.g. . The $`1^+`$ mode is expected to be dominantly excited by the isovector part of the $`M1`$ operator indicating its isovector character. The large $`M1`$ transition strength and its close correlation to the collective $`E2`$ excitation strength in deformed nuclei is usually considered an indication of the collective nature of the $`1_{\mathrm{sc}}^+`$ state. Another state with spin different from $`J^\pi =1^+`$ but of similar isovector character, the one-quadrupole phonon $`2_{\mathrm{ms}}^+`$ state, has been identified from $`M1`$ strengths, too. The first $`2_{\mathrm{ms}}^+`$ assignments to states at about 2 MeV in vibrational nuclei around the $`N=82`$ shell closure were based only on small $`E2/M1`$ multipole mixing ratios Lateron, several of these assignments were confirmed by the measurement of relatively large $`M1`$ and weakly-collective $`E2`$ transition strengths . Based on measured $`M1`$ and $`E2`$ strengths the experimentally observed $`1_1^+`$, $`2_3^+`$, and $`3_2^+`$ states in the nucleus <sup>94</sup>Mo were argued to belong to this type of excitations, too .
In nuclei not too far from shell closures the structure of low-lying valence excitations including the states, which are outside of the IBM configurational space, can be described using the shell model. It is the main aim of the present paper to describe the structure of the states observed in in the framework of the nuclear shell model, especially the structure of the $`1_1^+`$, $`2_3^+`$ and $`3_2^+`$ states. We present the results of shell model calculations for the low-lying positive parity states of <sup>94</sup>Mo with spin quantum numbers $`J^\pi =0^+`$$`4^+`$ and we discuss the isotensor character of their electromagnetic transitions.
## II Theoretical Approach
The shell model Hamiltonian is taken as $`H=H_0+V`$ where the mean field is given by
$$H_0=\underset{k}{\overset{n_{\mathrm{val}}}{}}\epsilon _ka_{\rho _k}^+a_{\rho _k},$$
(1)
where $`n_{\mathrm{val}}`$ is a number of single particle states in the adopted valence shells and the residual interaction is
$$V=\underset{\rho _a,\rho _b,\rho _c,\rho _d,J,M,T}{}<(\rho _a\rho _b)_{JT}|V_{12}|(\rho _c\rho _d)_{JT}>(a_{\rho _a}^+a_{\rho _b}^+)_{JM}^T(a_{\rho _c}a_{\rho _d})_{JM}^T.$$
(2)
Here $`a_\rho ^+`$ creates and $`a_\rho `$ annihilates a particle in the single particle orbital $`|\rho |n,l,j,m_j,t=1/2,t_z`$ and T is the isospin of the coupled particles. The first term $`H_0`$ is the Hamiltonian of the noninteracting particles. The residual interaction we used is the Surface Delta Interaction (SDI) . This interaction contains strong pairing and quadrupole parts and higher multipolarity components, which are weaker than the first two. The SDI is an extremely simple interaction from the mathematical point of view. The two body matrix elements of the SDI are:
$$<\rho _a\rho _b|V_{SDI}(𝐫_1,𝐫_2)|\rho _c\rho _d>_{JT}=4\pi A_T^{}<\rho _a\rho _b|\delta (\mathrm{\Omega }_{12})\delta (r_1R)\delta (r_2R)|\rho _c\rho _d>_{JT}$$
(3)
where $`\mathrm{\Omega }_{12}`$ is the angle between the interacting particles, $`R=1.2A^{1/3}`$ fm is the nuclear radius, and $`A_T^{}`$ is the strength constant of the SDI. There are three parameters $`A_{}^{}{}_{T=1}{}^{\rho ^{},\rho }`$ ($`\rho ,\rho ^{}\{p,n\}`$) that describe the interaction in the T=1 channel and one parameter $`A_{}^{}{}_{T=0}{}^{pn}`$ describing the interaction in the T=0 channel. The fitted interaction parameters $`A_T^{\rho ^{},\rho }`$ are connected to the parameters from Eq.(3) by the following expression:$`A_T^{\rho ^{},\rho }`$= $`A_T^{\rho ^{},\rho }<\delta (r_\rho R)\delta (r_\rho ^{}R)>`$, where the radial matrix elements $`<\delta (r_\rho R)\delta (r_\rho ^{}R)>`$ are supposed to be independent of the single particle states involved (see for details ).
For the shell model description of <sup>94</sup>Mo one may want to consider the $`N=Z=50`$ closed shell nucleus <sup>100</sup>Sn as the inert core. In order not to have to treat a too large model space we consider eight proton holes in the proton shells $`\pi g_{9/2}`$ and $`\pi p_{1/2}`$ for the description of $`{}_{42}{}^{}{}_{}{}^{94}`$Mo. Because of the Pauli principle this problem is equivalent to the consideration of four proton particles in these shells outside of the core $`{}_{38}{}^{}{}_{}{}^{88}`$Sr<sub>50</sub>. Therefore, we adopt <sup>88</sup>Sr as the inert core for the following shell model description of <sup>94</sup>Mo.
Single-particle energies $`\epsilon _j`$ were obtained from calculations for the neutron–odd nuclei <sup>89</sup>Sr, <sup>91</sup>Zr and <sup>93</sup>Mo and for the proton–odd nuclei <sup>89</sup>Y, <sup>91</sup>Y, <sup>91</sup>Nb, <sup>93</sup>Nb, and <sup>93</sup>Tc. These single-particle energies are close to those from . In order to get a rough estimate of the values of $`A_{T=1}^{nn}`$ and $`A_{T=1}^{pp}`$ parameters shell model calculations have been performed for the isobars <sup>90</sup>Sr and <sup>90</sup>Zr, which have either 2 neutrons or 2 protons, respectively, outside the core <sup>88</sup>Sr (see Fig.1). The results are $`A_{T=1}^{nn}=0.23`$ MeV and $`A_{T=1}^{pp}=0.35`$ MeV. The $`A_{T=1}^{pn}`$ parameter is taken as the average of $`A_{T=1}^{nn}`$ and $`A_{T=1}^{pp}`$, i.e. $`A_{T=1}^{pn}=(A_{T=1}^{nn}+A_{T=1}^{pp})/2`$. The value of 0.48 MeV for the $`A_{T=0}^{pn}`$ parameter was obtained by a fit to the excitation energies of the nucleus <sup>92</sup>Zr, which contains 2 protons and 2 neutrons outside the core <sup>88</sup>Sr. Experimental and calculated low-spin level schemes for the three even-even nuclei <sup>90</sup>Sr, <sup>90</sup>Zr, and <sup>92</sup>Zr are shown in Fig.1. The final values of the SDI parameters used for the description of <sup>94</sup>Mo were optimized by a fit to the low-spin level scheme of <sup>94</sup>Mo and are close to the values above. It is interesting to note that the excitation energies of low-lying states are much less sensitive to the $`A_{T=1}^{pn}`$ parameter than to the $`A_{T=0}^{pn}`$ parameter. However the electromagnetic transition strengths are very sensitive to both $`A_{T=1}^{pn}`$ and $`A_{T=0}^{pn}`$ parameters and the above discussed choice gives the best agreement between calculated and experimental strengths in <sup>94</sup>Mo nucleus. The single-particle energies and SDI parameters used for <sup>94</sup>Mo are presented in Table I. The SM calculations were performed on the Cologne Sun Ultra Enterprise 4000 workstation with two 166 MHz Ultra Sparc processors using the code RITSSCHIL .
## III Discussion
In the present shell model calculation for <sup>94</sup>Mo two neutrons can be distributed among five single particle orbitals: $`2d_{5/2}`$, $`3s_{1/2}`$, $`1g_{7/2}`$, $`2d_{3/2}`$ and $`1h_{11/2}`$. The lowest neutron orbital is $`2d_{5/2}`$. The fully occupied neutron orbital $`1g_{9/2}`$ forms the $`N=50`$ closed shell and is about 4 MeV below the neutron $`2d_{5/2}`$ orbital. Therefore, the influence of the $`1g_{9/2}`$ orbital on the structure of low–lying states of <sup>94</sup>Mo is expected to be small <sup>*</sup><sup>*</sup>*This expectation is supported by the first results of large scale shell model calculations for <sup>94</sup>Mo assuming the $`Z=N=40`$ nucleus <sup>80</sup>Zr as the inert core, which yield at the present stage almost negligible contributions of the $`\nu (g_{9/2})`$ orbital to low-lying low-spin states of <sup>94</sup>Mo .. This justifies the assumption of the $`N=50`$ neutron core. For protons we have included the two orbitals $`1g_{9/2}`$ and $`2p_{1/2}`$ in the configurational space. The closest higher–lying proton orbital to the proton $`1g_{9/2}`$ orbital appears more then 4 MeV above the Z=50 closed shell and is neglected. But the proton $`p_{1/2}`$ orbital is much closer to the proton $`1g_{9/2}`$ – about 1 MeV lower – and we have taken it into account by choosing $`Z=38`$ as the proton core. Within this configurational space we reproduce well many of the excited states in the spectrum of <sup>94</sup>Mo. This is shown in Fig. 2.
The main components of the low-lying states (see Table II) are seniority 2 and 4 with two protons in $`\pi (1g_{9/2})`$ and two neutrons in $`\nu (2d_{5/2})`$. But the influence of the $`\pi (2p_{1/2})`$ proton orbital is also significant - almost all states contain large components $`\pi (p_{1/2}^2g_{9/2}^4)_J`$ and this is the main component for the $`0_2^+`$ state. The contributions of the neutron orbitals $`\nu (3s_{1/2})`$, $`\nu (1g_{7/2})`$, $`\nu (2d_{3/2})`$ and $`\nu (1h_{11/2})`$ are smaller but for the $`1_2^+`$ the component $`\nu (g_{7/2}^1d_{5/2}^1)`$ is already the main one.
The $`M1`$ and $`E2`$ transition probabilities have been calculated and compared with the new experimental data. The results are shown in Tables III and IV. The reproduction of the data is in most cases very good. Tables III and IV include also IBM-2 predictions in the O(6) dynamical symmetry limit, where the states $`1_1^+,2_3^+,3_2^+`$ have mixed-symmetry. We note that the IBM-2 has only one free parameter - the effective quadrupole proton boson charge $`e_p`$, while the corresponding neutron charge was put to zero $`e_n=0`$. It is remarkable how well the shell model agrees also with the IBM-2 with only a few exceptions. The agreement of shell model calculations with IBM-2 results in different mass regions was noted also by other authors (see, for instance, ).
Let us discuss now the calculated $`M1`$ and $`E2`$ transition strengths in more detail. The MS assignments for the $`1_1^+,2_3^+,3_2^+`$ states of <sup>94</sup>Mo were based on the measurements of relatively large $`M1`$ transition strengths. For the calculation of $`M1`$ transitions between the shell model states we consider a nuclear magnetic dipole operator, which is the sum of proton and neutron one-body terms for orbital and spin contributions:
$$𝐓(M1)=\sqrt{\frac{3}{4\pi }}\left(\underset{i=1}{\overset{Z}{}}\left[g_p^l𝐥_i^p+g_p^s𝐬_i^p\right]+\underset{i=1}{\overset{N}{}}\left[g_n^l𝐥_i^n+g_n^s𝐬_i^n\right]\right)\mu _N,$$
(4)
where $`g_\rho ^l`$ and $`g_\rho ^s`$ are the orbital and spin g-factors and $`𝐥_i^\rho `$, $`𝐬_i^\rho `$ are the single particle orbital angular momentum operators and spin operators. For further discussion it is useful to decompose the $`M1`$ operator into an isoscalar part
$$𝐓_{IS}(M1)=\sqrt{\frac{3}{4\pi }}\left(g_J𝐉+g_S𝐒\right)\mu _N,$$
(5)
and an isovector part
$`𝐓_{IV}(M1)=\sqrt{{\displaystyle \frac{3}{4\pi }}}\left({\displaystyle \frac{g_p^lg_n^l}{2}}\left[𝐋_p𝐋_n\right]+{\displaystyle \frac{g_p^sg_n^s}{2}}\left[𝐒_p𝐒_n\right]\right).`$ (6)
where $`𝐋_\rho `$ and $`𝐒_\rho `$ denote the total orbital angular momentum and total spin operators for protons ($`\rho =p`$) and neutrons ($`\rho =n`$). $`𝐉=𝐋+𝐒`$ is the total angular momentum and does not generate $`M1`$ transitions. $`g_J=(g_p^l+g_n^l)/2=1/2`$ and $`g_S=[g_p^s+g_n^s(g_p^l+g_n^l)]/2=0.88\alpha _q1/2`$ with the quenching factor $`\alpha _q`$ defined by $`g_\rho ^s=\alpha _qg_\rho ^{s,\mathrm{free}}`$. The free spin $`g`$-factors are $`g_p^{s,\mathrm{free}}=5.58`$ and $`g_n^{s,\mathrm{free}}=3.82`$. Since $`g_p^s`$ and $`g_n^s`$ are of opposite sign and comparable, the isoscalar nondiagonal piece of an $`M1`$ matrix element is usually very small. It vanishes exactly for a quenching factor $`\alpha _q=0.57`$. For a good reproduction of the measured $`M1`$ transition strengths (see Table III) and for the sake of easy interpretation we used the quenching factor $`\alpha _q=0.57`$ which results in pure isovector $`M1`$ transitions.
The isovector $`M1`$ ground state excitation strength is calculated to be concentrated in the $`1_1^+`$ state. This agrees with the identification of the $`1_1^+`$ state with the scissors mode in the nucleus <sup>94</sup>Mo. This identification receives further support from the consistency of the data on <sup>94</sup>Mo with the systematics of the scissors mode extrapolated from the deformed rare earth nuclei. However, one cannot ascribe the collective scissor mode in the near-spherical nucleus <sup>94</sup>Mo as a pure orbital mode. Our calculations show that spin and orbital contributions are almost equal which agrees with the results of a single $`j`$ shell model .
The key signature for the lowest $`2_{\mathrm{ms}}^+`$ state, with a proton-neutron symmetry similar to that of the scissors mode, is a strong $`M1`$ transition to the $`2_1^+`$ state. Therefore, it is interesting to look to the $`2_1^+2_i^+`$ $`M1`$ strength distribution to judge the possible fragmentation of the $`2_{\mathrm{ms}}^+`$ state. ¿From the data it follows that the $`2_3^+2_1^+`$ is the strongest $`M1`$ transition from an excited $`2^+`$ state to the $`2_1^+`$ state which led to the MS assignment for the $`2_3^+`$ state. Comparing the calculated $`M1`$ strengths of the $`2_2^+2_1^+`$, $`2_3^+2_1^+`$ and $`2_4^+2_1^+`$ transitions one notes that the calculated $`B(M1;2_3^+2_1^+`$) value is about five times larger than the other two. The dominance of the $`B(M1;2_3^+2_1^+`$) value agrees with the data. The calculated $`B(M1;2_3^+2_1^+`$) value is also more than four times larger than the calculated $`B(M1;2_5^+2_1^+`$) value, which already overestimates the data. The shell model calculation agrees with the observation that the $`2_3^+2_1^+`$ transition concentrates the $`M1`$ strength between excited $`2^+`$ states and the $`2_1^+`$ state. This relatively strong isovector $`M1`$ transition agrees with the MS assignment for the $`2_3^+`$ state.
Also the experimental $`3_2^+`$ state decays by relatively strong $`M1`$ transitions to low-lying states. The calculated excitation energy of the $`3_2^+`$ state matches the experimental energy within 50 keV whereas the $`3_1^+`$ state lies 120 keV lower and the $`3_3^+`$ state lies 350 keV higher than the experimental $`3_{\mathrm{ms}}^+`$ state. Therefore we compare the $`3_2^+`$ shell model state with the observed $`3_2^+`$ state at 2965 keV. We note that the measured strong $`M1`$ transition from the $`3_2^+`$ state to the $`4_1^+`$ state is reproduced by the shell model within the experimental error bar. The shell model, however, underestimates the $`M1`$ strength of the $`3_2^+2_2^+`$ transition by a factor of two, while it overestimates the $`3_2^+2_1^+`$ $`M1`$ transition strength by an order of magnitude. The shell model results for the $`3_2^+`$ state disagree not only somewhat with the data but also with the prediction of the IBM in the O(6) dynamical symmetry limit for the $`3_{\mathrm{ms}}^+`$ state. However, the $`3_1^+`$ and $`3_2^+`$ states are close in energy, which renders the calculation of the wave functions more uncertain in the shell model, where no quantum number like the $`F`$-spin exists, which can assure the orthogonality of MS states to symmetric states. Moreover, the calculated $`3_3^+`$ state also shows $`M1`$ and $`E2`$ properties which are close to those of the experimental $`3_2^+`$ state. We conclude that the $`3_{\mathrm{ms}}^+`$ character is spread about the first three $`3^+`$ states in our shell model calculation with the surface delta residual interaction and the <sup>88</sup>Sr core. Thus for the $`3_2^+`$ state the MS assignment from the shell model results is less clear. We consider this fact not as an argument against the MS assignment of the experimental $`3_2^+`$ state of <sup>94</sup>Mo but as an indication of the limit of our present shell model approach for the description of the details (and the mixing) of wave functions for states, which lie close in energy. In total the $`M1`$ transition strengths calculated in the shell model support (or at least do not disagree with) the MS assignments for the $`1_1^+`$, $`2_3^+`$, and $`3_2^+`$ states of <sup>94</sup>Mo, if the existence of relatively strong isovector $`M1`$ transition is considered as a sufficient argument in favor of MS structures.
However, it should be stressed that the strongest M1 transition found in <sup>94</sup>Mo nucleus connects the $`4_2^+`$ and $`4_1^+`$ states that is in agreement with the present shell model calculations. This transition falls out from the $`sd`$-IBM-2 scheme and probably is related to the excitations of $`g`$-bosons in terms of the IBM.
Having found support for MS assignments from the large isovector $`M1`$ transition strengths it is interesting to turn to the $`E2`$ transition properties. While the existence of large isovector $`M1`$ transitions may indicate a different proton-neutron symmetry, one can, in contrast, judge a similar proton-neutron symmetry for two states from the existence of collective isoscalar $`E2`$ transition matrix elements between the two states. The $`E2`$ transition operator is the sum of proton and neutron parts:
$$𝐓(E2)=e_p𝐓_p(E2)+e_n𝐓_n(E2),$$
(7)
where $`e_p`$ and $`e_n`$ are the proton and the neutron effective quadrupole charges and $`𝐓_\rho (E2)=_i(r_i^\rho )^2𝐘_2(\theta _i^\rho ,\varphi _i^\rho )`$. It is again convenient to decompose the $`E2`$ operator in an isoscalar part
$$𝐓_{IS}(E2)=\frac{e_p+e_n}{2}\left[𝐓_p(E2)+𝐓_n(E2)\right],$$
(8)
and in an isovector part
$$𝐓_{IV}(E2)=\frac{e_pe_n}{2}\left[𝐓_p(E2)𝐓_n(E2)\right].$$
(9)
The calculated $`E2`$ transition strengths and matrix elements are compared to experiment and to schematic IBM-2 estimates in Tab. IV. In the low-lying low-spin level scheme of <sup>94</sup>Mo one expects from a simple quadrupole vibrator picture the existence of the collective isoscalar $`E2`$ transitions which are indicated in Fig. 3. The transition from the $`2_1^+`$ state to the $`0_1^+`$ ground state is a collective isoscalar E2 transition. Many components of the wave function of the $`2_1^+`$ state contribute coherently to the matrix element of the $`E2`$ transition operator. The isoscalar part of the $`2_1^+E20_1^+`$ matrix element is about ten times larger than the isovector part (see most right columns of Table IV). This is a well known general property of the lowest collective $`2^+`$ state. For <sup>94</sup>Mo the $`2_1^+`$ state is calculated to exhaust 97% of the total isoscalar $`E2`$ excitation strength of the ground state to the $`2^+`$ states up to 4 MeV. The $`4_1^+2_1^+`$ and $`2_2^+2_1^+`$ transitions are of similarly collective dominantly isoscalar $`E2`$ character. The strong collective isoscalar $`E2`$ transitions between the $`0_1^+`$, $`2_1^+`$, $`4_1^+`$ and $`2_2^+`$ states prove the proton-neutron symmetry of these states.
The E2 transitions from the $`2_2^+`$ state and the $`2_3^+`$ state to the ground state are much weaker than the $`2_1^+0_1^+`$ transition. The $`2_2^+0_1^+`$ transition has isovector character but it is weak. The $`2_3^+`$ state, however, carries 10$`\%`$ of the total $`E2`$ excitation strength of the $`2_1^+`$ state and it is the largest $`E2`$ excitation above the $`2_1^+`$ state. This supports the one-phonon character of the $`2_3^+`$ state, which is suggested from its interpretation as the one-quadrupole phonon $`2_{\mathrm{ms}}^+`$ state. Indeed, the transition from the $`2_3^+`$ state to the ground state is a mixture of isoscalar and isovector parts with a notable isovector component.
The $`E2`$ matrix elements between various states with MS assignments are very interesting. The calculated wave functions of the $`1_1^+`$ state and the $`3_2^+`$ state are dominated by basis states with seniority $`\nu =4`$ supporting their two-phonon interpretation. The $`1_1^+`$ state shows no collective $`E2`$ transitions to the $`2_1^+`$, $`2_2^+`$ states. The $`1_1^+2_3^+`$ transition is in contrast a collective isoscalar $`E2`$ transition which is comparable in strength with the $`2_1^+0_1^+`$ collective $`E2`$ transition. This indicates a similar proton-neutron symmetry of the $`1_1^+`$ state and the $`2_3^+`$ state and justifies to consider the $`1_1^+`$ state a two-phonon state formed by an isoscalar quadrupole phonon built on top of the MS $`2_3^+`$ state. The $`2_3^+`$ state has a proton-neutron symmetry similar to the $`1_1^+`$ state, but it is lower in energy and is of seniority $`\nu =2`$ like the collective lowest $`2_1^+`$ state. Also the calculated $`3_2^+`$ state decays by a strong collective isoscalar $`E2`$ transition to the $`2_3^+`$ state. The calculated $`3_2^+2_1^+`$ $`E2`$ transition strength is five times smaller than the $`3_2^+2_3^+`$ transition strength and the isovector part of the $`E2`$ matrix element is larger than the isoscalar part. Based on this comparison we can conclude that the $`3_2^+`$ state has qualitatively a similar proton-neutron symmetry as the $`1_1^+`$ state and the $`2_3^+`$ states. This conclusion supports the statement that the $`3_2^+`$ state of <sup>94</sup>Mo is one more representative of proton-neutron collective states with a “mixed-symmetry” character as it was argued before in Ref. . We note, that also the calculated $`3_2^+2_2^+`$ $`E2`$ transition has a large isoscalar part, which, however, overestimates the data. This disagreement may again be caused by a too large mixing of the calculated $`3_1^+`$ and $`3_2^+`$ states and is probably due to a too large symmetric three-phonon component in the calculated wave function of the $`3_2^+`$ state.
Of particular interest in this article is the proton-neutron structure of those states to which mixed-symmetry was previously assigned from the measurements of large $`M1`$ and $`E2`$ transition strengths . In order to simplify the discussion we consider now a schematic model for the analysis of the wave function of the $`2_{\mathrm{ms}}^+`$ state, which reflects the logic of the IBM-2. Let us assume that valence protons and neutrons couple separately to collective $`J_\rho ^\pi =2_p^+`$ and $`J_\rho ^\pi =2_n^+`$ configurations with seniorities $`\nu =2`$. In a two-level model one collective $`2_s^+`$ state is formed by symmetric linear combination of the $`2_p^+`$ proton and $`2_n^+`$ neutron configurations: $`|2_s^+>=(|2_p^+>+|2_n^+>)/\sqrt{2}`$. The orthogonal linear combination $`|2_{ns}^+>=(|2_p^+>|2_n^+>)/\sqrt{2}`$ has also seniority $`\nu =2`$ and is the collective “nonsymmetric” counterpart of the $`|2_s^+>`$ state. Furthermore, the $`2_p^+`$ proton and $`2_n^+`$ neutron configurations can couple to collective $`J^\pi =1^+,3^+`$ states with seniority $`\nu =4`$ and isovector character of decay to the $`2_s^+`$ state. The last three states should correspond to the lowest $`1^+,2^+,3^+`$ MS states in the IBM-2. The possible existence of such states in <sup>94</sup>Mo and in neighboring nuclei as a result of quadrupole surface vibrations in anti-phase and corresponding two-phonon excitations was discussed earlier in a geometrical approach by A. Faessler .
In the realistic shell model calculation for <sup>94</sup>Mo presented above MS states cannot so easily be identified from the wave functions. Isospin symmetry and seniority conservation are broken due to the interaction chosen and the single particle orbitals considered. The ground state contains 72% components with seniority $`\nu =0`$. The $`2_1^+`$ and $`2_3^+`$ states contain 70 $`\%`$ and 73 $`\%`$ components of seniority $`\nu =2`$. The relatively large components with seniority $`\nu =2`$ point at their predominantly one-quadrupole phonon nature . In contrast, the wave function of the $`2_2^+`$ state contains a large component of seniority $`\nu =4`$ (50% of the wave function) in agreement with its usual two-quadrupole phonon interpretation.
Let us now further analyze the calculated wave functions of the $`2_1^+`$ state and the $`2_3^+`$ state, which are considered to represent well the symmetric and the MS one-phonon $`2^+`$ states of the IBM-2. It is interesting that the components of the $`2_{1,3}^+`$ states with seniority $`\nu =2`$, $`|2_1^+,\nu =2`$ and $`|2_3^+,\nu =2`$, are approximately orthogonal. Their normalized scalar product
$$\frac{2_1^+,\nu =2|2_3^+,\nu =2}{\sqrt{2_1^+,\nu =2|2_1^+,\nu =2\times 2_3^+,\nu =2|2_3^+,\nu =2}}=0.07$$
is small. This fact is not a trivial consequence of the orthogonality of the $`|2_3^+`$ and $`|2_1^+`$ eigenstates, because their wave functions contain noticeable components with higher seniority. The seniority $`\nu `$=2 components resemble the schematic symmetric $`|2_s^+`$ and “nonsymmetric” $`|2_{ns}^+`$ states discussed above and can be used for further analysis. For this purpose it is interesting to decompose the seniority $`\nu `$=2 components of the $`2_1^+`$ and $`2_3^+`$ states into their proton and neutron parts. One obtains
$$|2_1^+,\nu =2=0.61|2_{1,p}^++0.80|2_{1,n}^+$$
(10)
and
$$|2_3^+,\nu =2=0.56|2_{3,p}^+0.82|2_{3,n}^+.$$
(11)
These decompositions cannot be directly compared to the schematic two level model mentioned above, because the basis states $`|2_{i,\rho }^+`$ are not identical. It turns out, however, that the normalized proton $`|2_{3,p}^+`$ basis state is rather similar to the normalized $`|2_{1,p}^+`$ basis state with a positive overlap of $`2_{3,p}^+|2_{1,p}^+`$=0.98. However, the overlapping of the neutron components of the states (10) and (11) is smaller. It amounts only to $`2_{3,n}^+|2_{1,n}^+`$=0.63 indicating considerable deviations from the pure IBM-2 picture. For a direct comparison between the structure of the $`|2_i^+,\nu =2`$ components it is more useful to express the components $`|2_{3,\rho }^+`$ in Eq. (11) by a linear combination of one part which is parallel to the $`|2_{1,\rho }^+`$ components and an orthogonal rest term $`|R`$ with $`R|2_{1,\rho }^+0`$. We obtain
$$|2_3^+,\nu =2=\gamma \left[0.72|2_{1,p}^+0.68|2_{1,n}^+\right]+|R$$
(12)
This result can be interpreted in the following way: The dominant seniority $`\nu =2`$ component of the $`2_3^+`$ state contains a fraction of $`\gamma ^2=58\%`$ of components that form the dominant seniority $`\nu =2`$ component of the $`2_1^+`$ state. Moreover, this fraction is almost orthogonal to the seniority $`\nu =2`$ component of the $`2_1^+`$ state, because the proton part and the neutron part contribute with a different sign while their amplitudes are almost equal. Based on this observation one can consider the $`|2_3^+`$ state as a good realization of the collective $`|2_{\mathrm{ms}}^+`$ MS state.
The schematic analysis given above helps to make some more general conclusions about the nuclear structure properties which can lead to the appearance of the “mixed symmetry” states in near-spherical nuclei. At first, we remind that the amplitudes of proton and neutron parts of symmetric and nonsymmetric states have to be approximately equal. These proton (neutron) parts of the wave functions of the symmetric state and its nonsymmetric counterpart must be rather similar, too. As it follows from our calculations this condition can be achieved if the configurational space for valence protons (neutrons) may be restricted to one high-j orbital (in our case proton $`1g_{9/2}`$) and one of the nearest orbitals with small single particle j quantum number (in our case proton $`2p_{1/2}`$). An ideal case is a single high-j orbital. Otherwise, if there are few neighboring neutron (proton) orbitals with the j value comparable to the single particle angular momentum of the selected leading neutron (proton) orbital and if the influence of these orbitals cannot be neglected, then the neutron (proton) parts of the symmetric and nonsymetric states can be rather different. We can observe it already in our case for the neutron parts of the $`2_1^+`$ state and the $`2_3^+`$ state: the overlapping $`2_{3,n}^+|2_{1,n}^+`$ is only 0.63. Therefore it can be expected that with the increase of the number of valence neutrons in the considered configurational space the neutron parts will be stronger fragmented and the neutron overlapping can be significantly reduced destroying the picture. On the contrary the increase of the number of valence protons in $`1g_{9/2}`$ orbital will keep the proton parts similar and the “mixed symmetry” $`2^+`$ state will probably survive. We hope that the above observations will be useful for the search for further mixed-symmetry phenomena in the neighboring nuclei.
We stress, however, that for the quantitative analysis of the relative proton-neutron symmetry of wave functions it can be more useful to analyze size and isotensor character of electromagnetic transition matrix elements between the calculated states, as was shown above, because of the presence of non-collective states in the shell model configurational space. Considering the radical truncation of the shell model problem which led to the formulation of the IBM it is remarkable how far the IBM and the shell model agree on the properties of mixed-symmetry states of <sup>94</sup>Mo.
## IV Conclusions
To summarize, we have performed shell model calculations for the near-spherical nucleus <sup>94</sup>Mo using the Surface Delta Interaction as the residual interaction. We calculated excitation energies of the low-spin positive-parity states and $`M1`$ and $`E2`$ transition strengths between them. In most cases the calculations agree well with the data. Calculated wave functions, $`M1`$, and $`E2`$ matrix elements support the previous mixed-symmetry assignments for the $`1_1^+`$ state, the $`2_3^+`$ state, and the $`3_2^+`$ state of <sup>94</sup>Mo. In particular, we find collective isoscalar $`E2`$ transitions between these three states and strong isovector $`M1`$ transitions to low-lying symmetric states. The strongest $`M1`$ transition is found between the $`4_2^+`$ state and the $`4_1^+`$ state. This transition is outside of the scope of the $`sd`$-IBM-2 approach. These findings indicate the common proton-neutron symmetry of the $`1_1^+,2_3^+,3_2^+`$ states showing that they form a class of states that differ to the lowest-lying ones by their proton-neutron structure. The analysis of the wave functions indicates in which neighboring nuclei these states can be most probably found.
## V Acknowledgments
We thank R.F. Casten, E. Caurier, A. Dewald, L. Esser, A. Gelberg, J. Ginocchio and T. Otsuka for discussions. R.V.J. thanks the Universität zu Köln for a Georg Simon Ohm guest professorship. This work was supported by the Deutsche Forschungsgemeinschaft under Contract Nos. Br 799/9-1 and Pi 393/1-1 and one of us (N.P.) got partly support by the US DOE under Contract No. DE-FG02-91ER-40609.
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# The symplectic and twistor geometry of the general isomonodromic deformation problem
### Introduction
The study of isomonodromic deformations of systems of ordinary differential equations in the complex plane was a significant topic at the beginning of the last century, when the classical work of Painlevé, Schlesinger, and Fuchs was published. It has come back into view in more recent years through connections with quantum field theory (Sato et al 1978, 1979, 1980, Dubrovin 1999), differential geometry (Hitchin 1996), and the theory of integrable systems (see, for example, Ablowitz and Clarkson 1991).
In this paper, I shall explore in detail one aspect of the modern theory, suggested by Hitchin (1995a). He considered the twistor space of a four-dimensional self-dual Riemannian manifold with $`\mathrm{SU}(2)`$ symmetry. This is a three-dimensional complex manifold $`𝒵`$, in which there is a four-dimensional family of projective lines corresponding to the points of the original manifold and on which the symmetry group acts holomorphically. The action is generated by three holomorphic vector fields which are independent in an open subset, but dependent on a special divisor $`S`$. By taking the vector fields as basis vectors, the tangent space to $`𝒵`$ is identified with $`\mathrm{sl}(2,)`$ at each point of the open subset. Thus the action determines a flat holomorphic connection on the trivial $`\mathrm{SL}(2,)`$ bundle over the open orbit.
The restriction of the connection to a twistor line is an $`\mathrm{sl}(2,)`$-valued meromorphic 1-form, with poles at the intersections with $`S`$. In the case that Hitchin considered, there are four poles and the 1-form determines a Fuchsian system, with four regular singularities. As the line is moved within the family, the poles move, but the monodromy of the system, which is the same as the holonomy of the flat connection, remains unchanged. By calling on the classical theory, therefore, one obtains from this geometrical picture a solution to the sixth Painlevé equation. Hitchin then goes on to exploit this correspondence to construct self-dual Einstein metrics from certain Painlevé transcendents.
Hitchin’s correspondence between twistor manifolds with symmetry and isomonodromic families of ordinary differential equations holds more generally. In this paper, I shall follow through the details of his suggestion for the class of isomonodromic deformations considered by Jimbo, Miwa, and Ueno (1981). This enables one to understand their results within the framework of the general deformation theory of Kodaira (1962).
In the general setting, we are given a complex Lie group $`G`$, a Riemann surface $`X`$, and a meromorphic 1-form $`\alpha `$ on $`X`$ with values in the Lie algebra $`𝔤`$. We pick a local coordinate $`z`$ and write $`\alpha =A\mathrm{d}z`$. Then the equation $`\mathrm{d}y+\alpha y=0`$ becomes a system of linear ordinary differential equations
$$\frac{\mathrm{d}y}{\mathrm{d}z}=Ay,$$
(1)
where $`y`$ is a fundamental solution, taking values in $`G`$, and $`A`$ is a meromorphic function of $`z`$. The first question concerns the existence of twistor spaces: this is answered by Proposition 1, which gives the existence of an embedding of $`X`$ in a complex manifold $`𝒵`$ on which $`𝔤`$ acts, with the generators independent on an open set, and from which $`\alpha `$ can be recovered by Hitchin’s construction. This structure is not unique, however, even if we restrict attention to a small neighbourhood of $`X`$ in $`𝒵`$. If $`A`$ has irregular singularities, then there are different choices for the way in which a divisor $`S`$ can be attached to the open set so that the whole of $`X`$ is embedded, including the singularities. Different choices give different possibilities for the normal bundle $`N`$ of $`X`$. By Kodaira’s theorem, the normal bundle determines the deformations of $`X`$ in $`𝒵`$: if
$$H^1(X,N)=0\text{and}\mathrm{dim}H^0(X,N)=d_0,$$
then $`X`$ is one of $`d_0`$-parameter family of compact curves $`X_t`$, on each of which the $`𝔤`$-action gives a linear system of differential equations. These are isomonodromic (Proposition 6). It is shown that there is a natural choice for the twistor space (the ‘full twistor space’ in Definition 2), for which the parameter space has the largest possible dimension in the generic case; a full twistor space exists generically (Proposition 2), and is unique in a neighbourhood of $`X`$ (Proposition 4). In the full case, $`N`$ can be constructed directly from $`\alpha `$; if $`H^1(X,N)=0`$, as is the case if $`X=_1`$, and generally if $`\alpha `$ has enough singularities, then every isomonodromic deformation arises from this construction (Proposition 8).
A second theme of this paper is the Hamiltonian nature of the isomonodromic deformation equations. A Fuchsian system on $`_1`$ is a system with regular singularities of the form
$$\frac{\mathrm{d}y}{\mathrm{d}z}=\frac{A_iy}{za_i},$$
where the residues $`A_i𝔤=\mathrm{sl}(2,)`$ are independent of $`z`$. Apart from gauge and coordinate transformations, the only possible deformations in the generic case are given by moving the poles $`a_i`$. The monodromy is then preserved if and only if
$$\frac{A_i}{a_j}=\frac{[A_i,A_j]}{a_ia_j},ij,\frac{A_i}{a_i}=\underset{ji}{}\frac{[A_i,A_j]}{a_ia_j},$$
(Schlesinger’s equations).<sup>1</sup><sup>1</sup>1One pole is fixed at infinity, and has residue $`A_i`$. Hitchin (1997) interpreted these as a Hamiltonian flow on the coadjoint orbits of the $`A_i`$s. This also generalises: when irregular singularities are present, the flows are on symplectic manifolds constructed from affine orbits in the loop algebra. The symplectic forms can be written down explicitly in terms of $`\alpha `$ and the Stokes’ matrices, and, at least in the case $`X=_1`$, one can also find explicit expressions for the Hamiltonians (Proposition 9). The symplectic structure is related to the structure of $`𝒵`$ in a neighbourhood of $`S`$.
An appendix outlines the theory of isomonodromic deformations for linear systems on a general Riemann surface.
Acknowledgements. This paper draws heavily on ideas on Nigel Hitchin as well as on Philip Boalch’s D.Phil. thesis (on which Boalch 2000 is based). I thank Majiec Dunajski and Guido Sanguinetti for many interesting discussions, and particularly Philip Boalch for his comments and for the corrections he made to the first draft.
### Twistor spaces
We suppose that we are given a complex Lie group $`G`$, a Riemann surface $`X`$, and a meromorphic 1-form $`\alpha `$ on $`X`$ with values in the Lie algebra $`𝔤`$. Our starting point is to interpret a solution $`y:X\{\text{poles}\}G`$ to the equation
$$\mathrm{d}y+\alpha y=0$$
(2)
as a complex curve in $`G`$ and to think of $`\alpha `$ as the pull-back of the Maurer-Cartan form on $`G`$—the $`𝔤`$-valued 1-form whose contraction with a left-invariant vector field is the corresponding element of $`𝔤`$.
In a local coordinate $`z`$ on $`X`$, (2) is a linear system of ordinary differential equations of the form
$$\frac{\mathrm{d}y}{\mathrm{d}z}=Ay,$$
where $`A`$ is meromorphic, with values in $`𝔤`$. Its poles are the singularities of the system—a pole of order $`r+1`$ is a singularity of rank $`r`$. Some familiar results about such systems are summarised in the appendix.
Of course $`y`$ is singular at the poles of $`\alpha `$, and it is multi-valued, so the embedding in $`G`$ is defined only locally. In the twistor picture, we seek to replace $`G`$ by a complex manifold $`𝒵`$ of the same dimension on which the right action of $`G`$ is retained and in which the singular points are included. In principle, the construction involves (i) taking a quotient by the monodromy group to make $`y`$ single-valued and (ii) attaching hypersurfaces on which the right-action of $`G`$ action is not free. The intersections of $`X`$ with the hypersurfaces will then correspond to the poles of the differential equation. Except for very special equations, however, the quotient is not Hausdorff. The best that we can do in general is to construct $`𝒵`$ as a neighbourhood of $`X`$, with the action of $`G`$ replaced by an action of its Lie algebra (which is enough to determine $`\alpha `$). The second step is generally straightforward in the regular case, but is more subtle when the linear system has irregular singularities.
In the context of the deformation problem, we shall adopt a special understanding of the meaning of the term ‘twistor space’.
###### Definition 1
A twistor space is a complex manifold $`𝒵`$ together with
> (i) a homomorphism from a complex Lie algebra $`𝔤`$ into the Lie algebra of holomorphic vector fields on $`𝒵`$; and
>
> (ii) a smooth compact complex curve $`X𝒵`$
such that the induced linear map $`\varphi _z:𝔤T_z𝒵`$ is an isomorphism for some $`zX`$.
Note that $`\mathrm{dim}𝒵=\mathrm{dim}𝔤`$. For the most part, we shall take $`𝔤=\mathrm{sl}(n,)`$, but other examples will also be considered.
Given a basis in $`𝔤`$, $`\mathrm{\Delta }=det\varphi `$ is a holomorphic section of $`^{\mathrm{dim}𝒵}T𝒵`$. We shall make the regularity assumptions that
$$S=\{\mathrm{\Delta }=0\}$$
is a complex hypersurface, that $`X`$ is transversal to $`S`$, that $`S`$ is the union of a finite set of components $`S_i`$, and that $`\mathrm{\Delta }`$ has a zero of order $`r_i+1`$ on $`S_i`$. These hold in all the twistor spaces constructed below.
At each $`z𝒵S`$, define $`\theta _zT_z^{}𝒵`$ by $`\theta _z=\varphi _z^1`$. Then $`\theta `$ is a holomorphic 1-form on $`𝒵S`$ with values in $`𝔤`$. It is meromorphic on $`𝒵`$ and satisfies the Maurer-Cartan equation
$$\mathrm{d}\theta +\theta \theta =0$$
(3)
Equivalently, $`\mathrm{d}+\theta `$ is a flat meromorphic connection on a trivial bundle over $`𝒵`$. The restriction $`\alpha =\theta |_X`$ determines a system of the form (2), with poles at the intersection points with $`S`$. We then say that $`(𝒵,X)`$ is a twistor space for the system.
Example. Let $`G=\mathrm{SL}(n,)`$ and let $`𝔱𝔤=\mathrm{sl}(n,)`$ denote the diagonal subalgebra. As an $`(n1)`$-dimensional additive group, $`𝔱`$ acts on itself by translation, and the action extends to the compactification $`_{n1}`$ when we add a hyperplane at infinity. We also have the left action of $`𝔱`$ on $`G`$, defined by
$$g\mathrm{exp}(A)g,A𝔱.$$
We put
$$𝒵=G\times _{n1}/𝔱.$$
Then the right action of $`G`$ on the first factor descends to the quotient.
We can think of $`𝒵`$ as being formed by attaching a single hypersurface $`S`$ to $`G`$ (the projection of the hyperplane at infinity in $`_{n1}`$). The effect is to compactify the one-parameter subgroups generated by the semisimple elements of $`𝔤`$. If $`A𝔱`$ generates a closed subgroup, then $`\{tA\}𝔱`$ compactifies to a projective line in $`_{n1}`$, and this in turn projects onto an embedded copy of $`_1`$ in $`𝒵`$. The corresponding system of linear equations is
$$\frac{\mathrm{d}y}{\mathrm{d}z}=Ay,$$
which has a singularity of rank $`1`$, a double pole, at infinity (the intersection with $`S`$).
Example. Suppose that $`X`$ has genus $`g`$. Let $`G=𝒵`$ be the Jacobi variety (an abelian group with Lie algebra $`𝔤=^g`$) and let $`X𝒵`$ be the standard embedding (see, for example, Farkas and Kra 1980, p. 87). The corresponding system is
$$\frac{\mathrm{d}y}{\mathrm{d}z}=A$$
where $`A=(\xi _1,\mathrm{},\xi _g)`$, with the $`\xi `$s a basis for the abelian differentials on $`X`$. Here there are no singularities.
#### Existence of twistor spaces
Does every system of ODEs of the form (2) have a twistor space? Since the restriction of $`\theta `$ to $`X`$ cannot vanish, a necessary condition is that $`\alpha 0`$ at every point of $`X`$. This condition holds in the generic case (since it fails only if the all the entries in the matrix $`A`$ have a coincident zero). It is also sufficient.
###### Proposition 1
Let $`\alpha `$ be a meromorphic $`𝔤`$-valued 1-form on $`X`$ with no zeros. Then the linear system of ODEs $`\mathrm{d}y+\alpha y=0`$ has a twistor space.
Proof. We construct $`𝒵`$ by taking a quotient of a neighbourhood of the identity section $`X`$ in $`X\times G`$ by a distribution $`F`$ constructed from the linear system.
Let $`D`$ be a neighbourhood of a pole $`a`$ not containing any of the other poles, and let $`z`$ be a coordinate on $`D`$ such that $`z=0`$ at $`a`$. Then $`\alpha =A\mathrm{d}z`$ in $`D`$, where $`A:D\{0\}𝔤`$ is holomorphic and has a pole of order $`r+1`$ at $`z=0`$.
Define $`F`$ to be the distribution on $`D\times G`$ tangent to the non-vanishing vector field
$$z^{r+1}(_zR_A)$$
where $`R_A`$ is the right-invariant vector field on $`G`$ generated by $`A(z)`$. If $`D^{}`$ is an open set not containing any other poles, then we define $`F`$ in the same way on $`D^{}\times G`$, but without the factor $`z^{r+1}`$; that is $`F`$ is tangent to $`_zR_A`$. The vector fields are proportional on $`DD^{}`$, so $`F`$ is well defined globally as a distribution on $`X\times G`$. Under the condition on $`\alpha `$, we have $`FT_xX=0`$ at every $`xX`$. So it is possible to choose an open neighbourhood $`N`$ of $`X`$ in $`X\times G`$ such that the quotient $`𝒵=N/F`$ is a Hausdorff complex manifold of the same dimension as $`G`$. We then have a double fibration
and a smooth curve $`\pi _2(X)𝒵`$, which we also denote by $`X`$.
Because we are looking only at a neighbourhood of the identity section, the right action of $`G`$ on $`X\times G`$ does not pass to $`𝒵`$; but the corresponding Lie algebra action does. Each $`v𝔤`$ can be identified with a left-invariant vector field on $`G`$, and hence with a vector field on $`X\times G`$ tangent to the fibres of $`\pi _1`$. Its projection $`V`$ by $`\pi _2`$ is a holomorphic vector field on $`𝒵`$, and the map $`vV`$ is a Lie algebra representation, satisfying the conditions in the definition of a twistor space. The singular hypersurface has components given by the poles of $`\alpha `$, and $`X`$ meets these transversally.
It remains to show that $`\theta |_X=\alpha `$. To do this, we note that the meromorphic vector field $`_zR_A`$ on $`_1\times G`$ is tangent to $`F`$, and so its projection into $`𝒵`$ vanishes. On the other hand, at the identity, the right- and left-invariant vector fields generated by an element of $`𝔤`$ coincide. Hence $`i__z^{\text{}}\theta =A`$. The proposition follows.
Remark. If we instead take the curve in $`𝒵`$ to be the projection under $`\pi _2`$ of $`X\times \{g\}`$ for some other constant element of $`g`$, then we obtain instead a twistor space for $`g^1\alpha g`$.
In the irregular case, the twistor space is not unique: the one that arises in Proposition 1 is minimal in a sense that will be explained later.
#### Full twistor spaces
The difference in structure between different twistor spaces of a system of ODEs can be understood by looking at the structure in a neighbourhood of a point of $`aS`$. By introducing a local coordinate $`z`$ that vanishes on $`S`$, we can choose a neighbourhood $`U`$ of the form $`S\times D`$, where $`D`$ is, say, the unit disc, and the $`XU`$ is $`\{(a,z)\}`$, $`zD`$.
Suppose that $`\alpha =A\mathrm{d}z`$ has a pole of order $`r+1`$ at $`z=0`$. Then corresponding system
$$\frac{\mathrm{d}y}{\mathrm{d}z}=Ay$$
(4)
has a singularity of rank $`r`$ at $`z=0`$.
Given a holomorphic vector field $`V`$ on $`𝒵`$ tangent to $`S`$, we can construct a holomorphic family of copies $`D_t`$ of $`D`$ in $`U`$ by moving $`D_0=XU`$ along $`V`$ (and if necessary restricting to a smaller neighbourhood of $`a`$): here $`t`$ is a complex parameter taking values in some neighbourhood of $`t=0`$. By restricting $`\theta `$ to each $`X_t`$, we get a one-parameter family of ODEs
$$\frac{\mathrm{d}y}{\mathrm{d}z}=A(z,t)y,$$
each with a singularity of rank $`r`$ at $`z=0`$ (the singularity does not move with $`t`$ because $`V`$ is tangent to $`S`$). It follows from the flatness of $`\mathrm{d}+\theta `$ that
$$\frac{A}{t}=\frac{\mathrm{\Omega }}{z}[A,\mathrm{\Omega }],$$
where $`\mathrm{\Omega }=i_V^{\text{}}\theta `$. This is the local deformation equation (see appendix).
At each fixed value of $`t`$, $`\mathrm{\Omega }`$ is a function of the coordinate $`z`$. Introducing the notation
$$\mathrm{\Omega }=_z\mathrm{\Omega }[A,\mathrm{\Omega }]$$
we have
$$\mathrm{\Omega }=O(z^{r1}),\mathrm{\Omega }=O(z^{r1})\text{as }z0\text{.}$$
(5)
When the singularity is irregular, the various twistor spaces differ in the extent to which the converse holds: in the minimal construction, an $`\mathrm{\Omega }`$ satisfying these conditions is of the form $`i_V^{\text{}}\theta `$ for some holomorphic $`V`$ only if $`\mathrm{\Omega }f(z)A`$ is holomorphic at $`z=0`$ for some holomorphic function $`f`$.
###### Definition 2
A twistor space $`𝒵`$ is full at $`aS`$ if for every $`\mathrm{\Omega }:D𝔤`$ such that (5) holds, there is a holomorphic vector field $`V`$ on $`U𝒵`$ such that $`\mathrm{\Omega }=i_V^{\text{}}\theta |_X`$. The twistor space is full if it is full at every point of $`S`$.
When the system has irregular singularities, and Rank$`(G)>2`$, the twistor space constructed in Proposition (1) is not full. We can see this by noting that, for any holomorphic $`V`$, $`i_V\theta `$ has singular part at each pole that is proportional (by a holomorphic function) to a multiple of $`A`$, and cannot therefore give rise to the most general $`\mathrm{\Omega }`$ satisfying (5). We shall put this more precisely below when we consider the normal bundle of $`X`$ in $`𝒵`$.
A full twistor space generates not only the ODE itself, but also its isomonodromic deformations. We shall see that it is possible to construct a full twistor space in the generic case, but there are some rather special exceptions. A necessary condition for existence is that if $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }^{}`$ both satisfy (5), and if
$$\mathrm{\Omega }=\frac{M}{z^{r+1}}+O(z^r),\mathrm{\Omega }^{}=\frac{M^{}}{z^{r+1}}+O(z^r),$$
as $`z0`$, with $`k>0`$, then $`[M,M^{}]=0`$. This fails in the following class of examples.
Example. Suppose that
$$A=z^2\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),$$
Then
$$\mathrm{\Omega }=\frac{M}{z},\mathrm{\Omega }^{}=\frac{M^{}}{z},$$
satisfy (5) for any constant $`M`$, $`M^{}`$; but both cannot be generated from vector fields in the same twistor space if $`[M,M^{}]0`$.
#### The singular hypersurface
We now look in detail at the structure of a twistor space in a neighbourhood of $`S`$. The fullness condition at a point of $`XS`$ is this: given $`\mathrm{\Omega }`$ such that (5) holds, is $`\varphi (\mathrm{\Omega })`$ holomorphic at $`z=0`$?
Let us write
$$A=\frac{p}{z^{r+1}}+H,\mathrm{\Omega }=\frac{\omega }{z^{r+1}}$$
(6)
where
$$p=p_0+p_1z+\mathrm{}+p_rz^r,\omega =\underset{0}{\overset{\mathrm{}}{}}\omega _iz^i$$
and $`H`$ is holomorphic on $`U`$. If $`\omega =O(z^{r+1})`$ as $`z0`$, then $`\mathrm{\Omega }`$ is holomorphic at $`z=0`$ and can be generated by the holomorphic vector field $`\varphi (\mathrm{\Omega })`$ in any twistor space for $`A`$; so the source of any difficulty lies in the first $`r`$ terms in the Taylor expansion of $`\omega `$.
When we separate out the coefficients of $`z^{2(r+1)}`$, $`z^{2r1}`$, …$`z^{(r+2)}`$ in (5), we obtain for $`r>0`$
$`[p_0,\omega _0]`$ $`=`$ $`0`$
$`[p_0,\omega _1]+[p_1,\omega _0]`$ $`=`$ $`0`$
$`\mathrm{}`$
$`[p_0,\omega _{r1}]+[p_1,\omega _{r2}]+\mathrm{}+[p_{r1},\omega _0]`$ $`=`$ $`0`$
$`[p_0,\omega _r]+[p_1,\omega _{r1}]+\mathrm{}+[p_r,\omega _0](r+1)\omega _0`$ $`=`$ $`0,`$
or in the case $`r=0`$,
$$[p_0,\omega _0]=\omega _0.$$
(7)
The generic case is that one or other of the following hold:
> (i) $`r>0`$ and the eigenvalues of $`p_0`$ are distinct; or
>
> (ii) $`r=0`$, the eigenvalues of $`p_0`$ are distinct, and no pair differ by an integer.<sup>2</sup><sup>2</sup>2To prove Lemma 1 and Proposition 2, we only need that no pair should differ by 1; however the stronger condition here is needed to construct $`g_\mathrm{f}^{\text{}}`$ (see Appendix), and is imposed here to avoid special cases in the presentation below.
###### Lemma 1
If $`A`$ is generic, then $`\mathrm{\Omega }`$ satisfies (5) if and only if $`\omega _0=0`$ and $`[\omega ,p]=O(z^{r+1})`$ as $`z0`$.
Proof. Under either of the conditions (i), (ii), the eigenvalues of $`p`$ can be assumed to be distinct, since they are distinct at $`z=0`$ and since we can, if necessary, replace $`U`$ by a smaller neighbourhood. So we can find a holomorphic gauge transformation $`g:UG`$ such that $`g^1pg`$ is the sum of a diagonal polynomial and a term that vanishes to order $`z^{r+1}`$ at $`z=0`$, and so can be absorbed into $`r`$.
If we assume first that $`p`$ is actually diagonal, then we deduce successively that $`\omega _0,\mathrm{}\omega _{r1}`$ are diagonal (for $`r>0`$) and that
$$[p_0,\omega _r]=(r+1)\omega _0.$$
For $`r>0`$, this implies that $`\omega _0=0`$ since the diagonal terms on the left-hand side vanish, and hence that $`\omega _r`$ is also diagonal. For $`r=0`$, it gives $`\omega _0=0`$ since $`p_0`$ has no pair of eigenvalues differing by 1. Thus, whether or not $`p`$ is diagonal, we have that $`\mathrm{\Omega }`$ is holomorphic at $`z=0`$ when $`r=0`$; and that when $`r>0`$,
$$\mathrm{\Omega }=\frac{gqg^1}{z^r}+O(z^0)$$
where $`q`$ is a diagonal polynomial of degree $`r1`$.
###### Proposition 2
Let $`\mathrm{d}y+\alpha y=0`$ be a generic system, with $`G=\mathrm{SL}(n,)`$. Then there exists a full twistor space.
Proof. Any twistor space is full at a singularity of rank $`r=0`$ since any $`\mathrm{\Omega }`$ satisfying (5) is then holomorphic at $`z=0`$, and can therefore be generated by a holomorphic vector field in any twistor space. In the irregular case, we construct $`𝒵`$ from the ‘minimal’ twistor space in Proposition 1, by cutting out and replacing a neighbourhood of each component of $`S`$ corresponding to an irregular singularity.
Suppose, to begin with, that the system has a singularity of rank $`r`$ at $`z=0`$ and that in a neighbourhood $`D`$ of $`z=0`$ we have $`\alpha =A\mathrm{d}z`$, where
$$A=\frac{p}{z^{r+1}}+H$$
(8)
with $`p`$ a diagonal polynomial of degree $`r`$ with distinct diagonal entries throughout $`D`$ and $`H`$ holomorphic. By making a diagonal gauge transformation, we can make $`H`$ off-diagonal.
Pick constant diagonal matrices $`q_1,\mathrm{},q_{n2}`$ which, together with $`p(0)`$, form a basis for the diagonal subalgebra of $`𝔤`$, and for each $`i`$ let $`H_i`$ be the off-diagonal matrix with entries
$$(H_i)_{ab}=\frac{zH_{ab}(q_{ia}q_{ib})}{p_ap_b},ab,$$
where $`q_{ia}`$ and $`p_a`$, $`a=1,\mathrm{},n`$, are the diagonal entries in $`q_i`$ and $`p`$. Thus
$$[p,H_i]=[q_i,zH],[q_i,H_j]=[q_j,H_i].$$
(9)
Now introduce evolution equations for the diagonal matrix $`p(z)`$ and the off-diagonal matrix $`H(z)`$ as functions of the complex variables $`t_1,\mathrm{},t_{n1}`$ by putting
$$_ip=rq_i_iH=_zH_i[H,H_i],$$
where $`_i=/t_i`$, $`_z=/z`$. The integrability of this system is established by showing that
$$_iH_j_jH_i=[H_i,H_j].$$
Since both sides are off-diagonal, this follows from (9) and
$`[p,_iH_j][p,_jH_i]`$ $`=`$ $`z[q_j,_iH]z[q_i,_jH]`$
$`=`$ $`z[q_j,_zH_i]z[q_i,_zH_j]z[q_j,[H,H_i]]+z[q_i,[H,H_j]]`$
$`=`$ $`z[H_i,[q_j,H]]z[H_j,[q_i,H]]`$
$`=`$ $`[p,[H_i,H_j]].`$
So the evolution equations extend $`H`$ and $`p`$ to functions of $`(z,t_1,\mathrm{}t_{n2})`$ on a neighbourhood $`W`$ of the origin in $`^{n2}`$.
It follows from the definitions that
$$=\mathrm{d}\frac{p\mathrm{d}z}{z^{r+1}}H\mathrm{d}z\underset{i}{}\left(\frac{q_i\mathrm{d}t_i}{z^r}+H_i\mathrm{d}t_i\right)$$
is flat meromorphic connection on the trivial bundle principal bundle $`P=G\times W`$.
Let $`𝒬`$ denote the quotient of a neighbourhood of the identity section in $`P`$ by the horizontal foliation. The foliation extends holomorphically to $`z=0`$ since it is spanned by
$$z^{r+1}_zpz^{r+1}H,z^r_iq_iz^rH_i(i=1,\mathrm{},n2),$$
where $`p,q_i,H,H_i`$ are interpreted as right-invariant vector fields on $`G`$.
The quotient is a ‘local twistor space’ for $`A`$ in the sense that it carries a holomorphic $`𝔤`$-action, which is free and transitive except on the hypersurface $`S^{}=\{z=0\}`$, and contains a copy of $`D`$ on which the induced system is $`\alpha `$. Moreover, the fullness condition holds at $`S`$ since $`p_0`$ and the $`q_i`$s span the diagonal subalgebra (in the case $`n=2`$, $`p_0`$ on its own does that). Any generic $`A`$ can be reduced to the form (8) by a holomorphic gauge transformation $`g(z)`$; so more generally a local twistor space can be constructed by applying the same gauge transformation to $``$.
By using the $`𝔤`$ action, we can identify $`𝒬S^{}`$ with $`VS`$, where $`V`$ is a neighbourhood in $`𝒵`$ of the $`z=0`$ intersection point of $`S`$ and $`X`$. Then the embedded copy of $`D\{0\}`$ is mapped onto a punctured neighbourhood of the singularity in $`X`$. The identification allows us to replace $`V`$ by $`𝒬`$. By repeating this for the other irregular singularities, we obtain a full twistor space.
Given the choice of coordinate $`z`$ in a neighbourhood of an irregular singularity, $`p_0`$ is a well-defined map $`S𝔤`$. Up to scale, $`p_0`$ is independent of the choice of $`z`$. We thus have a natural map $`[p_0]:S𝔤`$. It is equivariant with respect to the $`𝔤`$ action on $`S`$ and the adjoint action on $`𝔤`$. The following is immediate from the proof above.
###### Proposition 3
A twistor space for a generic system is full at an irregular singularity if and only if $`[p_0]:S𝔤`$ is regular at $`XS`$.
When the space is not full, $`[p_0](S)`$ has nonzero codimension in $`(𝔤)`$. In the generic case, the full twistor space is locally unique in the sense of the following proposition.
###### Proposition 4
Suppose that $`(𝒵,X)`$ and $`(𝒵^{},X)`$ are full twistor spaces for a generic linear system of ODEs on a Riemann surface $`X`$, with $`G=\mathrm{SL}(n,)`$. Then there are neighbourhoods $`UX`$ and $`U^{}X`$ in $`𝒵`$ and $`𝒵^{}`$ and a $`𝔤`$-equivariant biholomorphic map $`\rho :UU^{}`$ such that $`\rho (X)=X`$.
If we exclude the poles from $`X`$ and the corresponding hypersurfaces $`S`$ and $`S^{}`$ from $`𝒵`$ and $`𝒵^{}`$, then $`\rho `$ is determined in a straightforward way by the $`𝔤`$ actions on $`𝒵`$ and $`𝒵^{}`$. It is defined by choosing a (multivalued) solution $`y`$ to the system on $`X`$ and then extending $`y`$ to a (multivalued) map from a neighbourhood of $`XS`$ in $`𝒵S`$ to $`G`$ such that $`\mathrm{d}y+\theta y=0`$. Similarly $`y`$ extends to $`y^{}`$ on $`𝒵^{}`$. Then the required map is $`\rho =y^1y`$, where the branches are chosen so that $`\rho `$ is the identity on $`X`$ ($`\rho `$ is well defined since $`y`$ and $`y^{}`$ have the same holonomy). The fact that $`\rho `$ extends holomorphically to $`S`$ in the full case is a corollary of Proposition 6 below.
#### The structure of S
Suppose that $`G=\mathrm{SL}(n,)`$ and that $`(𝒵,X)`$ is full. We denote by $`\mathrm{\Gamma }`$ the space of parametrized curves
$$D𝒵:z\gamma (z),$$
where $`D`$ is the unit disc, $`\gamma (D)`$ meets some component of $`S`$ transversally at $`z=0`$, and $`\gamma `$ extends smoothly to $`|z|=1`$.
We shall now develop a picture in which $`\theta _\gamma =\gamma ^{}(\theta )`$ is regarded as an element of the dual of Lie algebra of the loop group $`LG`$.<sup>3</sup><sup>3</sup>3An element of loop algebra is a map $`B:S^1𝔤`$. We define $`\theta _\gamma ,B`$ by integrating $`\mathrm{tr}(B\theta _\gamma )`$ around the unit circle. Different elements of $`\mathrm{\Gamma }`$ give different points of an orbit in $`L𝔤^{}`$ of an affine action of $`LG`$. We shall construct a finite-dimensional complex symplectic manifold from the orbit which determines the singular part of $`\alpha `$ at $`z=0`$.
#### The Fuchsian case
If the singularity at $`z=0`$ is regular and generic, then any twistor space is full at $`z=0`$ and any holomorphic map $`z\mathrm{\Omega }(z)𝔤`$ generates a local holomorphic vector field tangent $`Y`$ to $`S`$. Moreover we can write
$$\theta =p_0\mathrm{d}(\mathrm{log}z)+\theta ^{}$$
where $`z=0`$ on $`S`$ and $`\theta ^{}`$ is holomorphic on $`S`$.<sup>4</sup><sup>4</sup>4Thus $`\theta `$ has a logarithmic pole in the sense of Malgrange (1982). Then $`p_0`$ is independent of the choice of the function $`z`$, and therefore determines a natural map $`\mu :S𝔤=𝔤^{}`$—the identification being given by the bilinear form $`\mathrm{tr}(\xi _1,\xi _2)`$, $`\xi _i𝔤`$. The image $`\mu (S)`$ is open subset of a coadjoint orbit.
By evaluating $`\mathrm{\Omega }`$ and $`Y`$ at $`z=0`$, we obtain a natural identification
$$T_aS=𝔤/[p_0(a)]aS.$$
We can therefore define a 2-form $`\sigma `$ on $`S`$ by
$$\sigma _a(Y,Y^{})=\mathrm{tr}(p_0[\mathrm{\Omega },\mathrm{\Omega }^{}]).$$
This is closed and presymplectic since it is the pull back to $`S`$ by $`\mu `$ of the symplectic form on $`\mu (S)`$.
#### The irregular case
In the irregular case, the analogous structure involves information from the higher formal neighbourhoods of $`S`$. It arises from the action on $`\mathrm{\Gamma }`$ of the group $`L_+G`$ of holomorphic maps $`g:DG`$ that extend smoothly to $`\overline{D}`$: if $`gL_+G`$ then $`\gamma \mathrm{\Gamma }`$, then $`(g\gamma )(z)=\gamma (z)g(z)`$.<sup>5</sup><sup>5</sup>5Of course this is well-defined only for $`g`$ close to the identity; what follows has to be qualified in a similar way.
A tangent vector $`Y`$ to $`\mathrm{\Gamma }`$ at $`[\gamma ]`$ is a section of $`T𝒵|_\gamma `$, tangent to $`S`$ at $`z=0`$. Put
$$\sigma _\gamma (Y,Y^{})=\frac{1}{2\pi \mathrm{i}}\mathrm{tr}(\mathrm{\Omega }\mathrm{\Omega }^{}),$$
(10)
where $`\mathrm{\Omega }=i_Y\theta `$, $`\mathrm{\Omega }^{}=i_Y^{}\theta `$, $`\mathrm{\Omega }=\mathrm{d}\mathrm{\Omega }+[\theta ,\mathrm{\Omega }]`$ and the integral is along a loop surrounding $`z=0`$. This form is closed, but degenerate (its closure follows from the construction below). Its characteristic distribution is integrable, by closure, and contains the $`Y`$s for which $`\mathrm{\Omega }=O(z^{r+1})`$ as $`z0`$. These span the orbits of the normal subgroup $`L_{r+1}GL_+G`$ of maps $`g:DG`$ such that $`g=1+O(z^{r+1})`$ as $`z0`$. Since $`\mathrm{\Gamma }/L_{r+1}G`$ is finite-dimensional, the quotient $`\mathrm{\Gamma }_r`$ of $`\mathrm{\Gamma }`$ by the characteristic distribition is a finite-dimensional symplectic manifold.
Since $`𝒵`$ is full, the tangent space to $`\mathrm{\Gamma }_r`$ at $`[\gamma ]`$ is the set of holomorphic maps
$$\mathrm{\Omega }:D\{0\}𝔤$$
such that (5) holds, modulo maps with zeros of order $`r`$ at $`z=0`$.
#### Affine orbits
Let $`LG`$ denote the loop group of smooth maps $`S^1G`$ (Pressley and Segal, 1986). Its Lie algebra is the space $`L𝔤`$ of smooth maps $`\mathrm{\Omega }:S^1𝔤`$.
A $`1`$-form $`\alpha `$ on $`S^1`$ with values in $`𝔤`$ determines an element of $`L𝔤^{}`$ by
$$\alpha ,\mathrm{\Omega }=\frac{1}{2\pi \mathrm{i}}\mathrm{tr}(\mathrm{\Omega }\alpha )$$
Let $`𝒜_rL𝔤^{}`$ denote the subspace of smooth $`\alpha `$s that extend meromorphically to $`D`$ with a pole at the origin of order at most $`r+1`$, and no other poles. Thus
$$L𝔤^{}𝒜_r𝒜_{r1}\mathrm{}𝒜_0.$$
Like any dual Lie algebra, $`(L𝔤)^{}`$ carries the standard Kostant-Kirillov-Souriau Poisson structure, which is preserved by the coadjoint action of $`LG`$. The natural symplectic arena for the isomonodromy problem is not, however, that of the corresponding coadjoint orbits, but rather that of the orbits of the affine action of $`L𝔤`$ on $`L𝔤^{}`$ given by the cocycle
$$c(\mathrm{\Omega },\mathrm{\Omega }^{})=\frac{1}{2\pi \mathrm{i}}\mathrm{tr}(\mathrm{\Omega }\mathrm{d}\mathrm{\Omega }^{}).$$
Each $`\mathrm{\Omega }L𝔤`$ determines a vector field on $`L𝔤^{}`$, its value at $`AL𝔤^{}`$ being given by
$$\delta \alpha ,=\alpha ,[\mathrm{\Omega },]c(\mathrm{\Omega },).$$
The first term on the right-hand side is the usual coadjoint action; the second is a translation introduced by Souriau (1970) (see also Woodhouse 1990).<sup>6</sup><sup>6</sup>6For a general Lie group, the affine action is symplectic, but not Hamiltonian when $`c`$ is not a coboundary. In fact the affine orbits are the models for the non-Hamiltonian transitive symplectic actions of the group in the same way that the coadjoint orbits are the models for the Hamiltonian actions. If $`c`$ is the obstruction to the existence of a moment map for some transitive symplectic action of the Lie algebra on a symplectic manifold $`M`$, then $`M`$ can be mapped equivariantly and symplectically to an affine orbit in the dual Lie algebra. The affine orbits have an alternative description in terms of the coadjoint orbits of the central extension determined by $`c`$, but this is less convenient for our purposes here. See Pressley and Segal (1986), p. 44. By integrating the flow on $`L𝔤^{}`$, we obtain the gauge action of $`LG`$:
$$\alpha g^1\alpha g+g^1\mathrm{d}g.$$
The symplectic structure on the corresponding orbits is<sup>7</sup><sup>7</sup>7 This is a good definition of $`\sigma `$ since the right-hand side vanishes whenever $`\mathrm{\Omega }_1`$ fixes $`A`$, and therefore whenever $`Y_1`$ vanishes at $`A`$. We shall not need to consider the precise sense in which $`\sigma `$ is nondegenerate, since we shall deal only with finite-dimensional submanifolds.
$$\sigma (Y,Y^{})=c(\mathrm{\Omega },\mathrm{\Omega }^{})\alpha ,[\mathrm{\Omega },\mathrm{\Omega }^{}]=\frac{1}{2\pi \mathrm{i}}\mathrm{tr}\left(\mathrm{\Omega }\mathrm{\Omega }^{}\right)=\frac{1}{2\pi \mathrm{i}}\mathrm{tr}\left(\mathrm{\Omega }\delta ^{}\alpha \right),$$
(11)
where $`Y,Y^{}`$ are the vector fields generating the actions of $`\mathrm{\Omega },\mathrm{\Omega }^{}L𝔤`$,
$$\mathrm{\Omega }=\mathrm{d}\mathrm{\Omega }+[\alpha ,\mathrm{\Omega }],$$
and $`\delta ^{}\alpha `$ is the variation induced by $`Y^{}`$. The flow of $`\mathrm{\Omega }L𝔤`$ is generated by
$$h(\alpha )=\frac{1}{2\pi \mathrm{i}}\mathrm{tr}(\alpha \mathrm{\Omega }).$$
However, $`[h_A,h_B]=h_{[A,B]}+c(A,B)`$, so the action is not Hamiltonian: the cocycle is the obstruction to the existence of a moment map.
#### The symplectic structure of $`M_r`$
Let $`𝒪`$ be the affine orbit of some generic element of $`𝒜_r`$; that is, an element of the form
$$g^1\mathrm{d}gg^1\left(\frac{p}{z^{r+1}}+H\right)g,$$
where $`gL_+G`$, $`HL_+𝔤`$ are holomorphic on $`D`$, and $`p`$ is a polynomial in $`z`$ with distinct eigenvalues for $`zD`$.
A general element $`\alpha 𝒪`$ is a smooth 1-form on $`S^1`$ with values in $`𝔤`$. For any $`\alpha ,\alpha ^{}𝒪`$, the two systems
$$\mathrm{d}y+\alpha y=0,\mathrm{d}y+\alpha ^{}y=0$$
on the circle have the same monodromy matrix $`M`$ up to conjugacy, since that is the condition that $`y`$ and $`\widehat{y}`$ can be chosen so that $`g=y(z)\widehat{y}(z)^1`$ is single valued; $`g`$ is then the element of $`LG`$ that maps one system into the other.
Since $`c`$ vanishes on $`L_+𝔤`$, the action of $`LG`$ on $`𝒪`$ is Hamiltonian when restricted to $`L_{r+1}G`$, the subgroup of loops of the form $`1+z^{r+1}h`$, where $`h`$ is holomorphic. We denote by $`_r`$ the Marsden-Weinstein reduction of $`\mu _{r+1}^1(0)`$, where $`\mu _{r+1}`$ is the moment map. That is, $`_r`$ is the quotient of $`𝒜_r𝒪(A)`$ by the action of $`L_{r+1}G`$. It a finite-dimensional complex symplectic manifold, with symplectic form that we shall again denote by $`\sigma `$. Its points are elements of $`𝒜_r𝒪(A)`$, modulo gauge transformations of the form
$$\alpha g^1\mathrm{d}gg^1\alpha g,gL_{r+1}G.$$
A tangent vector $`YT_\alpha _r`$ is of the form
$$\delta \alpha =\mathrm{\Omega }=\mathrm{d}\mathrm{\Omega }+[\alpha ,\mathrm{\Omega }]$$
where $`\mathrm{\Omega }`$ satisfies (5); two such $`\mathrm{\Omega }`$s define the same tangent vector whenever their difference is in $`L_{r+1}𝔤`$.
We can construct from $`\alpha 𝒜_r`$ the following objects (see appendix).
> (i) The singularity data $`(m,t)`$, where $`t`$ is a diagonal polynomial of degree $`r1`$ and $`m`$ is the exponent of formal monodromy.
>
> (ii) The formal series $`g_\mathrm{f}^{\text{}}=g_iz^i`$.
>
> (iii) The matrices $`C_i`$, defined as follows. For each $`\alpha `$, choose a solution $`y`$ to the corresponding system with fixed monodromy matrix $`M`$ and choose $`2r`$ sectors $`𝒮_i`$ at $`z=0`$, as in (19). Then put $`C_i=y_i(z)^1y(z)`$, where $`y_i`$ is the corresponding special solution and we continue $`y`$ in the positive sense around $`z=0`$ into the sector $`𝒮_i`$. If we put
>
> $$y_{2r+1}=y_1\mathrm{e}^{2\pi \mathrm{i}m}𝒮_{2r+1}=𝒮_1,$$
>
> then $`C_{2r+1}=\mathrm{e}^{2\pi \mathrm{i}m}C_1M`$ and we can define the Stokes’ matrices by $`S_i^{\text{}}=C_i^{\text{}}C_{i+1}^1`$ ($`i=1,\mathrm{}2r`$).
These are not quite uniquely determined: we are free to make the replacement
$$g_\mathrm{f}^{\text{}}g_\mathrm{f}^{\text{}}T,C_iT^1C_i,S_iT^1S_iT,$$
(12)
where T is diagonal and independent of $`z`$. We shall express the symplectic form on $`_r`$ in terms of these variables.
Given $`t`$ and the monodromy matrix $`M`$, the Stokes’ matrices and the exponent of formal monodromy satisfy two constraints.
> (C1) Let $`P_i`$ be the matrix of the permutation that puts the real parts of diagonal entries in $`z^rt`$ is increasing order in $`𝒮_i𝒮_{i+1}`$. Then for each $`i`$, $`P_i^1S_iP_i`$ is upper triangular and $`P_iS_{i+1}P_i^1`$ is lower triangular, both with ones on the diagonal. This follows from the fact that
>
> $$\mathrm{exp}(z^rt+m\mathrm{log}z)S_i\mathrm{exp}(z^rtm\mathrm{log}z)1$$
> (13)
> faster than any power of $`z`$ as $`z0`$ in $`𝒮_i𝒮_{i+1}`$.
>
> (C2) The product $`\mathrm{e}^{2\pi \mathrm{i}m}S_1\mathrm{}S_{2r}`$ is conjugate to $`M^1`$.
Denote by $`𝒞_r`$ the set matrices $`S_i\mathrm{SL}(n,)`$, $`m`$ diagonal and trace-free, satisfying these two constraints. Given a point $`𝒞_r`$, we choose $`C_1`$ such that
$$\mathrm{e}^{2\pi \mathrm{i}m}S_1\mathrm{}S_{2r}=C_1^{\text{}}M^1C_1^1$$
and define $`C_2,\mathrm{},C_{2r+1}`$ by $`C_{i+1}=S_i^1C_i`$. Put
$$\omega =\frac{1}{2\pi \mathrm{i}}\underset{1}{\overset{2r}{}}\mathrm{tr}\left(\mathrm{d}C_i^{\text{}}C_i^1\mathrm{d}S_i^{\text{}}S_i^1\right)+\pi \mathrm{i}\mathrm{tr}(\mathrm{d}m\mathrm{d}m)\mathrm{tr}(\mathrm{d}C_1^{\text{}}C_1^1\mathrm{d}m).$$
(14)
It is shown in the appendix that $`\omega `$ is skew-symmetric, and in fact a symplectic form on $`𝒞_r`$.
For each point of $`_r`$, we pick a representative in $`\alpha 𝒜_r`$. We then define 1-forms on $`_r`$ by
$$\mathrm{\Theta }=\mathrm{d}g^{\text{}}g^1+\frac{g\mathrm{d}tg^1}{z^r}\gamma =g_0^1\mathrm{d}g_0,$$
where $`\mathrm{d}`$ is the exterior derivative on $`_r`$ and $`g`$ is the polynomial obtained by truncating the formal power series $`g_\mathrm{f}`$ at some large power of $`r`$. With this notation, the symplectic form on $`_r`$ is given by the following proposition.
###### Proposition 5
The symplectic form on $`_r`$ is
$$\sigma =\frac{1}{2\pi \mathrm{i}}\mathrm{tr}\left(\mathrm{\Theta }\mathrm{\Theta }^{}\right)+\mathrm{tr}\left(\gamma \mathrm{d}m\right)\omega .$$
where $`\mathrm{\Omega }=_z\mathrm{\Omega }\mathrm{d}z+[\alpha ,\mathrm{\Omega }]`$.
The proof is by splitting the integral in the definition of $`\sigma `$ into sections lying in the various sectors, and then shrinking the contour to zero. The details are given in the appendix.
The formula for $`\sigma `$ is independent of the choices made in defining the variables on $`_r`$. In particular, it depends on the first $`r`$ terms in the formal series since the right-hand side of the formula is unchanged when $`\mathrm{\Theta }`$ is replaced by $`\mathrm{d}hh^1+h\mathrm{\Theta }h^1`$, where $`h=1+O(z^{r+1})`$. It is also unchanged when $`g`$, $`C_i`$ are replaced by $`gT`$, $`T^1C_i`$ where $`T`$ is diagonal and indpendent of $`z`$.
We can see the local structure of $`_r`$ from the proposition. The submanifolds on which $`g`$ and $`t`$ are constant (up to the freedom 12) are symplectomorphic to $`𝒞_r`$. While, those on which $`m`$ and $`S_i`$ are constant (up to 12) are symplectomorphic to a fixed manifold $`𝒫_r`$; by mapping $`[\gamma ]\mathrm{\Gamma }_r`$ to $`[\alpha ]`$, where $`\alpha =\gamma ^{}(\theta )`$, and by noting the coincidence of the formulas for the symplectic forms, we identify $`𝒫_r`$ with $`\mathrm{\Gamma }_r`$.
### Local uniqueness of the full twistor space
Let $`𝒫`$ denote the subset of $`𝒪𝒜_r`$ given by the fixing the values of the Stokes’ matrices and exponent of formal monodromy. By using the actions of $`L_{r+1}G`$ on $`\mathrm{\Gamma }`$ and $`𝒫`$ to pick represenetatives in $`[\gamma ]`$ and $`[\alpha ]`$, we can identify $`\mathrm{\Gamma }`$ with an open neighbourhood in $`𝒫`$ so that $`\gamma \mathrm{\Gamma }`$ corresponds to $`\alpha 𝒫`$ such at $`\alpha =\gamma ^{}(\theta )`$. We then deduce the following proposition.
###### Proposition 6
Suppose that $`(𝒵,X)`$ and $`(𝒵^{},X)`$ are full twistor spaces for a generic linear system of ODEs on a Riemann surface $`X`$, with $`G=\mathrm{SL}(n,)`$. Let $`aX`$ be a pole of order $`r+1>1`$. Then there are neighbourhoods $`U,U^{}`$ of $`a`$ in $`𝒵`$ and $`𝒵^{}`$ and a $`𝔤`$-equivariant biholomorphic map $`\rho :UU^{}`$ such that $`\rho (a)=a`$ and $`\rho (XU)=XU^{}`$.
Proof. Choose a coordinate $`z`$ on a small disc in $`X`$ such that the pole is at $`z=0`$, and extend this to a neighbourhood of $`a`$ in $`𝒵`$ so that $`S`$ is given by $`z=0`$. Then we can identify a neighbourhood of $`a`$ with $`S\times D`$, as before. For each $`sS`$, we have a holomorphic map $`\gamma _s:D𝒵`$ and hence an element $`\alpha _s`$ of $`𝒫`$ such that $`\gamma _s^{}=\alpha _s`$. Let $`\gamma _s^{}:D𝒵^{}`$ be the corresponding map into the second twistor space. Then the required biholomorphic map is $`\rho :(s,z)\gamma _s^{}(z)`$.
Proposition 4 above is an immediate corollary, since Proposition 6 implies that the map $`\rho `$ constructed there extends to $`S`$.
### Isomonodromic deformations
We have shown that a generic $`\mathrm{SL}(n,)`$ system of the form (2) on a Riemann surface can be generated from a twistor space $`(𝒵,X)`$, and that, if we require $`𝒵`$ to be full, then it is unique, at least in a neighbourhood of $`X`$. In the case of Fuchsian equations on $`_1`$, Hitchin (1995a) showed that the isomonodromic deformations of the system are given by the deformations of $`X`$ in $`𝒵`$ (every twistor space being full in the Fuchsian case). This is also true more generally, as we shall now see.
By Kodaira’s theorem, the deformations of a compact curve $`X𝒵`$ are determined by the properties of the normal bundle $`N=T𝒵|_X/TX`$. Put
$$d_1=\mathrm{dim}H^1(X,N)\text{and}d_0=\mathrm{dim}H^0(X,N).$$
When $`d_1=0`$, the theorem implies that $`X`$ is one of a complete $`d_0`$-parameter holomorphic family of embedded curves. The tangent space to the parameter space at $`X`$ is naturally identified with $`H^0(X,N)`$ (Kodaira 1962).
For each curve $`X`$ in the family, we have a meromorphic 1-form $`\theta |_X`$ and hence a system of differential equations of the form (2). As we vary the curve along a path $`X_t`$ in the family, $`t[0,1]`$, the tangent at $`t`$ is an element of $`H^0(X_t,N)`$. This we represent by local sections $`Y_i`$ of $`T𝒵|_{X_t}`$, chosen to be tangent to $`S`$ at the poles. Thus the $`Y_i`$s are uniquely determined up to the addition local tangent vector fields to $`X_t`$ that vanish at $`SX_t`$. Put $`\mathrm{\Omega }_i=i_{Y_i}\theta `$, $`\alpha _t=\theta |_{X_t}`$, and identify local neighbourhoods in the $`X_t`$s along $`Y_i`$. Then $`\mathrm{\Omega }_i`$ is meromorphic, with a pole of order $`r`$ at a singularity of rank $`r`$. By (3),
$$_t\alpha _t=_{\alpha _t}\mathrm{\Omega }_i=\mathrm{d}\mathrm{\Omega }_i+[\alpha _t,\mathrm{\Omega }_i].$$
Moreover on the overlap of their domains, $`\mathrm{\Omega }_i\mathrm{\Omega }_j=i_{T_{ij}}^{\text{}}\alpha _t\mathrm{d}z`$ for some tangent vector $`T_{ij}`$ to $`X`$, which must vanish at any poles in the overlap. By using the results in the appendix, we deduce the following proposition.
###### Proposition 7
Let $`G=\mathrm{SL}(n,)`$, let $`(𝒵,X)`$ be a twistor space, and let $`X^{}`$ be a deformation of $`X`$. Then the linear system of ordinary differential equations on $`X^{}`$ is an isomonodromic deformation of the linear system on $`X`$.
#### The minimal twistor space
In the minimal case, we can find the normal bundle of $`X`$ as follows. For $`xX`$, put $`L_\alpha (x)=\alpha (T_xX)𝔤`$ when $`x`$ is not a pole; and
$$L_\alpha (x)=z^{r+1}\alpha (T_xX),$$
when $`x`$ is a pole of rank $`r`$ and $`z=0`$ at $`x`$. Then $`L_AX_0`$ is a holomorphic line bundle. Moreover $`\alpha `$ is a global meromorphic section of $`L_\alpha K`$ with a pole of order $`r+1`$ at a singularity of rank $`r`$ and, by assumption, no zeros. Therefore
$$L_\alpha K=(r_i+1)a_\alpha $$
and so $`\mathrm{deg}L_\alpha =22g(r_i+1)`$, where the sum is over the singularities and $`g`$ is the genus.
Now a point $`xX`$ is the image of $`(x,e)X\times G`$ under the projection along $`F`$. So we have
$$N_x=T_{(x,e)}(X\times G)/(T_xL_\alpha ).$$
Thus there is a short exact sequence
$$0L_A𝔤N0,$$
where $`𝔤`$ is the trivial bundle $`𝔤`$-bundle over $`X`$.
When $`X=_1`$ and $`(r_i+1)4`$, we have $`H^1(𝔤)=0`$, $`H^0(L_A)=0`$, and
$$\mathrm{dim}H^1(L_A)=3+(r_i+1).$$
From the corresponding long exact sequence, therefore, $`H^1(N)=0`$, and
$$0H^0(𝔤)H^0(N)H^1(L_A)0.$$
It follows that
$$d_1=0,d_0=(r_i+1)3+\mathrm{dim}G,$$
We conclude that $`X`$ is one of a $`d_0`$-parameter of embedded copies of $`_1`$. If all the singularities are regular, then $`d_0`$ is the dimension of the space of configurations of poles ($`(r_i+1)3)`$, plus the dimension of $`G`$. The deformations of $`X`$ are parametrized by the positions of the poles, modulo the action of $`\mathrm{SL}(2,)`$ on $`_1`$, together with constant gauge transformations.
Example. The Schlesinger equations Suppose that $`X=_1`$ and that all the singularities are regular. We choose the domains of the $`\mathrm{\Omega }_i`$s so that each pole lies in only one. The $`\mathrm{\Omega }_i`$s are holomorphic and the deformation is determined by the holomorphic tangent vector fields $`T_{ij}`$ on the overlaps of the domains. Since $`H^1(X,TX)=0`$, we have $`T_{ij}=T_iT_j`$, where $`T_i`$ is holomorphic on the domain of $`\mathrm{\Omega }_i`$ (but possibly non-zero at the corresponding pole). So if we put $`Y=Y_i+T_i`$, then $`Y`$ is a global section of $`T𝒵|_X`$.
Let $`z`$ be a global stereographic coordinate on $`X`$, with $`z=\mathrm{}`$ not one of the poles, and use $`Y`$ to transfer $`z`$ along the deformations of $`X`$. Then
$$\alpha =\frac{A_i\mathrm{d}z}{za_i}$$
where the $`a_i`$s are positions of the poles and the coefficients $`A_i`$ are independent of $`z`$ and satisfy
$$\underset{i}{}A_i=0$$
(since there is no pole at infinity). Put $`\mathrm{\Omega }=i_Y\theta `$. Then $`\mathrm{\Omega }`$ is meromorphic with simple poles at the $`a_i`$s and
$$\frac{A}{t}=\frac{\mathrm{\Omega }}{z}[A,\mathrm{\Omega }].$$
In fact, since $`\mathrm{\Omega }_i`$ is holomorphic at $`a_i`$ and $`T_i(a_i)=a_i/t`$, we have
$$\mathrm{\Omega }=\frac{a_i}{t}\frac{A_i}{za_i}+O\left((za_i)^0\right).$$
It follows that
$$\mathrm{\Omega }=\underset{i}{}\frac{a_i}{t}\frac{A_i}{za_i}+k$$
where $`k`$ is a matrix independent of $`z`$. Therefore
$$\frac{A_i}{t}=\underset{ji}{}\frac{a_i}{t}\frac{[A_i,A_j]}{a_ia_j}+[k,A_i],$$
which is a form of the Schlesinger equations. The last term is simply a infinitesimal gauge transformation; the first gives the dependence of the $`A_i`$s on the configuration of the poles.
#### The full case
In a full twistor space, $`d_0`$ is generally larger than in the minimal construction. Here we can find $`N`$ in another way.
We suppose that all the singularities are generic. Then, by (5) and Lemma 1, a local section $`V`$ of $`T𝒵|_X`$ is a map $`\mathrm{\Omega }`$ from a neighbourhood in $`X`$ to $`𝔤`$ such that a singularity at of rank $`r`$,
$$\mathrm{\Omega }=O(z^r),[\alpha ,\mathrm{\Omega }]=O(z^{r1}),$$
where $`z`$ is a local coordinate that vanishes at the singularity. That is, in a local gauge in which the singular part of $`\alpha `$ is diagonal, the diagonal entries in $`\mathrm{\Omega }`$ have poles of at most order $`r`$, and the off-diagonal entries are holomorphic. Thus these algebraic conditions charaterize $`\mathrm{\Omega }`$ as a local section of a holomorphic bundle $`EX`$ with fibre $`𝔤`$ (and therefore rank $`n^21`$) and degree $`(n1)r_i`$. In the full case, therefore, we have that $`T𝒵|X=E`$ can be constructed directly from the positions and ranks of the singularities of the ODE on $`X`$.
Put
$$L=TX(a_i),$$
so that a local holomorphic section of $`L`$ is a tangent vector field that vanishes at the poles. Then we have a short exact sequence
$$0LEN0$$
with the second map given by contraction with $`\alpha `$. Hence there is an exact sequence
$$0H^0(L)H^0(E)H^0(N)H^1(L)H^1(E)H^1(N)0.$$
If the genus of $`X`$ is $`g`$, and if there are $`k`$ singularities in total, then
$$\mathrm{deg}(L)=22gk,\mathrm{dim}H^0(E)\mathrm{dim}H^1(E)=(n^21)(1g)+(n1)r_i,$$
(the latter identity coming from the Riemann-Roch theorem).
A global section of $`E`$ is a meromorphic map $`\mathrm{\Omega }:X𝔤`$ such that at a singularity of rank $`r`$,
$$\mathrm{\Omega }=z^rg_\mathrm{f}^{\text{}}qg_\mathrm{f}^1+O(z^0)$$
where $`z`$ is a local coordinate that vanishes at the singularity and $`q`$ is a diagonal polynomial of degree $`r1`$. When $`X=_1`$, $`\mathrm{\Omega }`$ is determined as a global rational map by the $`q`$s up to the addition of a constant element of $`𝔤`$, and the $`q`$s can be specified independently. In this case, therefore, $`\mathrm{dim}H^0(E)=(n1)r_i+\mathrm{dim}G`$, and $`\mathrm{dim}H^1(E)=0`$. Moreover if $`k4`$, then
$$\mathrm{dim}H^0(L)=0,\mathrm{dim}H^1(L)=k3.$$
It follows that
$$d_0=H^0(N)=(n1)r_i+\mathrm{dim}G+k3,d_1=\mathrm{dim}H^1(N)=0.$$
If either $`n=2`$ ($`G=\mathrm{SL}(2,)`$) or $`r_i=0`$ (all singularities regular), then $`d_0`$ is the same as in the minimal case; in either of these cases, the minimal twistor space is full and, by Proposition 8 below, it gives all possible isomonodromic deformations. In general, however, there are more isomonodromic deformations than are given by the minimal construction: the additional parameters are the coefficients of the diagonal polynomials $`t`$ (of degree $`r1`$) at the irregular singularities ($`r1`$).
When $`X`$ has higher genus, $`\mathrm{dim}H^1(E)=\mathrm{dim}H^0(E^{}K)^{}`$ is generically zero whenever $`(n1)r_i>n^2g+g2`$.
#### Twistor curves
Let $`\mathrm{d}y+\alpha y=0`$ be a generic system on a compact Riemann surface $`X`$ and suppose that $`H^1(E)=0`$. Then we can construct a full twistor space $`(𝒵,X)`$. Since $`H^1(X,N)=0`$, $`X`$ is one of a complete holomorphic family $`𝒦`$ of curves $`X𝒵`$.
###### Proposition 8
Let $`(X_t,\alpha _t)`$, $`t[0,1]`$, be an isomonodromic deformation of $`(X,\alpha )`$. Then for small $`t`$, there is a path $`X_t`$ in $`𝒦`$ such that $`\alpha _t=\theta |_{X_t}`$.
Proof. Let $`y_t`$ be solution to
$$\mathrm{d}y_t+\theta _ty_t=0$$
with constant monodromy, and with constant connection matrices $`C_i`$ to the special solutions at the poles.
Let $`z,z^{}X`$ be nearby points (neither a pole) and let $`g,g^{}G`$ be close to the identity. Then, by integrating the action of $`𝔤`$ on $`𝒵`$, we have two points $`zg`$, $`z^{}g^{}`$ in $`𝒵`$ near $`X`$. These are the same if
$$gg^1=y_0(z)y_0(z^{})^1.$$
(15)
Let $`z_tX_t`$ vary continuously with $`t`$, and suppose that $`z_t`$ is not a pole for any small $`t`$. Put
$$\rho _t(z_t)=z_0y_0(z_0)y_t(z_t)^1𝒵$$
(16)
(the right-hand side is interpreted by regarding $`z_0X`$ as a point of $`X𝒵`$ and by using the local action of $`G`$ on $`𝒵`$). This is independent of the choice of branch of $`y_t`$ and $`y_0`$ (so long as we make the choice of branch continuously) since $`y_0`$ and $`y_t`$ have the same monodromy. Moreover, $`\rho _t(z_t)`$ depends only on $`z_t`$, and not on the path, by (15). So if we exclude a small neighbourhood of each pole in $`X_t`$, then we can embed the complement in $`𝒵`$ by $`z_t\rho _t(z_t)`$. By fixing $`z_0`$ and moving $`z_t`$, we see from (16) that $`\rho _t^{}\theta =\alpha _t`$.
It remains to show that $`\rho _t`$ extends holomorphically to the poles. Consider one of the poles (a point of $`X_t`$, varying continuously with $`t`$). We can choose a coordinate $`z`$ in a neighbourhood $`D`$ of the pole on each $`X_t`$ so that $`D`$ is the unit disc and the pole is at $`z=0`$. Then, for small $`t`$, since $`𝒵`$ is full, there exists a holomorphic map $`\gamma _t:D𝒵`$ such that $`\alpha _t^{}=\gamma _t^{}\theta `$ has the same singularity data at $`z=0`$ as $`\alpha _t`$ has at $`z=0`$. Since $`\alpha _t`$ is an isomonodromic deformation, $`\alpha _t`$ and $`\alpha _t^{}`$ also have the same Stokes’ matrices.
Let $`y_t^{}`$ be a solution to
$$\mathrm{d}y_t^{}+\alpha _t^{}=0$$
with the same monodromy and connection matrices to the special solutions in the sectors at $`z=0`$ as $`y_0`$. Then
$$\gamma _t(z)=zy_0(z)y_t^{}(z)^1.$$
Further $`y_t^{}y_t^1`$ is holomorphic at $`z=0`$. This is because it is single-valued, since $`y_t`$ and $`y_t^{}`$ have the same holonomy, and bounded since in any sector $`𝒮`$ at $`z=0`$
$$y_t^{}y_t^1=y_𝒮^{}y_𝒮^1g_\mathrm{f}^{}g_\mathrm{f}^1$$
where $`y_𝒮^{}`$, $`y_𝒮^{\text{}}`$ are the corresponding special solutions and $`g_\mathrm{f}^{}`$ , $`g_\mathrm{f}^{\text{}}`$ are the formal gauge transformations to diagonal form. So the embedding $`\rho _t`$ extends by mapping $`zDX_t`$ to $`\gamma _t(z)y_t^{}y_t^1`$.
#### Isomonodromic flows for systems on the Riemann sphere
The number of independent isomonodromic deformations (the dimension of $`𝒦`$) of a generic system on $`X=_1`$,
$$\mathrm{dim}H^0(X,N)=(n1)r_i+\mathrm{dim}G+k3.$$
We shall now show that the deformations are given by Hamiltonian flows on symplectic manifolds constructed from the affine orbits in $`L𝔤^{}`$.
In this case, the twistor curves in $`𝒵`$ are copies of $`_1`$, and they can be parametrized by a global stereographic coordinate $`z\{\mathrm{}\}`$. We denote by $`\widehat{𝒦}`$ the space of parametrized curves, which has dimension
$$\mathrm{dim}\widehat{𝒦}=(n1)r_i+\mathrm{dim}G+k.$$
The points of $`\widehat{𝒦}`$ can be labelled by the positions of the poles ($`k`$ parameters), the polynomials $`t`$ at each pole ($`(n1)r_i`$ parameters) and a choice of gauge ($`\mathrm{dim}G`$ parameters).
An element of $`\widehat{𝒦}`$ is a mapping $`\rho :_1𝒵`$ from some fixed copy of the Riemann sphere. It determines a rational $`𝔤`$-valued $`1`$-form
$$\alpha =\rho ^{}\theta =A\mathrm{d}z,$$
where $`A`$ is rational, with poles of order $`r_i+1`$ at $`k`$ points $`a_1,\mathrm{},a_k`$ (we assume that none of the poles is at infinity, so $`A=O(z^2)`$ as $`z\mathrm{}`$). In a neighbourhood of $`a_i`$, we put $`z_i=za_i`$ and assume, without of loss of generality, that no other pole lies in the closure of the disc $`D_i=\{|z_i|<1\}`$. Then $`\alpha _i=\alpha |_{D_i}`$ determines a point of the symplectic manifold $`_{r_i}`$. Thus we have a map
$$\widehat{𝒦}=_{r_1}\times _{r_2}\times \mathrm{}\times _{r_k}.$$
It is not surjective, since a given point of $``$ is not, in general, given by the restrictions of a global 1-form $`\alpha `$. However $`\alpha `$ is uniquely determined by the positions of its poles and by its image in $``$, since the difference between two $`\alpha `$s with the same pole positions, and determining the same point of $``$, is a global holomorphic $`1`$-form, and therefore vanishes.
Given $`[\alpha _i]_{r_i}`$ and the points $`a_i_1`$, we put $`\alpha _i=A_i\mathrm{d}z_i`$ and denote by $`A_i`$ and $`A_{i+}`$ the negative and non-negative degree terms in the Laurent expansion of $`A_i`$ in powers of $`z_i`$ in a neighbourhood of $`a_i`$. Given also a diagonal polynomial $`q_i`$ of degree $`r_i1`$, we put
$$\mathrm{\Omega }_{q_i}=\left(z^{r_i}g_\mathrm{f}^1q_ig_\mathrm{f}^{\text{}}\right)_{}$$
where $`g_\mathrm{f}`$ is the formal gauge transformation to the diagonal form of $`\alpha _i`$ and again the minus subscript denotes the negative terms in the Laurent expansion in powers of $`z_i`$. Then $`A_i`$ and $`\mathrm{\Omega }_{q_i}`$ are global meromorphic functions on $`_1`$ with values in $`𝔤`$; they are holomorphic except at $`a_i`$, where they have poles of order $`r_i+1`$ and $`r_i`$, respectively. Moreover $`A_i`$ and $`\mathrm{\Omega }_i`$ are independent of the choice of representative in $`[\alpha _i]`$.
###### Proposition 9
The isomonodromic deformations of a generic $`\mathrm{SL}(n,)`$ system on $`_1`$ are generated by the Hamiltonians
$$h_v=\underset{j}{}\frac{1}{2\pi \mathrm{i}}\mathrm{tr}(\alpha _jv),h_i=\underset{ji}{}\frac{1}{2\pi \mathrm{i}}\mathrm{tr}(\alpha _jA_i),h_{q_i}=\underset{j}{}\frac{1}{2\pi \mathrm{i}}\mathrm{tr}\left(\alpha _j\mathrm{\Omega }_{q_i}\right),$$
on $``$, where $`v`$ is a constant element of $`𝔤`$ and the integrals are around small circles surrounding the poles.
The $`h_i`$s are time-dependent Hamiltonians, the ‘times’ being the positions $`a_i`$ of the poles.
Proof. First we note that the Hamiltonians $`h_v`$ generate the constant gauge transformations. Consider next the flow generated by $`h_i`$. We shall find the value of the Hamiltonian vector field at a point of $``$ constructed from a global meromorphic 1-form $`\alpha `$. To do this, we must we must evaluate the gradient of $`h_i`$ at such a point. We have
$`\delta h_i`$ $`=`$ $`{\displaystyle \underset{ji}{}}{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle \mathrm{tr}(\delta \alpha _jA_i+\alpha \delta A_i)}`$
$`=`$ $`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle \mathrm{tr}(\alpha _i\delta A_i)}+{\displaystyle \underset{ji}{}}{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle \mathrm{tr}(\delta \alpha _jA_i)}.`$
However
$`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle \mathrm{tr}(\alpha _i\delta A_i)}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle \mathrm{tr}(\alpha _{i+}\delta A_i)}`$
$`=`$ $`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle \mathrm{tr}(A_{i+}\delta \alpha _i)}.`$
We conclude that the value of the Hamiltonian vector field at such a point is
$$\delta \alpha _i=A_{i+},\delta \alpha _j=A_iji.$$
The claim is that this is tangent to an isomonodromic deformation. To see this, let $`y`$ be a solution to $`\mathrm{d}y+\alpha y=0`$, let $`D`$ be a disc containing $`a_i`$, but no other pole, and let $`D^{}`$ be a second disc not containing $`a_i`$ such that $`D,D^{}`$ is an open cover of $`_1`$. For small $`t`$, put $`F_t(z)=y(zt)y(z)^1`$. Then $`F:DD^{}G`$ is single-valued, holomorphic, and equal to the identity when $`t=0`$. By Birkhoff’s theorem, $`F_t=h_t^1h_t^{}`$ for some holomorphic maps $`h_t:DG`$, $`h_t^{}:D^{}G`$, with $`h_t^{}(\mathrm{})=1`$.
Put
$$y_t(z)=\{\begin{array}{cc}h_t(z)y(zt)\hfill & zD\hfill \\ h_t^{}(z)y(t)\hfill & zD^{}\hfill \end{array}$$
and $`\alpha _t=\mathrm{d}y_ty_t^1`$. Then the definitions agree on $`DD^{}`$ and $`\alpha _t`$ is a global meromorphic 1-form with poles at $`z=a_j`$ ($`ji`$) and $`z=a_i+t`$. Moreover $`_ty_t=\mathrm{\Omega }_ty_t`$, where
$$\mathrm{\Omega }_t=\{\begin{array}{cc}_th_t^{\text{}}h_t^1h_tA_i(zt)h_t^1\hfill & \text{in }D\hfill \\ _th_t^{}h_t^1\hfill & \text{in }D^{}\text{.}\hfill \end{array}$$
Since $`h_t`$ and $`h_t^{}`$ are holomorphic in $`D`$ and $`D^{}`$, it follows that the deformation is isomonodromic (see Proposition 11).
At $`t=0`$, we have $`h_t=h_t^{}=1`$ and
$$_th_t^{}=A_i_th_t=A_{i+};$$
we also have at $`t=0`$ that $`\delta \alpha _i=(_th_t)`$, $`\delta \alpha _j=(_th_t^{})`$ for $`ji`$. So the tangent to the deformation is the Hamiltonian vector field constructed above: these deformations move the poles, but leave the singularity data unchanged.
Now consider the flow generated by $`h_{q_i}`$. Proceeding as before to calculate the value of the Hamiltonian vector field at a point given by a global 1-form $`\alpha `$, we have
$`\delta h_{q_i}`$ $`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle \mathrm{tr}\left(\delta \alpha _j\mathrm{\Omega }_{q_i}+\alpha _j\delta \mathrm{\Omega }_{q_i}\right)}`$
$`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle \mathrm{tr}\left(\alpha _j\delta \mathrm{\Omega }_{q_i}\right)}.`$
So in this case, the value of the Hamiltonian vector field is
$$\delta \alpha _i=_\alpha \mathrm{\Omega }_{q_i},$$
which is clearly isomonodromic. These deformations change the singularity data at $`a_i`$, leaving the position of the poles unchanged.
Remark. The Hamiltonians $`H_{q_i}`$ vanish on the orbit of a global 1-form $`\alpha `$, while
$$h_i=\frac{1}{\tau }\frac{\tau }{a_i},$$
where $`\tau `$ is as in Jimbo et al (1981a).
Example. In the case of a generic system with only regular singularities, $`=_1\times \mathrm{}\times _1`$. The tangent space is spanned $`_1`$ is spanned by the generators of constant gauge transformations. If we write
$$\alpha =A\mathrm{d}z=\frac{A_i\mathrm{d}z}{za_i},\alpha _i=\frac{A_i\mathrm{d}z_i}{z_i},$$
then we can identify $``$ (as a symplectic manifold) with the product of the coadjoint orbits of the $`A_i`$ in $`𝔤^{}`$. In this case,
$$A_i=\frac{A_i}{za_i};$$
there are no $`h_{q_i}`$s, while
$$h_i=\underset{ji}{}\frac{\mathrm{tr}(A_iA_j)}{a_ia_j}.$$
## Appendix
### Singularities of systems of ODEs
Let $`A`$ be a meromorphic map from some neighbourhood $`D`$ of the origin in $``$ into $`𝔤`$ (the Lie algebra $`\mathrm{sl}(n,)`$), with a pole of order $`r+1`$ at $`z=0`$. Then the system
$$\frac{\mathrm{d}y}{\mathrm{d}z}=Ay,$$
(17)
has a singularity of Poincaré rank $`r`$ at the origin. It is regular or Fuchsian if $`r=0`$, and irregular if $`r>0`$. In this paper, $`y`$ will always be a matrix fundamental solution—that is, it will take values in $`G`$.
#### Gauge and point transformations
When we regard the system as a connection on a vector bundle, we must allow gauge transformations (changes of trivialization) of the form
$$Ag^1Agg^1\frac{\mathrm{d}g}{\mathrm{d}z},\text{ }g^1g,ygy,$$
where $`g:XG`$ is holomorphic. When $`g`$ is constant, $`A`$ transforms by conjugation. We also admit transformations of the coordinate $`z\widehat{z}`$, under which
$$A\widehat{A}=A\frac{\mathrm{d}z}{\mathrm{d}\widehat{z}}.$$
If this is to fix the singularity at the origin, then we must take $`\widehat{z}(0)=0`$.
#### Generic irregular singularities
An irregular singularity of rank $`r`$ is generic if the eigenvalues of
$$A_{r1}=z^{r+1}A|_{z=0}$$
are distinct. In this case, we can assume, without loss of generality, that the eigenvalues of $`z^{r+1}A`$ are distinct throughout the neighbourhood. If we choose an ordering for the eigenvalues, then we can find a holomorphic map $`g:DG`$ such that $`z^{r+1}g(z)A(z)g^1(z)`$ is holomorphic and diagonal, with the eigenvalues as diagonal entries.
It follows that we can find a holomorphic map $`DG`$ and a diagonal polynomial $`p(z)`$ of degree $`r`$ such that
$$Ag(z)^1\frac{p(z)}{z^{r+1}}g(z)$$
is holomorphic at the origin. Therefore we have the normal form:
$$A\mathrm{d}\left(\frac{t}{z^r}\right)+\frac{m\mathrm{d}z}{z}+R(z)$$
(18)
where ‘$``$’ denotes gauge equivalence, $`t`$ is a diagonal polynomial of degree $`r1`$, $`m`$ is a constant diagonal matrix called the exponent of formal monodromy, and the remainder $`R`$ is holomorphic at $`z=0`$. Given the local coordinate $`z`$ and the ordering of the eigenvalues, $`t`$ and $`m`$ are uniquely determined by $`A`$, independently of the choice of gauge. We call them the singularity data at $`0`$. If the ordering is changed, then the diagonal entries are permuted; if the coordinate is changed, then $`m`$ is unchanged, while $`tt^{}`$, where $`t^{}`$ is obtained from $`t`$ by making the coordinate transformation and truncating the Taylor series in $`z`$.
If one looks for a gauge transformation $`g(z)=g_0+g_1z+\mathrm{}`$ such that (18) holds with $`R=0`$, then the coefficients $`g_i`$ can be determined uniquely once a choice has been made for $`g_0`$ to diagonalize $`A_{r1}`$. For each such choice, one can therefore find a unique formal solution
$$y_\mathrm{f}^{\text{}}=g_\mathrm{f}^{\text{}}(z)\mathrm{exp}\left(Tz^r+m\mathrm{log}z\right).$$
In general, the formal series does not converge. However, by truncating, one can make $`R`$ vanish to arbitrarily high order at $`z=0`$.
#### Sectors and special solutions
The eigenvalues $`\lambda _i`$ of $`A_{r1}`$ determine a sequence of Stokes’ rays through the origin, on which $`\mathrm{Re}\left(z^r(\lambda _i\lambda _j)\right)`$ changes sign for some pair of eigenvalues. Given a pair of consecutive rays (in order around the unit circle), with arguments $`\theta _1,\theta _2`$, we define a sector $`𝒮`$ by
$$𝒮=\{z|\theta _1\pi /2r<\mathrm{arg}(z)<\theta _2+\pi /2r\}.$$
(19)
For each such sector $`𝒮`$, there is a unique special solution
$$y_𝒮^{\text{}}=g_𝒮^{\text{}}(z)\mathrm{exp}\left(tz^r+m\mathrm{log}z\right),$$
such that $`g_𝒮^{\text{}}g_\mathrm{f}`$ as $`z0`$ in $`𝒮`$. (See, for example, Boalch 2000 for a careful account of this proposition).
The solutions $`y_𝒮^{\text{}}`$ are independent of the choice of coordinate (as maps $`DG`$). They are uniquely determined in each sector by the choice of $`g_0`$. So, as sections of the principal bundle over each sector, they are determined uniquely by the choice of the frame at the origin in which $`A_{r1}`$ is diagonal.
#### Regular singularities
In the regular case ($`r=0`$), (18) still holds, with $`t=0`$, provided that the eigenvalues of $`A_1`$ are distinct. The formal series for $`g`$ can be found provided that, in addition, no two eigenvalues differ by an integer. It then necessarily converges, so $`y_\mathrm{f}`$ is a solution.
#### Global systems
Let $`X`$ be a compact Riemann surface. For each $`k`$-tuple $`𝐫=(r_1,\mathrm{},r_k)`$ of nonnegative integers, we denote by $`𝒟_𝐫(X,E)`$, or simply by $`𝒟_𝐫`$, the space of meromorphic $`\mathrm{sl}(n,)`$-connections $`=\mathrm{d}+\alpha `$ on the trivial vector bundle $`E=X\times ^n`$ with $`k`$ poles of order (at most) $`r_1+1,\mathrm{},r_k+1`$.<sup>8</sup><sup>8</sup>8Note that $`\alpha `$ determines a holomorphic map $`X𝔤`$
If we choose a local trivialization and a coordinate $`z`$ that vanishes at one of the poles, a connection $`𝒟_𝐫`$ is given in a neighbourhood of the pole by (17).<sup>9</sup><sup>9</sup>9 When $`X=_1`$, we have $`\alpha =A\mathrm{d}z`$, where $`z`$ is a stereographic coordinate and $`A`$ is defined globally as a rational section of $`𝔤𝒪(2)`$. If all its singularities are at finite values of $`z`$, then $`A`$ has $`N`$ poles and $`A=O(z^2)`$ as $`z\mathrm{}`$; but if one of the singularities is at infinity, then $`A=O(z^{r1})`$ as $`z\mathrm{}`$, where $`r`$ is the rank.
#### The monodromy representation
A local $`G`$-valued solution $`y`$ to the equation $`\mathrm{d}+\alpha y=0`$ can be continued analytically: it is singular at the poles, and multi-valued (single-valued on the covering space of the complement of the poles). If
$$\gamma :[0,1]_1\{a_1,\mathrm{},a_k\},\gamma (0)=\gamma (1)=z_0,$$
is a closed loop and $`z_0`$ is some fixed base point, then we have $`y(\gamma (1))=M_\gamma G`$ for some monodromy matrix $`M_\gamma `$ which depends only on the homotopy class of $`\gamma `$.
###### Definition 3
The monodromy representation of (2) is the homomorphism
$$\pi _1(_1\{a_1,\mathrm{},a_k\})G:[\gamma ]M_\gamma .$$
The monodromy representation is independent of $`z_0`$ and the choice of $`y`$ up to conjugation by a fixed element of $`G`$.
#### Deformations
###### Definition 4
Let $`_0,_1𝒟_𝐫(X)`$. A deformation of $`_0`$ into $`_1`$ is a smooth path $`_t𝒟_𝐫(X)`$, $`t[0,1]`$, from $`_0`$ to $`_1`$.
We are interested in the deformations of a connection $`_0`$ into a second $`_1`$ (or equivalently of $`\alpha _0`$ into $`\alpha _1`$) while preserving certain properties of the corresponding linear system.
###### Proposition 10
A deformation $`_t`$ has constant monodromy representation (up to conjugation) if and only if
$$\frac{\mathrm{d}_t}{\mathrm{d}t}=_t\mathrm{\Omega }_t$$
for some family of holomorphic maps $`\mathrm{\Omega }_t`$ (depending smoothly on $`t`$) from the complement of the poles of $`_t`$ in $`X`$ into $`𝔤`$.
There is an awkwardness in the terminology here: it is important to keep in mind that ‘having the same monodromy representation’ is not the same as ‘isomonodromic’ when irregular singularities are present.
Proof. By fixing a base point (disjoint from the poles) and a frame at the base point, we can find a solution $`y_t`$ for each $`t`$ which depends smoothly on $`t`$. If the monodromy representation is constant, then we can find a matrix $`K_tG`$ for each $`t`$ such that the monodromy matrices of $`y_tK_t`$ are constant. If we take $`t`$ close to $`t^{}`$ and exclude small discs around the poles of $`_t`$, then $`g_{tt^{}}=y_tK_tK_t^{}^1y_t^{}^1`$ is single-valued, and we can construct $`\mathrm{\Omega }_t`$ in (C1) by putting
$$\mathrm{\Omega }_t=\frac{g_{tt^{}}}{t^{}}|_{t^{}=t}.$$
(This is holomorphic except at the poles of $`_t`$).
Conversely, if we are given $`\mathrm{\Omega }_t`$, then $`\mathrm{d}A\mathrm{d}z\mathrm{\Omega }\mathrm{d}t`$ is a flat connection on the trivial bundle over
$$(X\text{poles})\times [0,1].$$
Its holonomy coincides with the monodromy $`A_t`$ for each $`t`$, and so the monodromy representation must be constant up conjugacy.
#### Isomonodromic deformations
We now consider deformations $`(X_t,_t)`$ in which we change both $`X`$ and the connection $``$.
Let $`a`$ be a pole of $``$ and $`𝒮`$ a sector at $`a`$. Then we have a solution $`y_{𝒮,a}^{\text{}}`$ to the system $`y=0`$, uniquely determined up to the choice of the frame at $`a`$ in which the leading coefficient of $`A`$ is diagonal. If $`a^{}`$ and $`𝒮^{}`$ are another pole and sector at $`a^{}`$, and if $`\gamma `$ is a path with initial point near $`a`$ in $`𝒮`$ and endpoint near $`a^{}`$ in $`𝒮^{}`$, then we can continue $`y_{𝒮,a}^{\text{}}`$ along $`\gamma `$. We shall have
$$y_{𝒮,a}^{\text{}}=y_{𝒮^{},a^{}}^{\text{}}C_{𝒮,𝒮^{},a,a^{},\gamma }^{\text{}}$$
where $`C_{𝒮,𝒮^{},a,a^{},\gamma }^{\text{}}`$ is a constant matrix. These matrices are uniquely determined by $`A`$ up to
$$C_{𝒮,𝒮^{},a,a^{},\gamma }^{\text{}}D_a^{\text{}}C_{𝒮,𝒮^{},a,a^{},\gamma }^{\text{}}D_a^{}^1,$$
where, for each pole $`a`$, $`D_a`$ is the product of a diagonal matrix and a permutation matrix. The matrices connecting the special solutions in adjacent sectors at the same pole are called Stokes’ matrices.
As we deform $``$ and $`X`$, we can vary $`𝒮,𝒮^{},\gamma `$, and the special solutions continuously.
###### Definition 5
A deformation is isomonodromic if the exponents of formal monodromy and the matrices $`C_{𝒮,𝒮^{},a,a^{},\gamma }^{\text{}}`$ are constant, for an appropriate choice of special solutions.
An isomonodromic deformation has constant monodromy representation, but the converse is not true except in the Fuchsian case (all singularities regular). The following characterization of the isomonodromy property is implicit in Jimbo, Miwa, and Ueno (1981a).
We can cover $`X_t`$ by discs $`D_i`$ varying continuously with $`t`$ such that each pole lies in just one disc. On each disc $`D_i`$, we can choose a coordinate $`z_i`$ such that, if $`a_iD`$ is a pole, then $`z_i=0`$ at $`a_i`$, independently of $`t`$. We shall use these coordinates to identify the discs as $`t`$ varies.
###### Proposition 11
A deformation $`(X_t,_t)`$ with constant monodromy representation is isomonodromic if and only if
> (i) in $`D_i`$, $`\mathrm{d}_t/\mathrm{d}t=_t\mathrm{\Omega }_{it}`$ for some meromorphic $`\mathrm{\Omega }_{it}:D_i𝔤`$ (depending smoothly on $`t`$); $`\mathrm{\Omega }_{it}`$ is holomorphic except possibly at a singularity of $`\alpha _t`$, where it has a pole of order at most $`r_i`$ if $`\alpha _t`$ has a singularity of rank $`r_i`$ in $`D_i`$;
>
> (ii) on the intersection $`D_iD_j`$ of two discs $`\mathrm{\Omega }_{it}\mathrm{\Omega }_{jt}=i_{T_{ij}}\alpha _t`$ for some holomorphic vector field $`T_{ij}`$.
Proof. We shall look at the proof in outline. Suppose that the deformation is isomonodromic. Let $`y_t`$ be a solution to $`_ty_t=0`$, depending continuously on $`t`$ and with constant monodromy (we have to keep in mind that $`y_t`$ is multi-valued and singular at the poles).
Let $`D_i`$ be a disc containing a pole (at $`z_i=0`$) and put
$$g_{itt^{}}(z_i)=y_𝒮^{\text{}}(z_i)y_𝒮^{}^1(z_i)$$
where $`y_𝒮^{\text{}}`$ and $`y_𝒮^{}^{}`$ are the special solutions at $`t`$ and $`t^{}`$ in the corresponding sectors at one of the poles. Then, for $`t^{}`$ close to $`t`$, $`g_{tt^{}}`$ is a single-valued holomorphic map $`D_i\{0\}𝔤`$; it is independent of sector, because the Stokes’ matrices are the same at $`t`$ and $`t^{}`$. Once it is established that it is possible to differentiate the asymptotic expansions term-by-term, it is immediate that $`\mathrm{\Omega }_{it}=_t^{}g_{tt^{}}|_{t^{}=y}`$ is meromorphic, and of the required form.
On a disc $`D_i`$ that does not contain a singularity, we put $`g_{itt^{}}=y_t(z_i)y_t^{}(z_i)^1`$, and define $`\mathrm{\Omega }_{it}`$ in the same way. If we choose the branch of $`y_t`$ to vary continuously with $`t`$, $`g_{itt^{}}`$ is independent of the choice of branch because the monodromy of $`y_t`$ is independent of $`t`$.
Given $`y_t`$, the only freedom in the construction of $`\mathrm{\Omega }_{it}`$ is in the choice of the coordinate $`z_i`$, and hence in the local identification of the discs on the different Riemann surfaces. A different choice for each $`t`$ will add $`i_T\alpha `$ to $`\mathrm{\Omega }_{it}`$ for some local holomorphic vector field $`T`$. Thus (iii) holds on the overlap of two discs.
To prove the converse, suppose that $`\mathrm{\Omega }_{it}`$ is meromorphic, as stated. Choose a continuously varying sector $`𝒮`$ at the pole, and let
$$y_𝒮^{\text{}}(z_i,t)=g_{𝒮,a}^{\text{}}(z_i)\mathrm{exp}\left(tz_i^r+m\mathrm{log}z_i\right),$$
be the corresponding special solution. Then, by writing $`_t=\mathrm{d}_zA_t\mathrm{d}z_i`$ and dropping the subscripts, we have
$$\frac{}{z}\left(\frac{y}{t}\mathrm{\Omega }y\right)=\frac{A}{t}y+A\frac{y}{t}\frac{\mathrm{\Omega }}{t}y\mathrm{\Omega }Ay=A\left(\frac{y}{t}\mathrm{\Omega }y\right).$$
It follows that
$$\frac{y}{t}\mathrm{\Omega }y=yK$$
for some matrix $`K`$, which can depend of $`t`$ but not $`z`$. Therefore
$$g_𝒮^1\left[\frac{g_𝒮^{\text{}}}{t}+g_𝒮^{\text{}}\frac{}{t}\left(\frac{t}{z^r}\right)\mathrm{\Omega }g_𝒮^{\text{}}\right]=\mathrm{e}^{tz^r+m\mathrm{log}z}K\mathrm{e}^{tz^rm\mathrm{log}z}.$$
The left-hand side is asymptotic to a power series, divided by $`z^{r+1}`$, as $`z0`$ in $`𝒮`$ (the same series for each sector at the pole). In the case $`r>0`$, each off-diagonal entry on the right-hand side has an exponential factor which must blow up as $`z0`$ along some directions in $`𝒮`$ since the angle of the sector $`𝒮`$ is more than $`\pi /r`$. This is a contradiction unless the off-diagonal entries in $`K`$ all vanish. Thus $`K`$ is a $`z`$-independent diagonal matrix. It can be absorbed into the special solutions to give that
$$\frac{y_𝒮^{\text{}}}{t}=\mathrm{\Omega }_ty_𝒮^{\text{}}$$
and hence that the $`C`$ matrices are constant. This is also true, more simply, in the regular case since the formal solutions then converge.
### Symplectic form on $`𝒞_r`$
We prove here that the tensor in (14) is a symplectic form on $`𝒞_r`$.
###### Proposition 12
$`\omega `$ is a symplectic form on $`𝒞_r`$, independent of the choice of $`C_1`$.
Proof. From the definitions, $`C_{2r+1}=\mathrm{e}^{2\pi \mathrm{i}m}C_1M`$ and, in variational notation,
$$\delta S_iS_i^1=C_i^{\text{}}(C_i^1\delta C_i^{\text{}}C_{i+1}^1\delta C_{i+1}^{\text{}})C_i^1.$$
(20)
We must show that $`\omega `$ is skew-symmetric, closed, and non-degenerate. From the first constraint, we have $`\mathrm{tr}(\delta S_i^{\text{}}S_i^1\delta ^{}S_i^{\text{}}S_i^1)=0`$. It follows that
$`0`$ $`=`$ $`{\displaystyle \underset{1}{\overset{2r}{}}}\mathrm{tr}\left(\delta S_i^{\text{}}S_i^1\delta ^{}S_i^{\text{}}S_i^1\right)`$
$`=`$ $`{\displaystyle \underset{1}{\overset{2r}{}}}\mathrm{tr}((C_i^1\delta C_i^{\text{}}C_{i+1}^1\delta C_{i+1}^{\text{}})(C_i^1\delta ^{}C_i^{\text{}}C_{i+1}^1\delta ^{}C_{i+1}^{\text{}}).`$
However
$`{\displaystyle \underset{1}{\overset{2r}{}}}\mathrm{tr}\left(C_{i+1}^1\delta C_{i+1}^{\text{}}C_{i+1}^1\delta ^{}C_{i+1}^{\text{}}\right)`$ $`=`$ $`{\displaystyle \underset{1}{\overset{2r}{}}}\mathrm{tr}\left(C_i^1\delta C_i^{\text{}}C_i^1\delta ^{}C_i^{\text{}}\right)4\pi ^2\mathrm{tr}(\delta m\delta ^{}m)`$
$`2\pi \mathrm{i}\mathrm{tr}\left(\delta m\delta ^{}C_1^{\text{}}C_1^1+\delta ^{}m\delta C_1^{\text{}}C_1^1\right).`$
The skew-symmetry follows. A similar calculation, starting from
$$\mathrm{tr}(\delta S_i^{\text{}}S_i^1\delta ^{}S_i^{\text{}}S_i^1\delta ^{\prime \prime }S_i^{\text{}}S_i^1)=0,$$
shows that $`\omega `$ is closed. To show that $`\omega `$ is nondegenerate, we note that if $`\omega (Y,)=`$, then $`\delta C_1^{\text{}}C_1^1\pi \mathrm{i}\delta m`$ is anti-diagonal, $`P_i^{\text{}}\delta C_i^{\text{}}C_iP_i^1`$ is lower triangular for each $`i`$, and $`P_i^{\text{}}\delta C_i^{\text{}}C_i^1P_i^1`$ is upper triangular. However, from (20),
$$\delta C_{i+1}^{\text{}}C_{i+1}^1=S_i^1\delta C_i^{\text{}}C_i^1S_iS_i^1\delta C_i^{\text{}}C_i^1S_i.$$
Therefore $`\delta C_i^{\text{}}C_i^1`$ is diagonal and so $`\delta C_i^{\text{}}C_i^1=\pi \mathrm{i}\delta m`$ for each $`i`$. It then follows from the second constraint (C2) in the defintion of $`𝒞_r`$ that $`\delta m=0`$.
If we make a different choice for $`C_1`$ at each point, then the effect is to replace $`C_i`$ by $`C_iK`$, where $`K`$ is independent of $`i`$. This adds
$$\frac{1}{2\pi \mathrm{i}}\underset{1}{\overset{2r}{}}\mathrm{tr}\left(C_i^{\text{}}\delta KC_i^1\delta ^{}S_i^{\text{}}S_i^1\right)\mathrm{tr}(C_1^{\text{}}\delta KC_1^1\delta ^{}m)$$
to $`\omega (Y,Y^{})`$, which vanishes by (20).
### Proof of Proposition 5
The manipulations are slightly more transparent in the classical variational notation, although it is straightforward to translate this into the language of differential forms.
We shall evaluate $`\sigma (Y,Y^{})`$ in (11) by putting $`\mathrm{\Omega }=\delta yy^1`$, $`\mathrm{\Omega }^{}=\delta ^{}yy^1`$. We shall then shrink the contour to the origin and use the asymptotic behaviour of the $`y_i`$s.
For each $`i`$, choose $`z_i𝒮_i𝒮_{i+1}`$ on the contour with $`z_{2r+1}=z_1`$, and define $`\mathrm{log}z`$ by making a cut along the ray through the origin and $`z_1`$. On each sector,
$$\mathrm{\Omega }=\mathrm{\Omega }_i+y_i\delta C_i^{\text{}}C_i^1y_i^1,$$
where $`\mathrm{\Omega }_i=\delta y_iy_i^1`$. Moreover,
$$\left(y_i^{\text{}}\delta C_i^{\text{}}C_i^1y_i^1\right)=0,$$
since $`C_i`$ is independent of $`z`$. Therefore, in the notation of (11),
$$\frac{1}{2\pi \mathrm{i}}\mathrm{tr}(\mathrm{\Omega }\mathrm{\Omega }^{})=\frac{1}{2\pi \mathrm{i}}\underset{1}{\overset{2r}{}}_{z_i}^{z_{i+1}}\mathrm{tr}\left(\mathrm{\Omega }_i^{\text{}}\mathrm{\Omega }_i^{}\right)+\frac{1}{2\pi \mathrm{i}}\underset{1}{\overset{2r}{}}\left(x_i(z_{i+1})x_i(z_i)\right)$$
(21)
where $`z_i`$ is some point on the contour in $`𝒮_i𝒮_{i+1}`$, $`x_i=\mathrm{tr}(\delta C_i^{\text{}}C_i^1y_i^1\delta ^{}y_i^{\text{}})`$, and the integrals are along segments of the contour. However,
$$x_{i+1}x_i=\mathrm{tr}\left(\delta C_i^{\text{}}\delta C_i^1\delta ^{}S_i^{\text{}}S_i^1S_i^1\delta S_iy_{i+1}^1\delta ^{}y_{i+1}\right).$$
(22)
This follows from the two relations $`y_i=y_{i+1}S_i^1`$ and $`C_i=S_iC_{i+1}`$, which imply that
$$\mathrm{tr}(\delta C_iC_i^1y_i^1\delta ^{}y_i)=\mathrm{tr}\left(\delta C_iC_i^1S_iy_{i+1}^1(\delta ^{}y_{i+1})S_i^1\delta C_iC_i^1\delta ^{}S_iS_i^1\right)$$
and
$$S_i^1\delta C_i^{\text{}}C_i^1S_i=\delta C_{i+1}^{\text{}}C_{i+1}^1S_i^1\delta S_i^{\text{}}.$$
As $`z0`$ in $`𝒮_i𝒮_{i+1}`$, the second term on the right-hand side of (22) goes to zero by (13). Moreover,
$$x_{2r+1}x_1=4\pi ^2\mathrm{tr}(\delta m\delta ^{}m)+2\pi \mathrm{itr}(\delta C_1C_1^1\delta ^{}m)2\pi \mathrm{i}(\delta my_1^1\delta y_1).$$
To deal with the first term in (21), we note that
$$\mathrm{\Omega }_i=\mathrm{\Theta }+g\delta mg^1\mathrm{log}z+O(z^N)$$
as $`z0`$ in $`𝒮_i`$ for some large $`N`$ (depending on the truncation of the formal power series). Therefore in $`𝒮_i`$
$`\mathrm{tr}(\mathrm{\Omega }_i\mathrm{\Omega }_i^{})=`$
$`\mathrm{tr}(\mathrm{\Theta }\mathrm{\Theta }^{})+\mathrm{tr}\left[\left(g^1\delta g+{\displaystyle \frac{\delta t}{z^r}}\right)\delta ^{}m\left(g^1\delta ^{}g+{\displaystyle \frac{\delta ^{}t}{z^r}}\right)\delta m\right]{\displaystyle \frac{\mathrm{d}z}{z}}`$
$`+_z\left(\mathrm{tr}\left(\delta ^{}m(g^1\delta g+z^r\delta t)\right)\mathrm{log}z\right)\mathrm{d}z+\mathrm{tr}(\delta m\delta ^{}m)\mathrm{log}z{\displaystyle \frac{\mathrm{d}z}{z}}+O(N^{}),`$
for some large $`N^{}`$. We therefore have
$`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle \underset{1}{\overset{2r}{}}}{\displaystyle _{z_i}^{z_{i+1}}}\mathrm{tr}(\mathrm{\Omega }_i\mathrm{\Omega }_i^{})`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle \mathrm{tr}(\mathrm{\Theta }\mathrm{\Theta }^{})}+\mathrm{tr}\left(\gamma \delta ^{}m\gamma ^{}\delta m\right)`$
$`+\mathrm{tr}(\delta ^{}my_1^1\delta y_1)_{z_1}+\pi \mathrm{itr}(\delta m\delta ^{}m)+ϵ`$
where $`ϵ0`$ as the contour is shrunk. The proposition follows by putting the two terms together, by shrinking the contour towards $`z=0`$, and by using the definition (14).
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# Abstract.
#### Abstract.
The Ginzburg–Landau (GL) equations of superconductivity provide a computational model for the study of magnetic flux vortices in type-II superconductors. In this article we show through numerical examples and rigorous mathematical analysis that the GL model reduces to the frozen-field model when the charge of the Cooper pairs (the superconducting charge carriers) goes to zero while the applied field stays near the upper critical field.
#### Key words:
Ginzburg–Landau equations, superconductivity, frozen-field approximation, asymptotic analysis.
## 1 Introduction
Superconducting materials hold great promise for technological applications. Especially since the discovery of the so-called high-temperature superconductors in the 1980s, much research has been devoted to understanding the behavior of these new materials. While conventional superconductors require liquid helium (3–4 degrees Kelvin) to remain in the superconducting state, high-temperature superconductors can be cooled with liquid nitrogen (76 degrees Kelvin)—a clear economic advantage. Unfortunately, high-temperature superconductors are ceramic materials, which are difficult to manufacture into films and wires, but progress is being made all the time.
High-temperature superconductors belong to the class of type-II superconductors. Unlike type-I superconductors, type-II superconductors can sustain a magnetic flux in their interior, but this flux is restricted to quantized amounts—filaments that are encircled by a current. The current shields the magnetic flux from the bulk, which is perfectly superconducting. The configuration resembles that of a vortex in a fluid, and the superconductor is said to be in the vortex state.
Figure 1 gives a sketch of the phase diagram of a type-II superconductor in the neighborhood of $`T_c`$, the critical temperature. The two-dimensional phase space is spanned by the temperature $`T`$ and the (magnitude of the) magnetic field $`H`$ and is roughly divided into three subregions. Each subregion corresponds to a particular state: the perfectly superconducting (Meissner) state below the lower critical field $`H_{c1}`$, where no magnetic field can penetrate the medium; the normal state above the upper critical field $`H_{c2}`$, where the superconductor behaves like a normal metal; and the intermediate vortex state. Above the critical temperature $`T_c`$ all superconducting properties are lost.
The vortices, and especially their dynamics, determine the current-carrying capabilities of a superconductor. Much effort, both experimental and theoretical, is therefore being spent on the study of vortex dynamics and, especially, mechanisms to inhibit vortex motion when the superconductor is subject to currents and fields. By “pinning” the vortices, one prevents energy dissipation and, hence, loss of superconductivity.
Vortices can be studied computationally at various levels of detail using different models. The Ginzburg–Landau (GL) model gives a field (continuum) description that, although phenomenological and not based on any microscopic quantum-mechanical theory, has been used successfully to study both the dynamics and the structure of vortex systems in realistic superconductor configurations . Figures 2 and 3 give two examples of computational results obtained with the GL equations. They illustrate both the effectiveness and the difficulties of such calculations.
Figure 2 shows a vortex configuration in a two-dimensional cross section of a twinned superconducting crystal, which was computed from a steady-state solution of the GL equations. The twin boundary (an irregularity in the structure of the crystal) is visible in the horizontal line through the center; it acts as a pinning site for the vortices. The field is perpendicular to the plane of the cross section, which measures $`128\times 192`$ coherence lengths (a characteristic length of the order of microns). Each dot corresponds to a vortex intersecting the plane of the cross section; the entire configuration has approximately 2,700 vortices. The figure shows the level of detail one can achieve with the GL model, given sufficient computing power. At the same time, it illustrates the level of computational complexity one faces if one uses the GL model.
Figure 3 shows a series of snapshots of a vortex configuration in three dimensions, also computed with the GL model. The objective of this computation was to simulate vortex motion through columnar defects and study the potential of the latter as pinning sites. The defects are visible as twisted straight lines. The vortices are the flexible tube-like structures; they move from one defect to another under the influence of external forces. The figure shows the motion of a vortex that is originally pinned on a defect. The vortex develops a loop, the loop peels off, the loop expands in both directions in a traveling-wave-like scenario, and gradually the entire vortex transfers to the next available defect.
Numerical simulations provide the only way to study vortex dynamics at this level of detail. They are an invaluable tool for fundamental research, complementing experiment and theory. Numerical simulations of realistic superconductors based on the GL model, like the ones illustrated in Figs. 2 and 3 are, however, extremely time consuming, and it is desirable to use simpler models whenever possible. Here, we focus on the “frozen-field model,” which is still a continuum model and the closest approximation to the full GL model. In the frozen-field model, the superconducting phenomena are decoupled from the electromagnetic field, and the latter is prescribed through a vector potential. The frozen-field model is much simpler and has been used successfully for numerical simulations of vortex systems .
In this article, we prove that the frozen-field model is obtained as the asymptotic limit of the GL model when the charge of the Cooper pairs (the superconducting charge carriers) goes to zero while the applied magnetic field stays near the upper critical field. (The upper critical field itself depends on the charge of the Cooper pairs and increases as the latter decreases.) Because the temperature is constant in the GL model, this limit corresponds to fixing the temperature $`T`$ and moving up vertically through the vortex regime to the curve labeled $`H_{c2}`$ in the phase diagram of Fig. 1. The convergence rate is second order in the small parameter.
For more background on the physics of superconductivity we refer the reader to the monograph by Tinkham . The original source for the GL equations of superconductivity is . A good introduction to the mathematics of the GL equations is . The dynamics of the GL equations have been studied by several authors; see and the references cited therein. The present investigation is closely related to the work of Du and Gray .
Section 2 introduces the Ginzburg–Landau equations, Section 3 contains the numerical results and Section 4 the analysis.
## 2 The Ginzburg–Landau equations
In the Ginzburg–Landau theory of superconductivity, the state of a superconducting medium is described by a complex scalar-valued order parameter $`\psi `$ and a real vector-valued vector potential $`𝑨`$. If the state varies with time, a third variable—the electric potential $`\varphi `$—is necessary to fully describe the electromagnetic field. The evolution of the state variables is governed by the time-dependent Ginzburg–Landau (TDGL) equations,
$$\gamma \mathrm{}\left(\frac{}{t}+\mathrm{i}\frac{q_s}{\mathrm{}}\varphi \right)\psi +\frac{1}{2m_s}\left(\frac{\mathrm{}}{\mathrm{i}}\frac{q_s}{c}𝑨\right)^2\psi +\alpha \psi +\beta |\psi |^2\psi =0,$$
(2.1)
$$\sigma \left(\frac{1}{c}\frac{𝑨}{t}\varphi \right)\frac{c}{4\pi }\times \times 𝑨+𝑱_s+\frac{c}{4\pi }\times 𝑯=0,$$
(2.2)
where the supercurrent density $`𝑱_s`$ is a nonlinear function of $`\psi `$ and $`𝑨`$,
$$𝑱_s=\frac{q_s\mathrm{}}{2\mathrm{i}m_s}(\psi ^{}\psi \psi \psi ^{})\frac{q_s^2}{m_sc}|\psi |^2𝑨=\frac{q_s}{m_s}\mathrm{}\left[\psi ^{}\left(\frac{\mathrm{}}{\mathrm{i}}\frac{q_s}{c}𝑨\right)\psi \right].$$
(2.3)
These equations are supplemented by the boundary conditions,
$$𝒏𝑱_s=0,𝒏\times (\times 𝑨)=𝒏\times 𝑯.$$
(2.4)
Here, $`𝑯`$ is the applied magnetic field, which we assume to be time independent. The constants $`m_s`$ and $`q_s`$ are the mass and charge, respectively, of a Cooper pair (the superconducting charge carriers, also referred to as superelectrons); $`c`$ is the speed of light; and $`\mathrm{}`$ is Planck’s constant divided by $`2\pi `$. A Cooper pair is made up of two electrons, each with charge $`e`$ ($`e`$ is the elementary charge); hence, $`q_s`$ is negative, $`q_s=2e`$.
The parameters $`\alpha `$ and $`\beta `$ are material parameters; $`\alpha `$ changes sign at the critical temperature $`T_c`$, $`\alpha (T)<0`$ for $`T<T_c`$ (superconducting state) and $`\alpha (T)>0`$ for $`T>T_c`$ (normal state); $`\beta `$ is only weakly temperature dependent and positive for all $`T`$. The remaining parameters are $`\sigma `$, the normal state conductivity, and $`\gamma `$, the mobility coefficient. The latter is dimensionless and related to the diffusion coefficient $`D`$, $`\gamma =\mathrm{}/2m_sD`$.
The boundary conditions (2.4) express the fact that superelectrons cannot leave the superconductor. Also, if no surface currents are present, the tangential components of the magnetic field must be continuous across the boundary.
The parameters $`\alpha `$ and $`\beta `$ are defined phenomenologically, but they can be expressed in terms of measurable quantities, such as the superconducting coherence length $`\xi `$ and the London penetration depth $`\lambda `$,
$$\xi =\left(\frac{\mathrm{}^2}{2m_s|\alpha |}\right)^{1/2},\lambda =\left(\frac{m_sc^2\beta }{4\pi q_s^2|\alpha |}\right)^{1/2}.$$
(2.5)
The coherence length and the London penetration depth define the respective characteristic length scales for the order parameter and the magnetic induction. Both depend on the temperature $`T`$ and diverge as $`T`$ approaches the critical temperature $`T_c`$, because of the factor $`|\alpha |^{1/2}`$. However, their ratio is, to a good approximation, independent of temperature. This ratio is the Ginzburg-Landau parameter,
$$\kappa =\lambda /\xi .$$
(2.6)
In high-$`T_c`$ superconductors, $`\kappa `$ is of the order of 50–100.
The electromagnetic variables are the magnetic induction $`𝑩`$, the current density $`𝑱`$, and the electric field $`𝑬`$; they are given in terms of $`𝑨`$ and $`\varphi `$ by the expressions
$$𝑩=\times 𝑨,𝑱=\frac{c}{4\pi }\times \times 𝑨,𝑬=\frac{1}{c}\frac{𝑨}{t}\varphi .$$
(2.7)
Equation (2.2) is essentially Ampère’s law, $`𝑱=(c/4\pi )\times 𝑩`$, where the current $`𝑱`$ is the sum of the supercurrent $`𝑱_s`$, the transport current $`𝑱_t=(c/4\pi )\times 𝑯`$, and a “normal” current $`𝑱_n=\sigma 𝑬`$ (Ohm’s law). Hence, the GL model uses a quasistatic version of Maxwell’s equations, where the time derivative of the electric field is ignored.
The TDGL equations were first given by Schmid in 1966 and subsequently derived from the microscopic theory of superconductivity by Gor’kov and Eliashberg . Our notation is the same as in Gor’kov and Kopnin .
The solution of the TDGL equations is not unique. Any solution $`(\psi ,𝑨,\varphi )`$ defines a family of solutions $`G_\chi (\psi ,𝑨,\varphi )`$ parameterized by a sufficiently smooth function $`\chi `$ of space and time,
$$G_\chi :(\psi ,𝑨,\varphi )(\psi \text{e}^{i(q_s/\mathrm{}c)\chi },𝑨+\chi ,\varphi \frac{1}{c}\frac{\chi }{t}).$$
(2.8)
This property is known as gauge invariance; the function $`\chi `$ is known as a gauge function. Gauge invariance does not affect the physically measurable quantities (the magnetic induction $`𝑩`$, the magnetization $`𝑴=𝑩𝑯`$, and the current density $`𝑱`$). Uniqueness requires an additional constraint, which is imposed through a gauge choice. The choice of a proper gauge for the TDGL equations has been a subject of considerable debate. The choice is a matter of convenience and may depend on the problem under investigation. In this article we adopt a gauge in which, at any time, the electric potential and the divergence of the vector potential satisfy the identity
$$\sigma \varphi +(c/4\pi )𝑨=0$$
(2.9)
everywhere in the domain, while $`𝑨`$ is tangential at the boundary. This choice is realized by identifying the gauge $`\chi `$ with a solution of the linear parabolic equation
$$\frac{\sigma }{c}\frac{\chi }{t}\frac{c}{4\pi }\mathrm{\Delta }\chi =\sigma \varphi +\frac{c}{4\pi }𝑨,$$
(2.10)
subject to the condition $`𝒏\chi =𝒏𝑨`$ on the boundary. In , it was shown that the TDGL equations, subject to the constraint (2.9), define a dynamical system under suitable regularity conditions on $`𝑯`$. (In the more general case, where $`𝑯`$ varies not only in space but also in time, the TDGL equations define a dynamical process.) This dynamical system has a global attractor, which consists of the stationary points of the dynamical system and the heteroclinic orbits connecting such stationary points. Furthermore, it was shown that every solution on the attractor satisfies the condition $`𝑨=0`$ (and, therefore, also $`\varphi =0`$). Thus, in the limit as $`t\mathrm{}`$, every solution of the TDGL equations satisfies the GL equations in the London gauge.
### 2.1 Nondimensional TDGL equations
In this section, we render the TDGL equations dimensionless by choosing units for the independent and dependent variables. Since we are interested in the collective behavior of vortices in the bulk of a superconductor in the limit of weak coupling ($`q_s0`$), we take care to choose the units in such a way that they remain of order one as $`q_s0`$. (We recall that $`q_s`$ is negative, $`q_s=2e`$.)
As $`q_s0`$, the coherence length $`\xi `$ remains of order one, while the penetration depth $`\lambda `$ increases like $`|q_s|^1`$; see Eq. (2.5). This suggests taking the coherence length $`\xi `$ as the unit of length.
To maintain the diffusion coefficient $`D=\mathrm{}/2\gamma m_s=\xi ^2(\gamma \mathrm{}/|\alpha |)^1`$ at order one, we measure time in units of $`\gamma \mathrm{}/|\alpha |`$.
The real and imaginary parts of the order parameter are conveniently measured in units of $`\psi _0=(|\alpha |/\beta )^{1/2}`$, which is the value of $`\psi `$ that minimizes the free energy in the absence of a field.
Next, consider the magnetic field. A fundamental quantity in the theory of type-II superconductors is the flux quantum $`\mathrm{\Phi }_0`$,
$$\mathrm{\Phi }_0=\frac{hc}{|q_s|}=2\pi \frac{\mathrm{}c}{|q_s|}.$$
(2.11)
The flux quantum is the unit of magnetic flux carried by a vortex. Together with the coherence length and penetration depth, it defines three characteristic field strengths: the lower critical field $`H_{c1}`$, the thermodynamical critical field $`H_c`$, and the upper critical field $`H_{c2}`$,
$$H_{c1}=\frac{\mathrm{\Phi }_0}{4\pi \lambda ^2\mathrm{ln}\kappa },H_c=\frac{\mathrm{\Phi }_0}{2\pi \xi \lambda \sqrt{2}},H_{c2}=\frac{\mathrm{\Phi }_0}{2\pi \xi ^2}.$$
(2.12)
Below $`H_{c1}`$, a superconductor is in the ideal superconducting (Meissner) state, where it does not support magnetic flux in the bulk; above $`H_{c2}`$, it is in the normal state, where the magnetic flux is distributed uniformly in the bulk; between $`H_{c1}`$ and $`H_{c2}`$, it is in the vortex state, where magnetic flux is quantized in vortex-like configurations (see Fig. 1). The thermodynamical critical field $`H_c`$ is intermediate between $`H_{c1}`$ and $`H_{c2}`$ and is defined by the identity $`H_c^2/8\pi =\frac{1}{2}\psi _0^2|\alpha |`$; $`H^2/8\pi `$ is the energy per unit volume associated with a field $`H`$, and $`\frac{1}{2}\psi _0^2|\alpha |`$ is the minimum condensation energy, which is attained when $`\psi =\psi _0`$, so these two quantities are in balance when $`H=H_c`$.
As $`q_s0`$, $`H_{c1}`$ goes to 0 like $`|q_s|`$, $`H_c`$ remains of order one, and $`H_{c2}`$ grows like $`|q_s|^1`$. This suggests that we define field strengths in terms of $`H_c`$. In fact, it is convenient to absorb a factor $`\sqrt{2}`$, so we adopt $`H_c\sqrt{2}`$ or, equivalently, $`\mathrm{}c/\xi \lambda |q_s|`$ as the unit for the magnetic field strength.
With the coherence length as the unit of length and $`H_c\sqrt{2}`$ as the unit of field strength, it follows that the vector potential is measured in units of $`\xi H_c\sqrt{2}`$. Furthermore, energy densities are measured in units of $`H_c^2/4\pi `$, which is the same as $`|\alpha |\psi _0^2`$.
Finally, we define the scalar potential $`\varphi `$ in units of $`(1/\gamma \psi _0^2\kappa |q_s|)(H_c^2/4\pi )`$. Notice that this unit remains of order one as $`q_s0`$, because $`\kappa |q_s|`$ is of order one. On the other hand, the product $`q_s\varphi `$, which represents an energy density, tends to zero as $`q_s0`$. (It remains finite on the scale of the penetration depth.)
Table 1 summarizes the relations between the original variables and their nondimensional (primed) counterparts. We adopt the latter as the new variables and work until further notice on the nondimensional problem. We omit all primes.
The nondimensional TDGL equations are
$$\left(\frac{}{t}\frac{\mathrm{i}}{\kappa }\varphi \right)\psi \left(+\frac{\mathrm{i}}{\kappa }𝑨\right)^2\psi (1|\psi |^2)\psi =0,$$
(2.13)
$$\sigma \frac{𝑨}{t}\mathrm{\Delta }𝑨\frac{1}{\kappa }𝑱_s\times 𝑯=\mathrm{𝟎},$$
(2.14)
where
$$𝑱_s=\frac{1}{2\mathrm{i}}(\psi ^{}\psi \psi \psi ^{})\frac{1}{\kappa }|\psi |^2𝑨=\mathrm{}\left[\psi ^{}\left(+\frac{\mathrm{i}}{\kappa }𝑨\right)\psi \right],$$
(2.15)
with the corresponding gauge condition,
$$\sigma \varphi +𝑨=0.$$
(2.16)
In deriving Eq. (2.14), we have made use of the gauge condition (2.16) and the vector identity
$$\mathrm{\Delta }𝑨=\times \times 𝑨+(𝑨).$$
(2.17)
If $`\mathrm{\Omega }`$ is the domain occupied by the superconducting material (measured in units of $`\xi `$), then Eqs. (2.13)–(2.16) must be satisfied everywhere $`\mathrm{\Omega }`$. At the boundary $`\mathrm{\Omega }`$ of $`\mathrm{\Omega }`$, we have the conditions
$$𝒏𝑱_s=0,𝒏\times (\times 𝑨)=𝒏\times 𝑯,𝒏𝑨=0.$$
(2.18)
Here, $`𝒏`$ is the local unit normal vector.
The electromagnetic variables are given by the expressions
$$𝑩=\times 𝑨,𝑱=\times \times 𝑨,𝑬=_t𝑨\varphi .$$
(2.19)
The values of the lower and upper critical fields are
$$H_{c1}=(2\kappa \mathrm{ln}\kappa )^1,H_{c2}=\kappa .$$
(2.20)
### 2.2 Link variables
The combination $`+(\mathrm{i}/\kappa )𝑨`$ plays a fundamental role; we refer to it as the $`𝑨`$-gradient and write
$$_𝐀=+\frac{\mathrm{i}}{\kappa }𝑨.$$
(2.21)
The $`𝑨`$-gradient defines the $`𝑨`$-Laplacian (or “twisted Laplacian”),
$$\mathrm{\Delta }_𝐀=_𝐀_𝐀=\left(+\frac{\mathrm{i}}{\kappa }𝑨\right)^2.$$
(2.22)
The relation between the $`𝑨`$-Laplacian and the ordinary Laplacian is most easily illustrated by means of the link variables,
$`U_x(x,y,z)`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{\mathrm{i}}{\kappa }}{\displaystyle ^x}A_x(\xi ,y,z)d\xi \right),`$
$`U_y(x,y,z)`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{\mathrm{i}}{\kappa }}{\displaystyle ^y}A_y(x,\eta ,z)d\eta \right),`$ (2.23)
$`U_z(x,y,z)`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{\mathrm{i}}{\kappa }}{\displaystyle ^z}A_z(x,y,\zeta )d\zeta \right).`$
(We omit the argument $`t`$.) The integrals are evaluated with respect to an arbitrary reference point. Each $`U_\mu `$ ($`\mu =x,y,z`$) is complex valued and unimodular, $`U_\mu ^{}=U_\mu ^1`$. The vectors $`𝑨`$ and $`𝑼`$ may be used interchangeably. With a slight abuse of notation, we have
$$𝑼=\mathrm{e}^{(\mathrm{i}/\kappa ){\scriptscriptstyle 𝑨}},_𝐀=U^{}U,\mathrm{\Delta }_𝐀=U^{}\mathrm{\Delta }U.$$
(2.24)
## 3 Numerical solution
A parallel code for solving Eqs. (2.13)–(2.18) has been developed as part of a project for large-scale simulations of vortex dynamics in superconductors. Details on these simulations and on the code will be published elsewhere; here, we give only a brief overview of the numerical methods and the results of numerical simulations showing the behavior of the solution as $`\kappa `$ increases.
The algorithm uses finite differences on a staggered grid, making all approximations accurate to second order in the mesh widths, and an implicit method for the time integration, making the algorithm (essentially) unconditionally stable. The code, written in C++, has been designed for a multiprocessing environment; it uses MPI for message passing.
We restrict the discussion to rectangular two-dimensional configurations that are periodic in one direction and open in the other. The configurations are assumed to be infinite in the third, orthogonal direction, which is also the direction of the applied magnetic field, $`𝑯=(0,0,H_z)`$.
### 3.1 Discretization
#### Computational grid.
The computational grid is uniform, with equal mesh sizes in the $`x`$ and $`y`$ direction, $`h_x=h_y=h`$. A vertex on the grid is denoted by $`𝒙_{i,j}=(x_i,y_j)`$ and is the point of reference for the grid cell shown in Fig. 4.
The indices run through the values $`i=1,\mathrm{},N_x`$ and $`j=1,\mathrm{},N_y`$. We assume periodicity in the $`x`$ direction and take the grid so the vertices with $`j=1`$ and $`j=N_y`$ are located on the open boundary of the superconductor. Thus, the size of the domain is $`S=N_x(N_y1)h^2`$.
#### Variables.
The discrete variables are introduced so that all derivatives are given by second-order accurate central-difference approximations. The scalar variables $`\psi `$ and $`\varphi `$ are defined on the vertices of the grid,
$$\psi _{i,j}=\psi (𝒙_{i,j}),\varphi _{i,j}=\varphi (𝒙_{i,j}).$$
(3.1)
(We use the same symbol for the original field and its discrete counterpart.) Vectors are defined at the midpoints of the links connecting adjacent vertices,
$$A_{x;i,j}=A_x(𝒙_{i,j}+\frac{1}{2}h_x𝒆_x),A_{y;i,j}=A_y(𝒙_{i,j}+\frac{1}{2}h_y𝒆_y).$$
(3.2)
Here, $`𝒆_x`$ and $`𝒆_y`$ denote the unit vectors in the $`x`$ and $`y`$ direction, respectively. The definition of the discrete supercurrent $`𝑱_s`$ is completely analogous. The link variables, defined in Eq. (2.24), are obtained from the vector potential,
$$U_{x;i,j}=\mathrm{e}^{(\mathrm{i}/\kappa )A_{x;i,j}h_x},U_{y;i,j}=\mathrm{e}^{(\mathrm{i}/\kappa )A_{y;i,j}h_y}.$$
(3.3)
They are therefore also defined on the links. Finally, the magnetic induction $`𝑩`$, which is a vector perpendicular to the plane and given by the curl of the vector potential, is defined at the center of a grid cell,
$$B_{z;i,j}=B_z(𝒙_{i,j}+\frac{1}{2}h_x𝒆_x+\frac{1}{2}h_y𝒆_y).$$
(3.4)
The definition of the discrete variables is also illustrated in Fig. 4.
Note that, because of the location of the grid relative to the boundaries, all scalar variables, as well as the $`x`$ components of all vectors ($`A_x`$, $`U_x`$, $`J_{s,x}`$, and so forth), are defined on a $`N_x\times N_y`$ grid, whereas the $`y`$ components of all vectors and the magnetic induction $`B_z`$ are defined on a $`N_x\times (N_y1)`$ grid.
#### Boundary conditions.
We assume periodicity in the $`x`$ direction, so we need to consider the boundary conditions (4.7) only at $`y=y_1`$ and $`y=y_{N_y}`$.
The boundary condition for the order parameter, $`𝒏_𝐀\psi =0`$, becomes
$$U_{y;i,1}\psi _{i,2}\psi _{i,1}=0,\psi _{i,N_y}U_{y;i,N_y1}^{}\psi _{i,N_y1}=0,$$
(3.5)
for $`i=1,\mathrm{},N_x`$. For the vector potential, we require that $`_yA_x=H_z`$ and $`A_y`$ is constant ($`A_y=0`$) on the boundary.
#### Operators.
The gradient of a scalar is a vector and is therefore defined at the midpoint of a link connecting two adjacent vertices. Thus,
$$\left(\varphi \right)_{x;i,j}=(_x\varphi )(𝒙_{i,j}+\frac{1}{2}h_x𝒆_x)=h_x^1(\varphi _{i+1,j}\varphi _{i,j}),$$
(3.6)
with an analogous definition for the $`y`$ component. The gauge-invariant $`𝐀`$-gradient $`_𝐀=+\mathrm{i}𝑨`$ is defined in a similar way, with
$$\left(_𝐀\psi \right)_{x;i,j}=h_x^1(\psi _{i+1,j}U_{x;i,j}\psi _{i,j}).$$
(3.7)
Thus, the discrete version of the twisted Laplacian $`\mathrm{\Delta }_𝐀`$ is
$`\left(\mathrm{\Delta }_𝐀\psi \right)_{i,j}`$ $`=`$ $`h_x^2(\psi _{i+1,j}U_{x;i,j}2\psi _{i,j}+\psi _{i1,j}U_{x;i1,j}^{})`$ (3.8)
$`+h_y^2(\psi _{i,j+1}U_{y;i,j}2\psi _{i,j}+\psi _{i,j1}U_{x;i,j1}^{}).`$
The discrete version of the (normal) Laplacian is defined in the usual way,
$$\left(\mathrm{\Delta }\psi \right)_{i,j}=h_x^2(\psi _{i+1,j}2\psi _{i,j}+\psi _{i1,j})+h_y^2(\psi _{i,j+1}2\psi _{i,j}+\psi _{i,j1}).$$
(3.9)
The magnetic induction, which is the curl of the vector potential, takes the form
$$B_{z;i,j}=h_x^1(A_{y;i+1,j}A_{y;i,j})h_y^1(A_{x;i,j+1}A_{x;i,j}).$$
(3.10)
We also need the divergence of the vector potential, which is given by
$$\left(𝑨\right)_{i,j}=h_x^1(A_{x;i,j}A_{x;i1,j})+h_y^1(A_{y;i,j}A_{y;i,j1}).$$
(3.11)
#### Algorithm.
For numerical purposes, it is useful to treat the TDGL equations (4.3) and (4.4) as two separate equations, which are coupled only through certain fields and variables. The electromagnetic potentials $`\varphi `$ and $`𝑨`$ are treated as static variables in the order parameter equation,
$$\left(_t(\mathrm{i}/\kappa )\varphi \right)\psi \mathrm{\Delta }_𝐀\psi (1|\psi |^2)\psi =0.$$
(3.12)
The local nonlinear part of this equation,
$$\left(_t(\mathrm{i}/\kappa )\varphi \right)\psi (1|\psi |^2)\psi =0,$$
(3.13)
is integrated in the simplest possible manner,
$$\psi _{i,j}(t+\mathrm{\Delta }t)=\mathrm{e}^{(i/\kappa )\varphi _{i,j}\mathrm{\Delta }t}\left\{\psi _{i,j}(t)+\mathrm{\Delta }t\left(1|\psi _{i,j}|^2\right)\psi _{i,j}\right\}.$$
(3.14)
The nonlocal part,
$$_t\psi \mathrm{\Delta }_𝐀\psi =0,$$
(3.15)
is integrated by using a backward Euler method, where the linear equation system is solved with the conjugate gradient method.
The equation for the vector potential,
$$\sigma _t𝑨\mathrm{\Delta }𝑨(1/\kappa )𝑱_s\times 𝑯=\mathrm{𝟎},$$
(3.16)
is linear and depends only indirectly on the order parameter through the supercurrent. If we treat the supercurrent as a static variable, we can integrate the equation easily, again using the backward Euler method. In the actual implementation, we use the fact that the domain is periodic to do a fast Fourier transform in the $`x`$ direction, which leaves us with a tridiagonal system to solve in the $`y`$ direction. This procedure is considerably faster than using an iterative method, such as the conjugate gradient method.
### 3.2 Numerical results
We use a rectangular sample, periodic in the $`x`$ direction and open in the $`y`$ direction, with $`N_x=N_y=128`$. We take $`h_x=h_y=\frac{1}{2}\xi `$, so the sample measures 64 coherence lengths in the periodic direction and 63.5 coherence lengths across. (The coherence length $`\xi `$ is defined in Eq. (2.5).)
First, we considered this system with $`\kappa =200`$ and an applied magnetic field $`H_z=0.088\kappa `$. With a relatively large value of $`\kappa `$, the surface barrier for vortex entry is low, and the system equilibrates relatively fast . The equilibration required $`5\times 10^4`$ time steps with $`\mathrm{\Delta }t=0.4`$. The magnetic field produces an almost perfect vortex lattice. Figure 5 gives a contour plot of the density of Cooper pairs $`|\psi |^2`$ at equilibrium; the zeros correspond to the centers of the vortices.
We then started from the configuration of Fig. 5 to find equilibrium configurations for other values of $`\kappa `$, varying $`\kappa `$ from $`\kappa _{\mathrm{min}}=40`$ to $`\kappa _{\mathrm{max}}=800`$. In this range, the ground states are comparable and similar to the one shown in Fig. 5. Since the magnetization of a sample is proportional to $`1/\kappa ^2`$, the vortex density decreases with $`\kappa `$; below $`\kappa _{\mathrm{min}}`$, the equilibrium state has fewer vortices, and a comparison becomes meaningless. Each equilibration required another $`3\times 10^4`$ time steps.
Figure 6 gives the computed values of the quantities
$$\delta \psi =\psi _\kappa \psi _{\kappa _{\mathrm{max}}}_{L^2},\delta B_z=\frac{B_{z,\kappa }B_{z,\kappa _{\mathrm{max}}}_{L^2}}{H_z_{L^2}},$$
(3.17)
for different values of $`\kappa `$. The data show a behavior like $`1/\kappa ^2`$ down to $`\kappa 40`$.
Figure 7 shows the average over $`x`$ of $`A_{x,\kappa }A_{x,\kappa _{\mathrm{max}}}`$ as a function of $`y`$ in the bulk of the sample, for different values of $`\kappa `$.
The numerical results show that the solution of the TDGL equations converges as $`\kappa `$ increases; in fact, they show quadratic convergence in the small parameter $`1/\kappa `$. Given the fact that the Ginzburg–Landau parameter of high-$`T_c`$ superconducting materials is of the order of 50–100, we conclude that the limiting equation is a practical alternative in many applications. The question thus becomes: What is the limiting equation, and can we confirm the numerical conclusions by rigorous arguments? We address this question in the next section.
## 4 Asymptotic analysis
We now return to the TDGL equations (2.13)–(2.18) and consider their limit as $`\kappa \mathrm{}`$. These are our standing hypotheses:
$`\mathrm{\Omega }`$ is bounded in $`𝐑^n`$ ($`n=2,3`$), with a sufficiently smooth boundary $`\mathrm{\Omega }`$, for example, $`\mathrm{\Omega }`$ of class $`C^{1,1}`$.
The parameters $`\kappa `$ and $`\sigma `$ are real and positive.
$`𝑯`$ is independent of time; as a function of position, it satisfies the regularity condition $`𝑯[W^{\alpha ,2}(\mathrm{\Omega })]^n`$ for some $`\alpha (\frac{1}{2},1)`$.
$`\kappa 1`$; $`\sigma =O(1)`$ and $`𝑯=O(\kappa )`$ as $`\kappa \mathrm{}`$.
The assumptions (H1)–(H3) suffice to prove that the TDGL equations define a dynamical system in the Hilbert space
$$𝒲^{1+\alpha ,2}=[W^{1+\alpha ,2}(\mathrm{\Omega })]^2\times [W^{1+\alpha ,2}(\mathrm{\Omega })]^n;$$
(4.1)
see . The space $`W^{1+\alpha ,2}(\mathrm{\Omega })`$ is continuously imbedded in $`W^{1,2}(\mathrm{\Omega })L^{\mathrm{}}(\mathrm{\Omega })`$, so $`\psi `$ and $`𝑨`$ are bounded and differentiable with square-integrable (generalized) derivatives. (H4) is the operative hypothesis for the asymptotic analysis.
### 4.1 Mathematical analysis
#### Scaling.
We begin by scaling the TDGL equations, taking into account the fact that we are interested in the limit as $`q_s0`$ (weak coupling), when the applied field is near the upper critical field. The scaling is done by means of the dimensionless GL parameter $`\kappa `$, which grows like $`|q_s|^1`$.
Since $`𝑯=O(\kappa )`$ as $`\kappa \mathrm{}`$, we begin by scaling $`𝑯`$ by a factor $`\kappa `$, $`𝑯=\kappa 𝑯^{}`$. By scaling the vector potential by the same factor $`\kappa `$, we achieve that the electromagnetic variables are all of the same order.
The scalar potential is proportional to the charge density of the Cooper pairs, which is $`O(|q_s|)`$ as $`q_s0`$. Hence, $`\kappa \varphi `$ remains of order one. This suggests scaling $`\varphi `$ by a factor $`\kappa ^1`$.
Table 2 summarizes the relation between the current (nondimensional) variables and their scaled (primed) counterparts. We adopt the latter as the new variables and work until further notice on the scaled problem. We omit all primes.
After scaling, the relevant parameter is $`\kappa ^2`$, rather than $`\kappa `$, so we introduce $`\epsilon `$,
$$\epsilon =\kappa ^2.$$
(4.2)
The scaled TDGL equations are
$$\left(_t\mathrm{i}\epsilon \varphi \right)\psi (+\mathrm{i}𝑨)^2\psi (1|\psi |^2)\psi =0,$$
(4.3)
$$\sigma _t𝑨\mathrm{\Delta }𝑨\epsilon 𝑱_s\times 𝑯=\mathrm{𝟎},$$
(4.4)
where
$$𝑱_s=\frac{1}{2\mathrm{i}}(\psi ^{}\psi \psi \psi ^{})|\psi |^2𝑨=\mathrm{}\left[\psi ^{}(+\mathrm{i}𝑨)\psi \right],$$
(4.5)
with the corresponding gauge condition,
$$\epsilon \sigma \varphi +𝑨=0.$$
(4.6)
The boundary conditions associated with Eqs. (4.3) and (4.4) are
$$𝒏(+\mathrm{i}𝑨)\psi =0,𝒏\times (\times 𝑨)=𝒏\times 𝑯,𝒏𝑨=0.$$
(4.7)
The electromagnetic variables are given by the expressions
$$𝑩=\times 𝑨,𝑱=\times \times 𝑨,𝑬=_t𝑨\epsilon \varphi .$$
(4.8)
#### Reduction to homogeneous form.
Next, we homogenize the problem. Let $`𝑨_0`$ be the (unique) minimizer of the convex quadratic form $`Q_1Q_1[𝑨]`$,
$$Q_1[𝑨]=_\mathrm{\Omega }\left[(𝑨)^2+|\times 𝑨𝑯|^2\right]d𝒙,$$
(4.9)
on dom$`(Q_1)=\{𝑨[W^{1,2}(\mathrm{\Omega })]^n:𝒏𝑨=0\text{ on }\mathrm{\Omega }\}`$. This minimizer satifies the boundary-value problem
$$\times \times 𝑨\times 𝑯=\mathrm{𝟎},𝑨=0\text{in }\mathrm{\Omega },$$
(4.10)
$$𝒏\times (\times 𝑨)=𝒏\times 𝑯,𝒏𝑨=0\text{on }\mathrm{\Omega },$$
(4.11)
in the dual of dom$`(Q_1)`$ with respect to the inner product in $`[L^2(\mathrm{\Omega })]^n`$. The mapping $`𝑯𝑨_0`$ is linear, time independent, and continuous from $`[W^{\alpha ,2}(\mathrm{\Omega })]^n`$ to $`[W^{1+\alpha ,2}(\mathrm{\Omega })]^n`$ . The contribution of the vector $`𝑨_0`$ to the magnetic field is
$$𝑩_0=\times 𝑨_0.$$
(4.12)
We substitute variables,
$$𝑨=𝑨_0+\epsilon 𝑨^{},$$
(4.13)
and rewrite the (scaled) TDGL equations (4.3)–(4.7) in terms of $`\psi `$ and $`𝑨^{}`$ (omitting the primes),
$$_t\psi +\mathrm{i}\sigma ^1((\epsilon 𝑨))\psi (+\mathrm{i}(𝑨_0+\epsilon 𝑨))^2\psi (1|\psi |^2)\psi =0\text{in }\mathrm{\Omega },$$
(4.14)
$$\sigma _t𝑨\mathrm{\Delta }𝑨𝑱_s=\mathrm{𝟎}\text{in }\mathrm{\Omega },$$
(4.15)
where
$$𝑱_s=\frac{1}{2\mathrm{i}}(\psi ^{}\psi \psi \psi ^{})|\psi |^2(𝑨_0+\epsilon 𝑨),$$
(4.16)
and
$$𝒏\psi =0,𝒏\times (\times 𝑨)=\mathrm{𝟎},𝒏𝑨=0\text{on }\mathrm{\Omega }.$$
(4.17)
#### Functional formulation.
We reformulate the system of Eqs. (4.14)–(4.17) as an ordinary differential equation for a vector-valued function $`u=(\psi ,𝑨)`$ from the time domain $`(0,\mathrm{})`$ to a space of functions on $`\mathrm{\Omega }`$,
$$u=(\psi ,𝑨):[0,\mathrm{})^2=[L^2(\mathrm{\Omega })]^2\times [L^2(\mathrm{\Omega })]^n.$$
(4.18)
The equation is
$$\frac{\mathrm{d}u}{\mathrm{d}t}+Au=f_0(u)+\epsilon f_1(u),$$
(4.19)
where $`A`$ is the linear operator in $`^2`$ associated with the quadratic form $`QQ[u]`$,
$$Q[u]=_\mathrm{\Omega }\left[|\psi |^2+\sigma ^1\left((𝑨)^2+|\times 𝑨|^2\right)\right]d𝒙,$$
(4.20)
on dom$`(Q)=\{u=(\psi ,𝑨)^2:𝒏𝑨=0\text{ on }\mathrm{\Omega }\}`$. The functions $`f_0`$ and $`f_1`$ are nonlinear,
$$f_i(u)=(\phi _i(\psi ,𝑨),\sigma ^1𝑭_i(\psi ,𝑨)),i=0,1,$$
(4.21)
where
$`\phi _0(\psi ,𝑨)`$ $`=`$ $`2\mathrm{i}𝑨_0(\psi )|𝑨_0|^2\psi +(1|\psi |^2)\psi ,`$ (4.22)
$`\phi _1(\psi ,𝑨)`$ $`=`$ $`\mathrm{i}(1\sigma ^1)(𝑨)\psi +2\mathrm{i}𝑨(\psi )(𝑨_0𝑨)\psi |𝑨|^2\psi ,`$ (4.23)
$`𝑭_0(\psi ,𝑨)`$ $`=`$ $`0,`$ (4.24)
$`𝑭_1(\psi ,𝑨)`$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{i}}}(\psi ^{}\psi \psi \psi ^{})|\psi |^2(𝑨_0+\epsilon 𝑨).`$ (4.25)
Given any $`f=(\phi ,\sigma ^1𝑭)^2`$, the equation $`Au=f`$ is equivalent with the system of uncoupled boundary-value problems
$`\mathrm{\Delta }\psi =\phi \text{ in }\mathrm{\Omega },`$ $`𝒏\psi =0\text{ on }\mathrm{\Omega },`$ (4.26)
$`\mathrm{\Delta }𝑨=𝑭\text{ in }\mathrm{\Omega },`$ $`𝒏\times 𝑨=\mathrm{𝟎},𝒏𝑨=0\text{ on }\mathrm{\Omega },`$ (4.27)
in the dual of dom$`(Q)`$ with respect to the inner product in $`^2`$. The operator $`A`$ is selfadjoint and positive definite in $`^2`$; hence, its fractional powers $`A^{\theta /2}`$ are well defined, they are unbounded if $`\theta 0`$, and dom$`(A^{\theta /2})`$ is a closed linear subspace of $`𝒲^{\theta ,2}=[W^{\theta ,2}(\mathrm{\Omega })]^2\times [W^{\theta ,2}(\mathrm{\Omega })]^n`$; see \[17, Section 1.4\].
The solution of Eq. (4.19) depends on $`\epsilon `$; we denote it by $`u_\epsilon `$. We compare $`u_\epsilon `$ with the solution $`u_0`$ of the reduced equation
$$\frac{\mathrm{d}u}{\mathrm{d}t}+Au=f_0(u).$$
(4.28)
###### Theorem 4.1
There exists a positive constant $`C`$ such that
$$u_\epsilon (t)u_0(t)_{𝒲^{1+\alpha ,2}}C\left(u_\epsilon (0)u_0(0)_{𝒲^{1+\alpha ,2}}+\epsilon \right),t[0,T].$$
(4.29)
Proof. Let $`B_R`$ be the ball of radius $`R`$ centered at the origin in $`𝒲^{1+\alpha ,2}`$. Let $`u_\epsilon B_R`$ and $`u_0B_R`$ satisfy Eqs. (4.19) and (4.28), respectively, with initial data $`u_\epsilon (0)`$ and $`u_0(0)`$. The difference $`v=u_\epsilon u_0`$ satisfies the differential equation
$$\frac{\mathrm{d}v}{\mathrm{d}t}+Av=f_0(u_\epsilon )f_0(u_0)+\epsilon f_1(u_\epsilon )$$
(4.30)
or, equivalently, the integral equation
$$v(t)=\mathrm{e}^{tA}v(0)+_0^t\mathrm{e}^{(ts)A}[f_0(u_\epsilon )f_0(u_0)+\epsilon f_1(u_\epsilon )](s)ds.$$
(4.31)
From the integral equation we obtain the estimate
$`v(t)_{𝒲^{1+\alpha ,2}}`$ $``$ $`\mathrm{e}^{tA}_{𝒲^{1+\alpha ,2}}v(0)_{𝒲^{1+\alpha ,2}}+{\displaystyle _0^t}A^{(1+\alpha )/2}\mathrm{e}^{(ts)A}_{𝒲^{1+\alpha ,2}}`$ (4.32)
$`\times \left[f_0(u_\epsilon )f_0(u_0)_{L^2}+\epsilon f_1(u_\epsilon )_{L^2}\right](s)\mathrm{d}s.`$
The operator norms satisfy the inequalities
$$\mathrm{e}^{tA}_{𝒲^{1+\alpha ,2}}1,A^{(1+\alpha )/2}\mathrm{e}^{(ts)A}_{𝒲^{1+\alpha ,2}}C(ts)^{(1+\alpha )/2};$$
(4.33)
see \[17, Theorem 1.4.3\]. Furthermore, adding and subtracting terms, we have
$`f_0(u_\epsilon )f_0(u_0)`$ $`=`$ $`(2i𝑨_0((\psi _\epsilon \psi _0))|𝑨_0|^2(\psi _\epsilon \psi _0)`$ (4.34)
$`+(1|\psi _\epsilon |^2|\psi _0|^2)(\psi _\epsilon \psi _0)\psi _\epsilon \psi _0(\psi _\epsilon ^{}\psi _0^{}),0),`$
where
$$2i𝑨_0((\psi _\epsilon \psi _0))_{L^2}2𝑨_0_L^{\mathrm{}}\psi _\epsilon \psi _0_{W^{1,2}}$$
$$C\psi _\epsilon \psi _0_{W^{1+\alpha ,2}}Cu_\epsilon u_0_{𝒲^{1+\alpha ,2}},$$
$$|𝑨_0|^2(\psi _\epsilon \psi _0)_{L^2}C𝑨_0_L^{\mathrm{}}^2\psi _\epsilon \psi _0_L^{\mathrm{}}$$
$$C\psi _\epsilon \psi _0_{W^{1+\alpha ,2}}Cu_\epsilon u_0_{𝒲^{1+\alpha ,2}},$$
and the other terms are estimated similarly. Here, $`C`$ is some (generic) positive constant, which may depend on $`𝑯`$ and $`\mathrm{\Omega }`$ but not on $`u_\epsilon `$ or $`u_0`$. (In these inequalities we have used the continuity of the imbedding of $`W^{1+\alpha ,2}`$ into $`W^{1,2}L^{\mathrm{}}`$.) The result is an inequality of the type
$$f_0(u_\epsilon )f_0(u_0)_{L^2}Cu_\epsilon u_0_{𝒲^{1+\alpha ,2}},$$
(4.35)
showing that $`f_0`$ is Lipschitz from $`𝒲^{1+\alpha ,2}`$ to $`^2`$.
Using similar estimates, we show that $`f_1`$ is bounded from $`𝒲^{1+\alpha ,2}`$ to $`^2`$, so there exists a positive constant $`C`$ such that
$$f_1(u_\epsilon )_{L^2}C.$$
(4.36)
Combining the estimates (4.33), (4.35), and (4.36) with the inequality (4.32), we conclude that there exist positive constants $`C_1`$ and $`C_2`$ such that
$$v(t)_{𝒲^{1+\alpha ,2}}v(0)_{𝒲^{1+\alpha ,2}}+\epsilon C_1t^{(1\alpha )/2}+C_2_0^t(ts)^{(1+\alpha )/2}v(s)_{𝒲^{1+\alpha ,2}}ds.$$
(4.37)
Applying Gronwall’s inequality, we obtain the estimate
$$v(t)_{𝒲^{1+\alpha ,2}}C\left(v(0)_{𝒲^{1+\alpha ,2}}+\epsilon \right),t[0,T],$$
(4.38)
for some positive constant $`C`$.
It follows from Theorem 4.1 that, if the initial data are such that $`u_\epsilon (0)u_0(0)_{𝒲^{1+\alpha ,2}}=o(1)`$ as $`\epsilon 0`$, then
$$\underset{\epsilon 0}{lim}u_\epsilon =u_0$$
(4.39)
in $`C([0,T];𝒲^{1+\alpha ,2})`$ for any $`T>0`$. In particular, if $`u_\epsilon (0)u_0(0)_{𝒲^{1+\alpha ,2}}=O(\epsilon )`$, then the convergence in Eq. (4.39) is $`O(\epsilon )`$.
### 4.2 Interpretation and final remarks
It remains to translate the results back in terms of the original variables. We denote the solution of the TDGL equations, Eqs. (2.13)–(2.18), by $`\psi _\kappa `$, $`𝑨_\kappa `$, $`\varphi _\kappa `$. The variables $`𝑨_\kappa `$ and $`\varphi _\kappa `$ are related by the gauge condition $`\sigma \varphi _\kappa +𝑨_\kappa =0`$ at all times. Let $`𝑩_\kappa =\times 𝑨_\kappa `$.
Let $`𝑨_{\mathrm{}}`$ be the solution of the boundary-value problem
$$\times \times 𝑨\times 𝑯=\mathrm{𝟎},𝑨=0\text{in }\mathrm{\Omega },$$
(4.40)
$$𝒏\times (\times 𝑨)=𝒏\times 𝑯,𝒏𝑨=0\text{on }\mathrm{\Omega },$$
(4.41)
and put $`𝑩_{\mathrm{}}=\times 𝑨_{\mathrm{}}`$. The vector $`𝑨_{\mathrm{}}`$ and, hence, $`𝑩_{\mathrm{}}`$ do not vary with time. Let $`\psi _{\mathrm{}}`$ satisfy the equations
$$_t\psi \mathrm{\Delta }_𝐀_{\mathrm{}}\psi (1|\psi |^2)\psi =0\text{ in }\mathrm{\Omega },𝒏_{𝐀_0}\psi =0\text{ on }\mathrm{\Omega }.$$
(4.42)
Then it follows from Theorem 4.1 that there exists a positive constant $`C`$ such that
$`\psi _\kappa (t)\psi _{\mathrm{}}(t)_{W^{1+\alpha ,2}}+{\displaystyle \frac{𝑩_\kappa (t)𝑩_{\mathrm{}}_{W^{\alpha ,2}}}{𝑯_{W^{\alpha ,2}}}}`$ (4.44)
$`C\left(\psi _\kappa (0)\psi _{\mathrm{}}(0)_{W^{1+\alpha ,2}}+{\displaystyle \frac{𝑩_\kappa (0)𝑩_{\mathrm{}}_{W^{\alpha ,2}}}{𝑯_{W^{\alpha ,2}}}}+{\displaystyle \frac{1}{\kappa ^2}}\right),`$
for all $`t[0,T]`$, $`T>0`$.
The approximation $`(\psi _{\mathrm{}},𝑩_{\mathrm{}})`$ is the “frozen-field approximation.” Hence, the analysis shows that the solution of the TDGL equations converges to the frozen-field approximation, uniformly on compact time intervals $`[0,T]`$ in the topology of $`[W^{1+\alpha ,2}(\mathrm{\Omega })]^2\times [W^{\alpha ,2}(\mathrm{\Omega })]^n`$, as soon as the initial data satisfy the asymptotic estimates $`\psi _\kappa (0)\psi _{\mathrm{}}(0)_{W^{1+\alpha ,2}}=o(1)`$ and $`𝑩_\kappa (0)𝑩_{\mathrm{}}_{W^{\alpha ,2}}=o(\kappa )`$ as $`\kappa \mathrm{}`$. Under slightly sharper conditions we obtain the order of convergence.
###### Corollary 4.1
If
$$\psi _\kappa (0)\psi _{\mathrm{}}(0)_{W^{1+\alpha ,2}}=O\left(\frac{1}{\kappa ^2}\right)\text{and}\frac{𝑩_\kappa (0)𝑩_{\mathrm{}}_{W^{\alpha ,2}}}{𝑯_{W^{\alpha ,2}}}=O\left(\frac{1}{\kappa ^2}\right)$$
as $`\kappa \mathrm{}`$, then
$$\psi _\kappa (t)\psi _{\mathrm{}}(t)_{W^{1+\alpha ,2}}+\frac{𝑩_\kappa (t)𝑩_{\mathrm{}}_{W^{\alpha ,2}}}{𝑯_{W^{\alpha ,2}}}=O\left(\frac{1}{\kappa ^2}\right),$$
(4.45)
uniformly on compact intervals.
This result explains the numerical results presented in Section 3.
#### Remark 1.
The asymptotic approximation procedure can be continued to higher order, as can be seen from a formal expansion. The equations for the order parameter and the vector potential decouple, and at each order one finds first the vector potential, then the order parameter. The vector potential satisfies a linear heat equation; for example, the first correction beyond $`𝑨_{\mathrm{}}`$ is $`\kappa ^1𝑨`$, where $`𝑨`$ satisfies the equation
$$\sigma _t𝑨+\mathrm{\Delta }𝑨=\mathrm{}\left[\psi _{\mathrm{}}^{}_𝐀_{\mathrm{}}\psi _{\mathrm{}}\right].$$
(4.46)
#### Remark 2.
The analysis given here differs at several points from the analysis of Ref. . First, our scaling is slightly different and, we believe, more in tune with the physics; second, our regularity assumptions on the applied field are weaker; third, our proofs are more direct; and fourth, our results hold in a stronger topology.
### Acknowledgments
We thank Professor Todd Dupont (University of Chicago) for stimulating discussions throughout the course of this investigation. We also acknowledge the work of Damien Declat (student from ISTIL, Lyon, France), who assisted in the development of an earlier version of the parallel computer program.
This work was supported by the Mathematical, Information, and Computational Sciences Division subprogram of the Office of Advanced Scientific Computing Research, U.S. Department of Energy, under Contract W-31-109-Eng-38. The second author was partially supported by the University of Chicago/Argonne National Laboratory Collaborative Grant No. 96-011.
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# On the Fairlie’s Moyal formulation of M(atrix)-theory
## 1 Introduction
Matrix model formulation of M-theory, introduced three years ago by Banks et al , and studied later by many authors\[2-3\], have shown to be an important area of research. In this model, recall for instance that the fundamental degrees of freedom of the theory are the D0-branes whose interaction processus are described by space-time matrix coordinates. Actually, we know that this matrix model is described by the maximally supersymmetric U(N)Yang-Mills gauge theory, where N is the light-like momentum or the number of D0-branes interpreted in as the number of Green Schwartz strings (gas of N strings) in the light cone gauge. Its also known that these strings have differents lenghts depending on the winding number $`m_k`$ around the compact direction, a fact which can be described by the following expression
$$N=m_kN_k.$$
(1)
The large N limit is then shown to corresponds to long strings in the infinite momentum frame (IFM). Several ways to investigate this large N limit have been followed. Among these issues one focus on the work of Fairlie which purposes to clarify the meaning of this infinite limit in 4 and 10 dimensions in terms of Moyal bracket formalism. As claimed by this author, this Moyal bracket description is motivated by the fact that the matrix theory takes on an aspect analogous to a Moyal formulation of quantum mechanics. Note by the way, that the Moyal bracket is defined via the star product known to be an essential ingredient towards setting a non commutative geometry framework \[5-6\]. Recently this old mathematical idea have been a subject of a revived interest in the study of string and matrix theory \[7-8-9-10\].
Before going into the description of the main lines of the present work, we will try to summarise the Fairlie’s Moyal bracket formulation in M-theory as done in . As signalled above and in accordance with , the author examine the Moyal version of matrix theory in the belief that this is the most appropriate for the discussion of the large N limit, and for the investigation of parallels with quantum mechanics. Indeed, it is consistently shown that the Moyal bracket description of infinite limit of matrix theory in 2 and 8 transverse dimensions leads to a class of solutions for the second order matrix theory equations obtained from the first order ones.
Using this correspondence between the second and the first order equations, the author of proceed then by solving the original theory by inverting the problem which means finding a solution to first order equations which guarantees the integrability of the second order ones. This remarkable feature is due to the existence of the Bogomol’nyi bound to the Euclideanised version of the above equations in 4 and 10 dimensions. The general construction of the obtained solutions is also given in terms of a representation of the target space coordinates as non local spinor bilinears, which are generalisations of the standard wigner functions on phase space.
In the present work, we consider an alternative way to reexamine the matrix theory equations of motion in 10 dimensions in the infinite limit, formulated in terms of Poisson bracket instead of Moyal bracket. Among our results we rewrite in a first step the Nahm’s equations into a simple form. This consist precisely in reducing the 14 coupled Nahm’s equations (11,12) leading just to five ones (13,14). These reduced equations are convenient in the sense that their derivative with respect to $`\sigma `$ gives exactly the second order equations of motion (16) which we derive starting from the $`SU(\mathrm{})`$ pure Yang Mills lagrangian density. We show also that the second order equations of motion exhibit a special invariance property as given in (21).
The next step consist to understand the meaning of the first order Nahm’s equations and their relation with the equations of motion. For this, we give an interpretation of the former as the vanishing condition of some conserved currents that we consider and interpret the above correspondence as the conservation property of these currents. Finally, we propose to solve the derived second order equations of motion. We have accomplish this task by considering the Fourier modes decomposition, using some algebraic manipulations and setting the important ansatz (48) leading to build our membrane solution.
Next, we try to look for the possible connection of our solution with other solutions known in literature. We start first from the important remark that our membrane solution (71) exhibits an oscillation behaviour which can be reduced to give rise to other solutions (see for example ) once particular choices of our parameters are performed.
We organise this paper as follows. In section 2, we review the large N limit in terms of Moyal bracket as done in and present our alternative approach in section 3 to derive the second order membrane equations of motion in 10 dimensions. In section 4, we give a breif comment about Nahm’s equations and present in section 5 the analysis leading to construct a membrane solution in 10 dimensions. Section 6 is devoted to our conclusion.
## 2 Moyal bracket and large N limit
Since matrix model formulation of M-theory was introduced by Banks et al , a lot of stimulated papers elaborating the issue was done \[2-3\]. Focusing for instance one of these papers , when the author describes the large N limit of matrix theory in 4 and 10 dimensions using the Moyal bracket formalism. The starting point is the matrix theory with an action of a two dimensional $`N=8`$ supersymmetric U(N) Yang-Mills theory namely
$$S=\frac{1}{2\pi \alpha ^{}}Tr((D_\mu X)^2+\theta ^T\mathit{}\theta +g_{s}^{}{}_{}{}^{2}F_{\mu \nu }^{}{}_{}{}^{2}\frac{1}{g_{s}^{}{}_{}{}^{2}}[X^i,X^j]^2+\frac{1}{g_s}\theta ^T\gamma _i[X^i,\theta ]d\sigma d\tau ,$$
(2)
where $`X^i`$ and $`(\theta _L^\alpha ,\theta _R^\alpha )`$ are the 8 scalar and 8 fermionic fields respectively, which are $`N\times N`$ hermetian matrices transforming as the 8v vector and (8s, 8c) spinor representations of the SO(8) R-symmetry group of transversal rotations.
The coordinates $`\sigma `$ and $`\tau `$ with $`0\sigma 2\pi `$ parametrise a cylinder (world sheet).
The passage to the large N limit consist in a first step in considering the matrices $`X^i`$ as functions of two phase space variables $`\alpha ,\beta `$ as well as $`\sigma ,\tau `$ such that $`X^iX^i(\alpha ,\beta ,\sigma ,\tau )`$ . The associated matrix elements may be regarded as the Fourier components of $`X^i`$ . The second step consist in substituting commutators of the action (2) by the Moyal bracket in the following way
$$Tr[X^i,X^j]𝑑\sigma \frac{1}{\lambda ^2}𝑑\sigma 𝑑\alpha 𝑑\beta \{X^i,X^j\}_{MB}^2,$$
(3)
where $`\{X^i,X^j\}_{MB}=\mathrm{sin}\{X^i,X^j\}`$ is the sine or Moyal bracket, with deformation parameter $`\lambda `$ defined as the imaginary part of the star product \*. Recall that the star product of two functions $`X^i`$ and $`X^j`$ is defined as:
$$X^iX^j=\underset{\stackrel{\alpha ^{}\alpha }{\beta ^{}\beta }}{lim}\mathrm{exp}^{i\lambda (_\alpha _\beta ^{}_\alpha ^{}_\beta )}X^i(\alpha ,\beta ,\sigma )X^j(\alpha ^{},\beta ^{},\sigma ).$$
(4)
The Moyal bracket is then defined as the antisymmetric part of the star product such that
$$\{X^i,X^j\}_{MB}=\underset{\stackrel{\alpha ^{}\alpha }{\beta ^{}\beta }}{lim}\mathrm{sin}\lambda (_\alpha _\beta ^{}_\alpha ^{}_\beta )X^i(\alpha ,\beta ,\sigma )X^j(\alpha ^{},\beta ^{},\sigma ).$$
(5)
As quoted in , the point of this construction is that in the limiting points $`\lambda \frac{2\pi }{N}`$, the Moyal bracket (5) reproduces the commutators of $`N\times N`$ matrices $`X^i`$(bosonic coordinates) through the association of the components $`X_{mn}^i`$ of $`X^i`$ with the Fourier modes of a function $`X^i(\alpha ,\beta ,\sigma )`$ periodic in $`\alpha ,\beta `$ . The remaining terms in the action (2) involving fermionic coordinates are replaced by
$$\begin{array}{ccc}\{X^\mu ,\theta \}_{MB}\hfill & =& \hfill \underset{\stackrel{\alpha ^{}\alpha }{\beta ^{}\beta }}{lim}\mathrm{sin}\lambda (_\alpha _\beta ^{}_\alpha ^{}_\beta )X^i(\alpha ,\beta ,\sigma )\theta (\alpha ^{},\beta ^{},\sigma )\\ \{\theta ^T,\mathit{}\theta \}_{MB}\hfill & =& \hfill \underset{\stackrel{\alpha ^{}\alpha }{\beta ^{}\beta }}{lim}\mathrm{cos}\lambda (_\alpha _\beta ^{}_\alpha ^{}_\beta )\theta ^T(\alpha ,\beta ,\sigma )\mathit{}\theta (\alpha ^{},\beta ^{},\sigma ).\end{array}$$
(6)
Thus the action becomes
$$\begin{array}{ccc}S_{MB}=\frac{1}{2\pi \alpha ^{}}((D_\mu X)^2+\mathrm{cos}\{\theta ^T,\mathit{}\theta \}+g_{s}^{}{}_{}{}^{2}TrF_{\mu \nu }^{}{}_{}{}^{2})𝑑\alpha 𝑑\beta 𝑑\sigma 𝑑\tau \hfill & & \\ ((\frac{1}{\lambda g_{s}^{}{}_{}{}^{2}}\mathrm{sin}\{X^i,X^j\})^2\frac{1}{g_s}\mathrm{cos}\{\theta ^T\gamma _i,\frac{1}{\lambda }\mathrm{sin}\{X^i,\theta \}\})𝑑\alpha 𝑑\beta 𝑑\sigma 𝑑\tau .\hfill & & \end{array}$$
(7)
Moreover, once the following correspondence $`\lambda \frac{2\pi }{N}`$ is performed, the final form of the action , in the large N-limit ($`\lambda 0`$), is then expressed in terms of ordinary Poisson brackets as
$$\begin{array}{ccc}S_{PB}\hfill & =& \hfill \frac{1}{2\pi \alpha ^{}}((D_\mu X)^2+\theta ^T\mathit{}\theta +g_{s}^{}{}_{}{}^{2}TrF_{\mu \nu }^{}{}_{}{}^{2})𝑑\alpha 𝑑\beta 𝑑\sigma 𝑑\tau \\ & & \hfill ((\frac{1}{g_{s}^{}{}_{}{}^{2}}\{X^i,X^j\})^2\frac{1}{g_s}\theta ^T\gamma _i\{X^i,\theta \})𝑑\alpha 𝑑\beta 𝑑\sigma 𝑑\tau .\end{array}$$
(8)
The obtained action defines a membrane SU($`\mathrm{}`$) pure Yang Mills theory. Later on, we will consider the longitudinal as well as the timelike coordinates in order to treat the system dynamically .
## 3 Membranes in 10 dimensions: an alternative approach
The aim of this section is to derive the equations of motion associated with the SU($`\mathrm{}`$) Yang-Mills theory describing the membrane in 8 transverse directions and show how these second order equations are connected to the first order Nahm’s equations. The latter’s are shown to play a central role as they provide a way to emphasise the similarity to the phase space formulation of quantum mechanics .
A remarkable feature of the situation with 8 transverse directions is that, due to the existence of a self-dual (antisymmetric) 4-tensor $`T_{\mu \nu \rho \sigma }`$ in 10 dimensions, the theory admits a class of solutions which we can obtain from a first order formulation. The Lagrangian density describing the 10-dimensional membrane have then the form of an SU($`\mathrm{}`$) pure Yang-Mills theory (8).
The situation is further simplified to one of dependence upon only one of the variables $`\sigma `$ (apart from the $`\alpha ,\beta `$ dependence of the gauge potential $`X^i(\alpha ,\beta ,\sigma )`$ ). Using the gauge choice $`X^0=constant`$, the Lagrangian density is simply written as
$$\begin{array}{ccc}L=\frac{1}{2}(_\sigma X^\mu )^2+\frac{1}{4}\{X^\mu ,X^\nu \}^2,\mu =1,\mathrm{}\mathrm{..9},\hfill & & \end{array}$$
(9)
where $`_\sigma =/\sigma `$ and where the dependence upon $`\tau `$ is ignored. The symbol$`\{X^\mu ,X^\nu \}`$ deals with the Poisson bracket given by
$$\begin{array}{ccc}\{X^\mu ,X^\nu \}=_\alpha X^\mu _\beta X^\nu _\beta X^\mu _\alpha X^\nu .\hfill & & \end{array}$$
(10)
Note by the way that (9) describes a 10-dimensional membrane with $`X^0=constant`$, living inside 8 transverse directions. This situation is motivated by the fact that the matrix theory admits a class of solutions obtainable from a first order formulation (the Nahm equations). This is a set of 9 first order differential equations with 6 constraint equations given by
$$\begin{array}{ccc}_\sigma X^1+\{X^2,X^9\}=0\hfill & & \\ _\sigma X^2+\{X^9,X^1\}=0\hfill & & \\ _\sigma X^3+\{X^4,X^9\}=0\hfill & & \\ _\sigma X^4+\{X^9,X^3\}=0\hfill & & \\ _\sigma X^5+\{X^6,X^9\}=0\hfill & & \\ _\sigma X^6+\{X^9,X^5\}=0\hfill & & \\ _\sigma X^7+\{X^8,X^9\}=0\hfill & & \\ _\sigma X^8+\{X^9,X^7\}=0\hfill & & \\ _\sigma X^9+\{X^1,X^2\}+\{X^3,X^4\}+\{X^5,X^6\}+\{X^7,X^8\}=0,\hfill & & \end{array}$$
(11)
with
$$\begin{array}{ccc}\{X^1,X^3\}+\{X^4,X^2\}+\{X^5,X^7\}+\{X^8,X^6\}=0\hfill & & \\ \{X^1,X^4\}+\{X^2,X^3\}+\{X^8,X^5\}+\{X^7,X^6\}=0\hfill & & \\ \{X^1,X^5\}+\{X^4,X^8\}+\{X^7,X^3\}+\{X^6,X^2\}=0\hfill & & \\ \{X^1,X^6\}+\{X^2,X^5\}+\{X^3,X^8\}+\{X^4,X^7\}=0\hfill & & \\ \{X^1,X^7\}+\{X^3,X^5\}+\{X^8,X^2\}+\{X^6,X^4\}=0\hfill & & \\ \{X^1,X^8\}+\{X^5,X^4\}+\{X^2,X^7\}+\{X^6,X^3\}=0.\hfill & & \end{array}$$
(12)
Focusing for the moment to reduce much more these equations, we have used some algebraic manipulations and showed that (11-12) can be simply written in the following way
$$\begin{array}{ccc}_\sigma X^\mu +()^{\mu +1}\{X^{\mu +()^{\mu +1}},X^9\}=0,\mu =1,\mathrm{},8\hfill & (a)& \\ _\sigma X^9+_{\mu =1}^8\frac{1}{2}()^{\mu +1}\{X^\mu ,X^{\mu +()^{\mu +1}}\}=0,\hfill & (b)& \end{array}$$
(13)
and
$$\begin{array}{ccc}\{X^1,X^j\}+()^{ȷ+1}\{X^{ȷ+()^{ȷ+1}},X^2\}=\{X^7,X^{ȷ+2()^{ȷ+1}}\}+()^{ȷ+1}\{X^{ȷ+3()^{ȷ+1}},X^8\},j=3,6\hfill & & \\ \{X^1,X^j\}+()^{ȷ+1}\{X^{ȷ+()^{ȷ+1}},X^2\}=\{X^5,X^{ȷ4()^{ȷ+1}}\}+()^{ȷ+1}\{X^{ȷ3()^{ȷ+1}},X^6\},j=4,7\hfill & & \\ \{X^1,X^j\}+()^{ȷ+1}\{X^{ȷ+()^{ȷ+1}},X^2\}=\{X^3,X^{ȷ+2()^{ȷ+1}}\}+()^{ȷ+1}\{X^{ȷ+3()^{ȷ+1}},X^4\},j=5,8\hfill & & \end{array}$$
(14)
Indeed, setting for example $`\mu =1,2`$ , we recover respectively from (13-a) the first two equations of (11) namely
$$\begin{array}{ccc}_\sigma X^1+\{X^2,X^9\}=0\hfill & & \\ _\sigma X^2+\{X^9,X^1\}=0.\hfill & & \end{array}$$
(15)
Next note the important remark that there exist an intriguing correspondence between the first order Nahm’s equations (13) and the second order derived equations of motion. Indeed consider the lagrangian (9); and applying the Euler Lagrange equations we find the equations of motion for the membrane, namely:
$$\begin{array}{ccc}_\sigma ^2X^\mu +_{\nu =1}^8\{X^\nu ,\{X^\nu ,X^\mu \}\}=0,\hfill & & \end{array}$$
(16)
with $`\mu =1,\mathrm{},9`$ and $`X_0=constant`$. The above correspondence is then established by derivating the Nahm’s equations with respect to the coordinate $`\sigma `$ and using the constraints (14). As an example consider (13.b) such that
$$\begin{array}{ccc}_\sigma ^2X^9+\underset{\nu =1}{\overset{8}{}}\frac{1}{2}()^{\nu +1}(\{_\sigma X^\nu ,X^{\nu +()^{\nu +1}}\}+\{X^\nu ,_\sigma X^{\nu +()^{\nu +1}}\})=0,\hfill & & \end{array}$$
(17)
which becomes upon using (13.a)
$$\begin{array}{ccc}_\sigma ^2X^9+\underset{\nu =1}{\overset{8}{}}\frac{1}{2}(\{X^{\nu +()^{\nu +1}},\{X^{\nu +()^{\nu +1}},X^9\}\}+\{X^\nu ,\{X^\nu ,X^9\}\})=0,\hfill & & \end{array}$$
(18)
and reproducing then exactly the equations of motion (16) for $`\mu =9`$ with the following algebraic property
$$\begin{array}{ccc}\underset{\nu =1}{\overset{8}{}}\{X^{\nu +()^{\nu +1}},\{X^{\nu +()^{\nu +1}},X^9\}\}=\underset{\nu =1}{\overset{8}{}}\{X^\nu ,\{X^\nu ,X^9\}\}.\hfill & & \end{array}$$
(19)
Note by the way that by virtue of this equality, the equations of motion (16) can be equivalently written as:
$$\begin{array}{ccc}_\sigma ^2X^\mu +\underset{\nu =1}{\overset{8}{}}\{X^{\nu +()^{\nu +1}},\{X^{\nu +()^{\nu +1}},X^\mu \}\}=0,\hfill & & \end{array}$$
(20)
which shows an invariance property with respect to the following transformation
$$\nu \nu +()^{\nu +1}\nu +\mathrm{exp}^{i\pi (\nu +1)},$$
(21)
affecting the repeated index $`\nu `$ .
Consequently we note that the equation of motion concerning the longitudinal coordinates $`X^9`$ is simply obtained when derivating with respect to $`\sigma `$. The remaining second order equations of motion, for transverse directions $`\mu =1,\mathrm{},8`$, arise thanks to the existence of constraint equations.
## 4 About the Nahm’s equations
Having shown explicitly how the second order equations of motion are derived from the Nahm’s first order ones, we will try now to search for the meaning of the relation between these two kind of equations. For this task, we will not ignore for instance about the coordinate $`\tau `$ and assume that our fields $`X^\mu `$ depend on the full set of coordinates $`\alpha ,\beta ,\sigma `$ and $`\tau `$ defining a space containing the world volume of the membrane.
This space can be structured as follows. Let $`\sigma `$ and $`\tau `$ define a complex two-dimensional world-sheet $`\mathrm{\Sigma }`$ with local coordinates $`Z=\sigma +i\tau `$ and $`\overline{Z}=\sigma i\tau `$ .
We use $``$ and $`\overline{}`$ to denote $`/z`$ and $`/\overline{z}`$ respectively. The coordinates $`\alpha ,\beta `$ parametrise a phase space which we denote by $`P(\alpha ,\beta )`$ such that the fields $`X^\mu `$ are shown to live inside the space $`\mathrm{\Sigma }P(\alpha ,\beta )`$.
The lagrangian describing the SU($`\mathrm{}`$) Yang-Mills theory of the membrane can be written formally as
$$\begin{array}{ccc}L=\frac{1}{2}(X^\mu )(\overline{}X^\mu )+\frac{1}{4}\{X^\mu ,X^\nu \}^2,\mu =1,\mathrm{..9},\hfill & & \end{array}$$
(22)
in such a way that for $`\tau =0`$ , one recover directly the standard theory (9) for which $`=\overline{}=_\sigma `$ .
The equation of motion associated to the lagrangian (22) is given by
$$\begin{array}{ccc}\overline{}X^\mu +_{\nu =1}^8\{X^\nu ,\{X^\nu ,X^\mu \}\}=0.\hfill & & \end{array}$$
(23)
This equation looks like a standard 2d conformal field theory equation of motion, for which we are usually interested in deriving the conserved currents, discussing the underlying conformal symmetries and integrability. A part from being interesting for the above reasons, our equations contains further informations concerning the structure of the membrane in 8 transverse directions and the ”matrix” behaviour of the fields $`X^\mu `$.
Using our knowledge on conformal fields theories and integrable systems for which 2d conserved currents are defined such that their conserved law reproduces in some sense the equation of motion, we are for instance interested in cheking what happens for our equation (23). In fact, we remark from our previous analysis that its possible to associate conserved currents to (23). Denoting these objects by $`J^\mu (J^i,J^9),i=1,\mathrm{},8`$; we can write
$$\begin{array}{ccc}J^i=X^i+()^{i+1}\{X^{i+()^{i+1}},X^9\},i=1,\mathrm{},8\hfill & (a)& \\ J^9=X^9+_{\nu =1}^8\frac{1}{2}()^{\nu +1}\{X^\nu ,X^{\nu +()^{\nu +1}}\},\hfill & (b)& \end{array}$$
(24)
and their conservation properties are
$$\begin{array}{ccc}\overline{}J^i=\overline{}X^i+_{\nu =1}^8\{X^\nu ,\{X^\nu ,X^i\}\}=0,i=1,\mathrm{},8\hfill & & \end{array}$$
(25)
and
$$\begin{array}{ccc}\overline{}J^9=\overline{}X^9+_{\nu =1}^8\{X^\nu ,\{X^\nu ,X^9\}\}=0,\hfill & & \end{array}$$
(26)
by virtue of (23).
Now, the point is that the conserved currents $`J^\mu ,\mu =1,\mathrm{}9`$ given by (24) are nothing but the objects defining the first order Nahm’s equations (13) which we can rewrite as follows
$$\begin{array}{ccc}J^\mu =0,\mu =1,\mathrm{..8}\hfill & & \\ J^9=0.\hfill & & \end{array}$$
(27)
This property can be traced to the fact that when the coordinate $`\tau `$ is ignored, which means setting simply $`\tau =0`$, the currents $`J^\mu (\mu =1,\mathrm{}9`$) become then vanishing objects and give rise then to the first order Nahm’s equations (13). We can then interpret these kind of equations, as the vanishing property of the conserved currents of the theory (22) and interpret the conservation law equations (25, 26) as the property giving the correspondence between the first and the second order equations as discussed previousely.
## 5 Solving the equations of motion
In this section we will return back to our discussion of section 3 in which we ignore the parameter $`\tau `$ and consider the following equations of motion(16)
$$\begin{array}{ccc}_\sigma ^2X^\mu +\underset{\nu =1}{\overset{8}{}}\{X^\nu ,\{X^\nu ,X^\mu \}\}=0.\hfill & & \end{array}$$
(28)
Now, having shown how these equations describe the membrane in 8 transverse directions, we are now in position to look for the explicit solution of the model. To start, we assume that the coordinates $`X^\mu (\alpha ,\beta ,\sigma )`$ can be written in terms of the Fourier modes as follows
$$\begin{array}{ccc}X^\mu (\alpha ,\beta ,\sigma )=\underset{mnZ}{}X_{mn}^\mu (\sigma )L_{mn}(\alpha ,\beta ),\hfill & & \end{array}$$
(29)
which suppose the periodicity of $`X^\mu (\alpha ,\beta ,\sigma )`$ in $`\alpha ,\beta `$
$$\begin{array}{ccc}X^\mu (\alpha +2\pi ,\beta +2\pi ,\sigma )=X^\mu (\alpha ,\beta ,\sigma ),\hfill & & \end{array}$$
(30)
with
$$\begin{array}{ccc}L_{mn}(\alpha ,\beta )=\mathrm{exp}i(m\alpha +n\beta ).\hfill & & \end{array}$$
(31)
Furthermore, the modes $`X_{mn}^\mu `$ as we will show bellow, can be considered as operator entries of the coordinates $`X^\mu `$ satisfying
$$\begin{array}{ccc}\{X^\mu ,X^\nu \}=\underset{\stackrel{m_1,n_1}{m_2,n_2}}{}X_{m_1n_1}^\mu (\sigma )X_{m_2n_2}^\nu (\sigma )\{L_{m_1n_1},L_{m_2n_2}\},\hfill & & \end{array}$$
(32)
where
$$\begin{array}{ccc}\{L_{m_1n_1},L_{m_2n_2}\}=(m_2n_1m_1n_2)L_{m_1+m_2,n_1+n_2},\hfill & & \end{array}$$
(33)
which reproduces in some sense the Poisson Bracket algebra in the large N limit of the area preserving diffeomorphism on the torus . Note by the way that (33) share a striking resemblance with the Antoniadis et al Lie algebra . This is an infinite dimensional algebra which was generalised to include the Virasoro algebra, the Frappat et al symmetries as well as their W and central charges extensions .
An important question about (33) is to look for the meaning of the corresponding generalisations in our context. Moreover, the appearence of this algebra indicates also how the SU(N) symmetry (Moyal) of the supersymmetric Yang-Mills matrix theory becomes the area preserving diffeomorphism group namely $`SU(\mathrm{})`$ describing the membrane.
Now, before solving explicitly our problem, we will present some algebraic properties related to (28). To start, consider the antisymmetry property of the Poisson bracket (32) namely
$$\begin{array}{ccc}\{X^\mu ,X^\nu \}=\{X^\nu ,X^\mu \},\hfill & & \end{array}$$
(34)
from which we can write
$$\begin{array}{ccc}\{X^\mu ,X^\nu \}\hfill & =& \hfill \underset{\stackrel{m_1,n_1}{m_2,n_2}}{}X_{m_1n_1}^\mu (\sigma )X_{m_2n_2}^\nu (\sigma )\{L_{m_1n_1},L_{m_2n_2}\}\\ & =& \hfill \underset{\stackrel{m_1,n_1}{m_2,n_2}}{}X_{m_1n_1}^\nu (\sigma )X_{m_2n_2}^\mu (\sigma )\{L_{m_1n_1},L_{m_2n_2}\}.\end{array}$$
(35)
On the other hand
$$\begin{array}{ccc}\{X^\mu ,X^\nu \}\hfill & =& \hfill \underset{\stackrel{m_1,n_1}{m_2,n_2}}{}X_{m_1n_1}^\mu (\sigma )X_{m_2n_2}^\nu (\sigma )\{L_{m_1n_1},L_{m_2n_2}\}\\ & =& \hfill \underset{\stackrel{m_1,n_1}{m_2,n_2}}{}X_{m_1n_1}^\mu (\sigma )X_{m_2n_2}^\nu (\sigma )\{L_{m_2n_2},L_{m_1n_1}\}\\ & =& \hfill \underset{\stackrel{m_1,n_1}{m_2,n_2}}{}X_{m_2n_2}^\mu (\sigma )X_{m_1n_1}^\nu (\sigma )\{L_{m_1n_1},L_{m_2n_2}\}.\end{array}$$
(36)
The above formulas, lead then to set the following results
$$\begin{array}{ccc}X_{m_1n_1}^\mu (\sigma )X_{m_2n_2}^\nu (\sigma )=X_{m_1n_1}^\nu (\sigma )X_{m_2n_2}^\mu (\sigma ),\hfill & & \end{array}$$
(37)
$$\begin{array}{ccc}X_{m_1n_1}^\mu (\sigma )X_{m_2n_2}^\nu (\sigma )=X_{m_2n_2}^\mu (\sigma )X_{m_1n_1}^\nu (\sigma ),\hfill & & \end{array}$$
(38)
showing the antisymmetry feature of the components $`X_{mn}^\nu `$ with respect to both the indices $`\mu ,\nu `$ and the combined index $`(m_in_i)`$, i=1,2 once the following mapping are performed
$$\begin{array}{ccc}(m_1n_1)\hfill & & \hfill (m_2n_2)\\ \mu \hfill & & \hfill \nu .\end{array}$$
(39)
Indeed, combining (37) and (38) we obtain
$$\begin{array}{ccc}X_{m_1n_1}^\mu (\sigma )X_{m_2n_2}^\nu (\sigma )=X_{m_2n_2}^\nu (\sigma )X_{m_1n_1}^\mu (\sigma ).\hfill & & \end{array}$$
(40)
This result shows explicitly the bosonic behaviour of the coordinates $`X_{mn}^\nu `$ , a natural property which is suspected already at the level of the Fourier modes decomposition (29). Indeed, the latter’s suppose from the begining that the coordinates $`X^\nu `$ are ”functions” of two phase space variables $`\alpha ,\beta `$ as well as $`\sigma `$.
In the spirit to understand much more the equations of motion and their underlying symmetries, we remark that a realisation of the component $`X_{mn}^\nu `$ is provided by the following expression
$$\begin{array}{ccc}X_{mn}^\mu (\sigma )=\gamma ^\mu X_{mn}(\sigma ),\hfill & & \end{array}$$
(41)
where $`\gamma ^\mu `$, $`\mu =1,\mathrm{},9`$ are the gamma matrices satisfying the Clifford algebra
$$\begin{array}{ccc}\gamma ^\mu \gamma ^\nu =\gamma ^\nu \gamma ^\mu \hfill & & \\ \gamma _{}^{\mu }{}_{}{}^{2}=1.\hfill & & \end{array}$$
(42)
Injecting for example (41) into (40), one find by virtue of (42) a non commutative property of the component $`X_{mn}`$ namely:
$$\begin{array}{ccc}X_{m_1n_1}(\sigma )X_{m_2n_2}(\sigma )=X_{m_2n_2}(\sigma )X_{m_1n_1}(\sigma ).\hfill & & \end{array}$$
(43)
In summary we can write (29) as follows
$$\begin{array}{ccc}X^\mu (\alpha ,\beta ,\sigma )=\underset{m,nϵZ}{}\gamma ^\mu X_{mn}(\sigma )L_{mn}(\alpha ,\beta ),\hfill & & \end{array}$$
(44)
which gives after some computations
$$\begin{array}{ccc}\underset{\nu =1}{\overset{8}{}}\{\{X^\mu ,X^\nu \},X^\nu \}=\underset{\nu =1}{\overset{8}{}}\underset{\stackrel{m_i,n_i}{i=1,2,3}}{}X_{m_1n_1}^\mu (\sigma )X_{m_2n_2}^\nu (\sigma )X_{m_3n_3}^\nu (\sigma )\{\{L_{m_1n_1},L_{m_2n_2}\},L_{m_3n_3}\},\hfill & & \end{array}$$
(45)
with
$$\begin{array}{ccc}\{\{L_{m_1n_1},L_{m_2n_2}\},L_{m_3n_3}\}=(m_2n_1m_1n_2)[m_3(n_1+n_2)n_3(m_1+m_2)]L_{\underset{i=1,2,3}{}m_i,\underset{i=1,2,3}{}n_i},\hfill & & \end{array}$$
(46)
and
$$\begin{array}{ccc}\underset{\nu =1}{\overset{8}{}}X_{m_1n_1}^\mu (\sigma )X_{m_2n_2}^\nu (\sigma )X_{m_3n_3}^\nu (\sigma )=8\gamma ^\mu X_{m_1n_1}(\sigma )X_{m_2n_2}(\sigma )X_{m_3n_3}(\sigma ).\hfill & & \end{array}$$
(47)
Recall that our aim is to solve the membrane equations of motion which consist to give an explicit solution for the coordinates $`X_{mn}`$ satisfying the non commutative rule (43). In order to solve the above problem we propose the following ansatz.
$$\begin{array}{ccc}X_{m_1n_1}(\sigma )X_{m_2n_2}(\sigma )X_{m_3n_3}(\sigma )=ϵ_{123}F_{sym}(\{m_i\},\{n_i\})X_{m_1+m_2+m_3,n_1+n_2+n_3}(\sigma ),\hfill & & \end{array}$$
(48)
where
$$\begin{array}{ccc}ϵ_{123}=ϵ_{(m_1n_1)(m_2n_2)(m_3n_3)},\hfill & & \end{array}$$
(49)
stands for the Levi-Cevita tensor given by
$$ϵ_{123}=\{\begin{array}{cc}\text{-1}\hfill & \text{for odd permutation}\hfill \\ \text{+1}\hfill & \text{for even permutation .}\hfill \end{array}$$
(50)
and where $`F_{sym}(\{m_i\},\{n_i\})`$ is for the moment an arbitrary function required to be symmetric with respect to the following permutation of integer values
$$\begin{array}{ccc}(m_in_i)(m_jn_j),i,j=1,2,3.\hfill & & \end{array}$$
(51)
Later on, we will denote simply this symmetric function as $`F_{(123)}`$.
The principal idea in setting the ansatz (48) is based on the fact that we need to write the bi-Poisson bracket $`\{\{X^\mu ,X^\nu \},X^\nu \}`$ as a simple function of $`X_{mn}`$ a fact which means that we should linearise the cubic matrix product $`X_{m_1n_1}X_{m_2n_2}X_{m_3n_3}`$ to give rise to (48). The apparition of the fully antisymmetric tensor $`ϵ_{123}`$ in this ansatz is justified by the non commutative behaviour of the components $`X_{mn}`$ as given by (43). On the other hand, the function $`F_{sym}`$ introduced in this ansatz and which we will discuss later, is chosen to be symmetric in agreement with the above non commutativity property.
Actually with this ansatz; one can easily identify both the left and hand sides terms of (28). In order to do this, remark first that we have
$$\begin{array}{ccc}\underset{\nu =1}{\overset{8}{}}\{\{X^\mu ,X^\nu \},X^\nu \}=\hfill & & \\ \underset{\stackrel{m_i,n_i}{i=1,2,3}}{}(\gamma ^\mu )8ϵ_{123}F_{(123)}\omega (\{m_i\},\{n_i\})X_{\underset{i=1,2,3}{}m_i,\underset{i=1,2,3}{}n_i}(\sigma )L_{\underset{i=1,2,3}{}m_i,\underset{i=1,2,3}{}n_i},\hfill & & \end{array}$$
(52)
with
$$\begin{array}{ccc}\omega (\{m_i\},\{n_i\})=(m_2n_1m_1n_2)[m_3\underset{i=1}{\overset{3}{}}n_in_3\underset{i=1}{\overset{3}{}}m_i].\hfill & & \end{array}$$
(53)
which leads then to write the membrane equations of motion as follows
$$\begin{array}{ccc}\underset{p,q}{}L_{p,q}_\sigma ^2X_{p,q}=\underset{\stackrel{m_i,n_i}{i=1,2,3}}{}8ϵ_{123}F_{(123)}\omega (\{m_i\},\{n_i\})X_{{\scriptscriptstyle m_i},{\scriptscriptstyle n_i}}(\sigma )L_{{\scriptscriptstyle m_i},{\scriptscriptstyle n_i}}.\hfill & & \end{array}$$
(54)
Now, in order to do identification in (54) in a consistent way, one should sum over the same indices $`p,q`$ in both the sides of this equation. This is possible, since we can set $`p=m_i`$ and $`q=n_i`$ or equivalently
$$\begin{array}{ccc}m_3=p(m_1+m_2)\hfill & & \\ n_3=q(n_1+n_2).\hfill & & \end{array}$$
(55)
This leads to consider the following transformations
$$\begin{array}{ccc}\underset{\stackrel{m_in_i}{i=1,2,3}}{}\hfill & & \hfill \underset{pq}{}\underset{\stackrel{m_i,n_i}{i=1,2}}{}\\ \omega (\{m_i\},\{n_i\})_{i=1,2,3}\hfill & & \hfill \omega (\{m_i\},\{n_i\},p,q)_{i=1,2}=\stackrel{~}{\omega }\\ ϵ_{123}\hfill & & \hfill ϵ_{12\overline{3}}\stackrel{~}{ϵ}\\ F_{(123)}\hfill & & \hfill F_{(12\overline{3})}=\stackrel{~}{F}_{sym},\end{array}$$
(56)
where $`ϵ_{12\overline{3}}`$ is the Levi-Cevita tensor introduced previousely and which is given by
$$\begin{array}{ccc}ϵ_{12\overline{3}}=ϵ_{(m_1n_1)(m_2n_2)(pq)}.\hfill & & \end{array}$$
(57)
With these simple transformations, the equations of motion (28) become
$$\begin{array}{ccc}_\sigma ^2X_{pq}=\underset{\stackrel{m_i,n_i}{i=1,2}}{}8\stackrel{~}{ϵ}\stackrel{~}{\omega }\stackrel{~}{F}_{sym}X_{pq}(\sigma ),\hfill & & \end{array}$$
(58)
with
$$\begin{array}{ccc}\stackrel{~}{\omega }=(m_2n_1m_1n_2)[p(n_1+n_2)q(m_1+m_2)].\hfill & & \end{array}$$
(59)
Note by the way that, as $`\stackrel{~}{\omega }`$ is antisymmetric with respect to the following mapping
$$\begin{array}{ccc}(m_1,n_1)\hfill & & \hfill (m_2,n_2),\end{array}$$
(60)
$`ϵ_{12\overline{3}}\stackrel{~}{\omega }`$ remains invariant with respect to the above change of integers values $`(m_i,n_i),i=1,2`$. We have
$$\begin{array}{ccc}ϵ_{12\overline{3}}\omega (\{m_i\},\{n_i\},p,q)ϵ_{21\overline{3}}\omega (m_1m_2,n_1n_2,p,q)\hfill & & \end{array}$$
(61)
since
$$\begin{array}{ccc}ϵ_{12\overline{3}}=ϵ_{21\overline{3}}\hfill & & \\ \omega (m_i,n_i,p,q)=\omega (m_1m_2,n_1n_2,p,q)\hfill & & \end{array}$$
(62)
Next, denoting by $`\mathrm{\Omega }`$ the following function
$$\mathrm{\Omega }(p,q)=\underset{\stackrel{m_i,n_i}{i=1,2}}{}8\stackrel{~}{ϵ}\stackrel{~}{\omega }\stackrel{~}{F}_{sym},$$
(63)
whose convergence is closely connected to how we can choice the function $`\stackrel{~}{F}_{sym}`$.
In fact the introduction of this symmetric function in our ansatz is really a crucial step as we can easily check that for arbitrary values of $`p,q`$ the following function
$$\begin{array}{ccc}\mathrm{\Omega }^{}(p,q)=\underset{\stackrel{m_i,n_i}{i=1,2}}{}8\stackrel{~}{ϵ}\stackrel{~}{\omega },\hfill & & \end{array}$$
(64)
always diverge. To avoid this divergency problem one should then consider $`\stackrel{~}{F}_{sym}`$ to play the role of the regularisation function which justify in some sense the introduction of this function in our ansatz (48). Following this discussion, a natural choice of this regularisation function is given by
$$\begin{array}{ccc}\stackrel{~}{F}_{sym}=F_{(12\overline{3})}=exp(|f_{12\overline{3}}|),\hfill & & \end{array}$$
(65)
where $`f_{12\overline{3}}`$=$`f(\{m_i\},\{n_i\},p,q)`$ is some function of the integer values $`m_i,n_i,p,q`$ required to be either symmetric or antisymmetric with respect to permutations of indices $`1(m_1n_1)`$, $`2(m_2n_2)`$ and $`\overline{3}(pq)`$.
The previous choice of $`\stackrel{~}{F}`$ as given in (65) is motivated by our requirement that the function $`\mathrm{\Omega }(p,q)`$ (63) should converge in the infinte limit of the integer values $`m_i,n_i,i=1,2`$.
Several realisations of the function $`f_{(12\overline{3})}`$ can occur. We give herebellow a typical example namely:
a) The symmetric choice
$$\begin{array}{ccc}f_{(12\overline{3})}=\underset{i=1}{\overset{2}{}}(m_i+n_i)+p+q.\hfill & & \end{array},$$
(66)
b) The antiymmetric choice
$$\begin{array}{ccc}f_{(12\overline{3})}=ϵ_{12\overline{3}}(\underset{i=1}{\overset{2}{}}(m_i+n_i)+p+q).\hfill & & \end{array}$$
(67)
The equation of motion (58) reads
$$\begin{array}{ccc}_\sigma ^2X_{pq}\mathrm{\Omega }X_{pq}=0,\hfill & & \end{array}$$
(68)
whose solution is shown to be
$$\begin{array}{ccc}X_{pq}(\sigma )=A_{pq}e^{\sqrt{\mathrm{\Omega }}\sigma }+B_{pq}e^{\sqrt{\mathrm{\Omega }}\sigma },\hfill & & \end{array}$$
(69)
where $`A_{pq},B_{pq}`$ are arbitrary constants which should satisfy the following relations
$$\begin{array}{ccc}A_{pq}A_{rs}=A_{rs}A_{pq}\hfill & & \\ B_{pq}B_{rs}=B_{rs}B_{pq}\hfill & & \\ A_{pq}B_{rs}+B_{rs}A_{pq}=(A_{rs}B_{pq}+B_{pq}A_{rs}),\hfill & & \end{array}$$
(70)
originated from the non commutativity properties of $`X_{pq}`$ (43). Finally the matrix model membrane solution read as
$$\begin{array}{ccc}X^\mu (\alpha ,\beta ,\sigma )=\underset{pq}{}\gamma ^\mu (A_{pq}e^{\sqrt{\mathrm{\Omega }}\sigma }+B_{pq}e^{\sqrt{\mathrm{\Omega }}\sigma })L_{pq}.\hfill & & \end{array}$$
(71)
The following significant question is in order: how one can compare or connect our solution with those already found in the same context in \[15-16\]?
One way to do this, is to discuss at the level of the derived solution (71) some particular cases related to the constants $`A_{pq}`$ , $`B_{pq}`$ and the regularised number $`\mathrm{\Omega }(p,q)`$.
$`\mathrm{𝟏}.𝛀>\mathrm{𝟎}`$
a)$`A_{pq}=B_{pq}`$
$$\begin{array}{ccc}X^\mu (\alpha ,\beta ,\sigma )=\underset{pq}{}2\gamma ^\mu A_{pq}ch(\sqrt{\mathrm{\Omega }}\sigma )e^{i(p\alpha +q\beta )}.\hfill & & \end{array}$$
(72)
b) $`A_{pq}=B_{pq}`$
$$\begin{array}{ccc}X^\mu (\alpha ,\beta ,\sigma )=\underset{pq}{}2\gamma ^\mu A_{pq}sh(\sqrt{\mathrm{\Omega }}\sigma )e^{i(p\alpha +q\beta )}.\hfill & & \end{array}$$
(73)
$`\mathrm{𝟐}.𝛀<\mathrm{𝟎}`$
a) $`A_{pq}=B_{pq}`$
$$\begin{array}{ccc}X^\mu (\alpha ,\beta ,\sigma )=\underset{pq}{}2\gamma ^\mu A_{pq}\mathrm{cos}(\sqrt{\mathrm{\Omega }}\sigma )e^{i(p\alpha +q\beta )}.\hfill & & \end{array}$$
(74)
b) $`A_{pq}=B_{pq}`$
$$\begin{array}{ccc}X^\mu (\alpha ,\beta ,\sigma )=\underset{pq}{}2i\gamma ^\mu A_{pq}\mathrm{sin}(\sqrt{\mathrm{\Omega }}\sigma )e^{i(p\alpha +q\beta )}.\hfill & & \end{array}$$
(75)
The above examples correspond then to typical membrane solutions whose oscillation behaviour is the same for all the values of the space time index $`\mu `$, with $`\mu =1,\mathrm{},9`$ once the signe of $`\mathrm{\Omega }`$ as well as the values of $`A_{pq}`$ and $`B_{pq}`$ are fixed.
Note also that we can consider from our solution (71), other examples for which the oscillation behaviour change when changing the index $`\mu `$ as done by the author of . We can then conclude that the solution we have derived is important in the sense that it exhibits among others an oscillation behaviour having a striking resemblance with the solution presented for example in . The claim is to remark that the coordinates p,q and t used for example in the Kim’s work, coincide respectively with $`\alpha ,\beta `$ and $`\sigma `$ in our construction.
## 6 Conclusion
This paper focus to solve some non-linear differential equations decribing the membrane in 10 dimensions by means of Poisson bracket formalism.
Among the results obtained, we derive the equations of motion associated with $`SU(\mathrm{})`$ Yang Mills theory describing the membrane in 8 transverse directions and show how the second order equations of motion are related to the first order ones (the Nahm’s equations). This provides a way to emphasise the similarity to the phase space formulation of quantum mechanics as signalled by the author of .
We rewrite the Nahm’s equations in a simple form reducing then their number from 14 to 5. This is convenient in the sense that their derivative with respect to $`\sigma `$ gives in a natural way the equations of motion.
We interpret the Nahm’s equations as the vanishing property of some conserved currents associated with the $`SU(\mathrm{})`$ Yang Mills theory and the correspondence with the equations of motion as the conservation property of these currents. The parameter $`\tau `$ is shown to play an important role in this sense.
To solve our problem, which means find an explicit expession for the fields $`X^\mu `$, we develop an algebraic analysis and propose an ansatz leading to construct the membrane solution given in (71).
In the spirit to compare our solution to some of the well known ones, one performe special choices on the constants $`A_{pq},B_{pq}`$ as well as on the regularised number $`\mathrm{\Omega }(p,q)`$ (63) a fact which leads to recover a general oscillation behaviour (see our previous examples) shared by our solution and the other ones already established in literature.
Moreover, we guess that other non trivial solutions can be obtained if one forget about the ansatz (48) and know how to solve in general the non-linear differential equation (28).
Acknowledgements
Two of the authors M.Hssaini and M.B.Sedra would like to thank the Abdus Salam International Centre for Theoretical Physics and the considerable help of the High Energy Section where this work was accomplished. The authors are especially grateful to Prof. K.S. Narain for reading the manuscript and for his important remarks. We thank also M. Alif Postdoc at the mathematics section at ICTP for useful conversations. Finally we thank the PARS program: Phys 27.372/98 CNR, for scientific help.
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# The influence of interactions and minor mergers on the structure of galactic disks Based on observations obtained at the European Southern Observatory (ESO, La Silla, Chile), Calar Alto Observatory operated by the MPIA (DSAZ, Spain), Lowell Observatory (Flagstaff/AZ, USA), and Hoher List Observatory (Germany).
## 1 Introduction
Although major galaxy mergers seem to be more spectacular and have therefore received the most attention, there is evidence that tidal interactions or the accretion of small, low-mass satellites (“minor mergers”) occur more frequently in the local universe (e.g. Frenk et al. frenk1988 (1988); Carlberg & Couchman carlberg1989 (1989)). Considering the high density of galaxies in groups and clusters and the fact that a large fraction of their members belongs to the dwarf galaxy population it is not unexpected that we know a large number of galaxies influenced by environmental effects of this order of magnitude. That includes both galaxies with clear signs of tidal interactions that happend in the recent past and those currently involved merging processes (Arp arp1966 (1966); Arp & Madore arp1987 (1987); Fried fried1988 (1988); Zaritsky et al. zaritsky1993 (1993), zaritsky1997 (1997)). A good example is our own Galaxy, that forms a future minor merger together with the Large Magellanic Cloud ($`M_{\mathrm{LMC}}2\times 10^{10}M_{}`$; $`M_{\mathrm{LMC}}/M_{\mathrm{MW}}0.1`$) and with several other satellites (e.g. Sgr dwarf elliptical) having smaller mass ratios (Irwin et al. irwin1985 (1985); Schommer et al. schommer1992 (1992); Toth & Ostriker toth1992 (1992); Ibata & Lewis ibata1998 (1998)). However, little is known about the rate of minor merging events, their influence on the structure and kinematics of galactic disks, and the efficiency of evoked perturbations.
In recent years this problem has been addressed by several numerical N-body simulations as well as analytical estimations. It was found that minor mergers and accretion events in the range $`M_{\mathrm{sat}}/M_{\mathrm{disk}}0.05\mathrm{}0.2`$ must be a more common processes in the local universe than previously argued, representing an important mechanism for driving the evolution of galaxies. In particular, it was concluded that one of the most striking structural changes produced by a single merger is a vertical disk thickening by a factor of 2–4, dependending noticeably on the initial disk properties such as the ratio between scale length and -height $`h/z_0`$ (Ostriker ostriker1990 (1990); Toth & Ostriker toth1992 (1992); Quinn et al. quinn1993 (1993); Mihos et al. mihos1995 (1995); Walker et al. walker1996 (1996)). However, the self-consistent simulations recently carried out by Velazquez & White (velazquez1999 (1999)) indicate that these results tend to overestimate the vertical disk heating. They find that the heating factor is closer to 1.5–2, depending mainly on the satellite orbit (i.e. prograde or retrograde) and the mass of the bulge component. Furthermore, the obtained results might still be influenced by counteracting processes such as dissipative gas cooling, subsequent star formation, the presence of several stellar disk components, etc. On the other hand, the observed thinness of typical late-type galaxy disks without indications of tidal interaction/accretion constrains the value of vertical disk heating. It is thus unlikely that such “superthin” galaxies have absorbed more than a few percent of their mass within their lifetime (Toth & Ostriker toth1992 (1992)).
At present there are only a few observational studies that aim at proving the effects predicted by the simulations. Zaritsky (zaritsky1995 (1995)) analyzed observations of nearby spiral galaxies based on magnitude residuals from the Tully-Fisher relationship, chemical abundance gradients, and asymmetries in their stellar disks. He concluded that even relatively isolated spiral galaxies have experienced accretion of companion galaxies over the last few Gigayears. In their series of studies Reshetnikov et al. (reshetnikov1993 (1993)) and Reshetnikov & Combes (reshetnikov1996 (1996), reshetnikov1997 (1997)) investigated the effects of tidally-triggered disk thickening between galaxies of comparable mass. They used optical photometric data of a sample of 24 interacting/merging and a control sample of 7 non-interacting disk galaxies. As a main result they find that the ratio $`h/z_0`$ of the radial exponential disk scale length $`h`$ to the constant scale height $`z_0`$ is about two times smaller for interacting galaxies. However, the relatively small galaxy samples used in these studies make it difficult to derive reliable estimates on the actual size of the structural changes. This also prevents a consistent check with the results from simulations.
Therefore we started a detailed statistical study in order to investigate systematically the influence of interactions and minor mergers on the radial and vertical disk structure of spiral galaxies in both optical and near infrared (NIR) passbands. Our study is based on a sample of 110 highly-inclined/edge-on disk galaxies, consisting of two subsamples of 61 non-interacting galaxies and 49 interacting/minor merger candidates. Additionally, 41 of these galaxies were observed in the NIR.
In Paper I (Schwarzkopf & Dettmar 2000a ) a detailed description of the project structure and its main questions was given. We reported on the sample selection, observations, and data reduction as well as on the disk modelling- and fitting procedure.
In Sect. 2 of this paper (Paper II) the sample and applied corrections are briefly summarized. In Sect. 3 we analyze the radial and vertical disk structure of both subsamples. The global disk parameters, their ratios, and the vertical brightness distribution are investigated. The derived colour gradients are also analyzed. We discuss the obtained results in Sect. 4 and summarize and conclude the paper in Sect. 5.
## 2 The data
### 2.1 Sample and observations
In Paper I we found that it is crucial for this study to have two subsamples of highly-inclined/edge-on galaxies ($`i85\mathrm{°}`$) that were selected carefully in order to diminish overlapping effects, i.e. a contamination introduced by an uncertain allocation of galaxies to the non-interacting or interacting/merging sample. For the latter sample we therefore used a classification scheme that was introduced by Arp & Madore (arp1987 (1987)). Additionally, for most of the minor merger candidates the mass ratio between the companion and the main body was checked. Finally, we have shown that the distribution of the morphological types between both subsamples is statistically indistinguishable over the whole range studied, i.e. between $`0T9`$ (Paper I). Although we cannot exclude overlapping effects completely, the remaining uncertainties in the classification of the subsamples were thus reduced to a minimum. A wrong allocation of objects would only lead to an underestimation of the actual differences between both samples.
Since a large sample of galaxies was needed for this study the observations were obtained with different telescopes and during several observing runs between February 1996 and June 1998. Details of the observations and the data reduction can also be found in Paper I.
### 2.2 Distances and corrections
The distances to the observed galaxies with known redshifts were calculated using a Hubble flow with a Hubble constant of $`H_0=75\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$, corrected for the “Virgocentric Flow” model predicted by Kraan-Korteweg (kraan-korteweg1986 (1986)). The model chosen is characterized by a local infall motion of the Local Group towards the Virgo cluster with an adopted velocity of $`v_{\mathrm{vc}}=220\mathrm{km}\mathrm{s}^1`$. It describes the motions of galaxies in the environment of the cluster by a non-linear flow-model. Assuming this model and the adopted Hubble constant, the distance of the Virgo cluster is $`r_{\mathrm{vir}}=15.8`$ Mpc.
The heliocentric velocities $`v_0`$ needed for this model are listed in Table 5, column (5). They were calculated from optical/HI-velocities taken from NED<sup>1</sup><sup>1</sup>1The NASA/IPAC Extragalactic Database (NED) is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration.. The heliocentric velocities $`v_0`$ were corrected for the system velocity $`v_{\mathrm{LG}}`$ of the Local Group via formula (2) in Richter et al. (richter1987 (1987)). The velocity corrections $`\mathrm{\Delta }v`$ required for this model were derived from Fig. 3a in Kraan-Korteweg (kraan-korteweg1986 (1986)).
The distance distribution of all sample galaxies with known redshifts ($`n=97`$) is shown in Fig. 1a. The nearest galaxy is NGC 4244 at 3.8 Mpc, the most distant galaxy is ESO 379-G20 at 193 Mpc (Table 5). Fig. 1b shows the distance distributions for both subsamples of 43 interacting/merging and 54 non-interacting galaxies. According to the test of Kolmogorov & Smirnov (Darling darling1957 (1957); Sachs sachs1992 (1992)) – hereafter KS – both distributions are statistically indistinguishable: the test result of 0.09 is significantly lower than the value necessary for the 20%-limit (0.22), which is the strongest of the KS-criteria. This is – in addition to the indistinguishable distribution of morphological types (Paper I, Fig. 1) – essential in order to avoid selection biases and thus to derive reliable disk parameters.
Since seeing effects become significant for small disk structures such as the scale height of flat disks (for seeing conditions around $`2\mathrm{}`$ and for features $`4\mathrm{}`$ the error amounts to 20%) all disk parameters listed in Table 5 were corrected in the following way: using the values of predominant seeing conditions (derived from averaged FWHM values of stars in the resulting frames, given in Paper I, Table 2) the images were de-convoluted with the Lucy-Richardson algorithm (standard MIDAS routine). In order to profit from the full vertical resolution – important to distinguish between different vertical disk models used for fitting; Paper I, Sect. 4) – no vertical binning was applied.
### 2.3 Disk models and fitting procedure
In order to analyze and to compare the structure of disk components of a large sample of highly-inclined/edge-on spiral galaxies we developed an improved disk modelling- and fitting procedure that is based on a 3-dimensional luminosity distribution proposed by van der Kruit & Searle (1981a ,b ; 1982a ), hereafter KS I–III. The results presented in the next Section were derived using these disk models. An detailed description of their properties and the determination of global disk parameters was given in Paper I.
## 3 Results
The following statistical analysis compares the radial and vertical disk parameters of both galaxy samples based on the optical ($`R`$-band) data set. The derived main disk parameters are: inclination angle $`i`$, cut-off radius $`R_{\mathrm{max}}`$, scale length $`h`$, scale height $`z_0`$, and best-fitting vertical model $`f(z)`$. They are listed in Table 5, columns (7)-(11).
### 3.1 The radial and vertical disk structure
#### 3.1.1 The cut-off radius “$`R_{\mathrm{max}}`$
The distribution of disk cut-off radii derived from the values in Table 5 is shown in Fig. 2 for the samples of non-interacting and interacting/merging galaxies. Both distributions cover a wide range between 4 kpc $`R_{\mathrm{max}}`$ 45 kpc, with the same global maximum around $`14`$ kpc. Slight differences can be detected in a region of large cut-off radii: many of the interacting/merging galaxy disks are concentrated in a strong peak between 8 kpc $`R_{\mathrm{max}}`$ 16 kpc, followed by a noticable drop-off towards larger radii. The distribution ends abruptly at $`R_{\mathrm{max}}32`$ kpc.
Disks of non-interacting galaxies show a more regular distribution, decreasing from the maximum at $`R_{\mathrm{max}}=14`$ kpc towards larger radii. However, the median derived for both distributions is almost identical at $`(R_{\mathrm{max}})_{\mathrm{norm}.}`$ 17.2 kpc and $`(R_{\mathrm{max}})_{\mathrm{merg}.}17.0`$ kpc, respectively.
Since the slight differences detected between the two distributions are caused by only a few galaxies, the KS-test shows that both samples are close to unity and thus statistically indistinguishable (the result of 0.14 is clearly below the critical 20%-limit of 0.22).
#### 3.1.2 The disk scale length “$`h`$
In Fig. 3 the distribution of disk scale lengths is shown for the non-interacting and interacting/merging galaxy samples. As in the cut-off statistics, both distributions have approximately the same global maxima, located at $`h4`$ kpc. The disks of non-interacting galaxies possess scale lengths in a wide range between 1.5 kpc $`h`$ 16 kpc, with a regular decrease towards larger values. The distribution of disk scale lengths of interacting spirals shows a similar behaviour, but with a more sharply truncated end at $`h`$ 10.5 kpc. This results in a slightly different median for both distributions at $`h_{\mathrm{norm}.}5.6`$ kpc and $`h_{\mathrm{merg}.}4.9`$ kpc, respectively.
In analogy to the cut-off-statistics, the differences between both $`h`$-distributions are only due to a few galaxies (most probably reflecting the slightly different number of galaxies in both samples) and thus marginal. The KS-test therefore classifies these differences as statistically insignificant with a result of 0.11, which is clearly below the mentioned 20%-limit of 0.22.
#### 3.1.3 The ratio of radial parameters “$`R_{\mathrm{max}}/h`$
Although several studies of edge-on spiral galaxies were focussed on an investigation of disk scale parameters of non-interacting galaxies, only a few have studied the ratio of disk cut-off radius to the scale length $`R_{\mathrm{max}}/h`$ (KS I-III; Barnaby & Thronson barnaby1992 (1992); Barteldrees & Dettmar barteldrees1994 (1994); Schwarzkopf schwarzkopf1999 (1999); Schwarzkopf & Dettmar schwarzkopf1997 (1997)). One of the reasons might be the fact that for most purposes a disk model consisting of a radial exponentially decreasing luminosity without a cut-off may be a good approximation (Shaw & Gilmore shaw1989 (1989); de Grijs & van der Kruit grijs1996 (1996); de Grijs et al. grijsetal1997 (1997); Reshetnikow & Combes reshetnikov1996 (1996), reshetnikov1997 (1997)). However, the importance of a cut-off radius as a reasonable step towards a more realistic description of the properties of galactic disks and its necessity for a precise quantitative description of the observed radial disk profiles were impressively confirmed by the results of the former studies. The existence of a disk cut-off, although still objectionable within the framework of galaxy evolution, seems therefore now well accepted.
This is consistent with the results of our disk modelling procedure (Paper I), showing that the shape of radial disk profiles and thus the derived value for the disk scale length is also influenced by the size of the cut-off radius. Using a disk model without a cut-off would increase the (existing) uncertainties in estimating reliable disk scale lengths.
The ratios $`R_{\mathrm{max}}/h`$ found in this study for both samples of non-interacting and interacting/merging galaxies are summarized in Table 1. They are compared with the data from literature of published samples of non-interacting galaxies. Fig. 4 shows the $`R_{\mathrm{max}}/h`$-distribution for both the total galaxy sample and the two subsamples. Due to its statistically significance (108 galaxies) the distribution of the total galaxy sample is very regular and, though the slightly steeper drop-off towards smaller values, close to a Gaussian. The median of the total sample is at $`(R_{\mathrm{max}}/h)_{\mathrm{tot}}=3.59`$, that of the non-interacting and interacting/merging galaxy sample at $`(R_{\mathrm{max}}/h)_{\mathrm{norm}.}=3.66`$ and $`(R_{\mathrm{max}}/h)_{\mathrm{merg}.}=3.53`$, respectively (typical errors are given in Table 1). According to the KS-test both distributions are statistically indistinguishable with a result of 0.12 (critical 20%-limit is at 0.21).
The relatively high scatter of the merger-sample, which becomes apparent in a wider basis of the distribution (Fig. 4b), is mainly due to a radially perturbed disk structure of these galaxies. These perturbations may cause asymmetries in the radial disk profiles and thus substantially errors in the scale length and/or the cut-off radius. Unlike this result it is a remarkable fact that the $`R_{\mathrm{max}}/h`$-distribution of interacting/merging galaxies shows a very sharp peak at $`R_{\mathrm{max}}/h3.4`$, while non-interacting galaxies are spread in a much wider plateau.
Within the estimated errors the obtained $`R_{\mathrm{max}}/h`$-ratios of all (sub-)samples are consistent with other studies (Table 1). However, this study – dealing with the by far largest galaxy sample for a cut-off statistics and hence with smaller errors – indicates that the $`R_{\mathrm{max}}/h`$-ratios are lower than previously argued. The mean ratio now seems to be closer to $`(R_{\mathrm{max}}/h)_{\mathrm{tot}}=3.6`$ than to the value 4.2 often used in the literature. This fact may be explained by selection effects of studies dealing with statistically very small samples of relatively nearby ($`D20`$ Mpc) and medium-sized ($`R_{\mathrm{max}}25`$ kpc) galaxies, which is the case for some of these data. Therefore, some additional information on these studies are included in Table 1. A possible correlation between disk parameters and distance will be studied in Sect. 3.2.
#### 3.1.4 The disk scale height “$`z_0`$
The distribution of exponential disk scale heights $`z_0`$ – calculated by using the absolute $`z_0`$-values in Table 5, column (10), and the corresponding transformations in (5) of Paper I – is shown in Fig. 5a for the non-interacting galaxy sample, and in Fig. 5b for the interacting/merging galaxies. By comparing these diagrams clear differences concerning both the trend and the median of the distributions can be detected:
The majority ($`60\%`$) of normal galaxies is concentrated in a region $`z_01.1`$ kpc, with a clear maximum between 400 pc $`z_0`$ 800 pc. The distribution shows a sharp truncation towards extremely thin disks at $`z_0=300`$ pc, followed by a very regularly decreasing part towards thicker disks up to $`z_03`$ kpc.
Unlike this, the distribution of interacting/merging galaxies shows a trend that is completely contrary to that found for normal galaxies: the scale height increases rapidly towards thicker disks, having a maximum between 1.2 kpc $`z_0`$ 1.6 kpc, and a higher frequency of disks thicker than 2 kpc. The median for both distributions is $`(z_0)_{\mathrm{norm}.}1.0`$ kpc and $`(z_0)_{\mathrm{merg}.}1.5`$ kpc. It is also remarkable that very thin disks in the range 250 pc $`z_0`$ 450 pc, as they are frequently observed in the sample of non-interacting galaxies (e.g. UGC 231, UGC 4278, UGC 4943, see Table 5), are completely missing in the interacting sample. Since this is not due to the sample selection criteria (Paper I) it is most likely a result of the various collective instabilities of such thin disks, in particular against external perturbations as they are commonly evoked by tidal interactions or even minor mergers (Gerin et al. gerin1990 (1990); Toth & Ostriker toth1992 (1992)).
The KS-test confirms that the differences between both distributions are statistically significant: the result of 0.24 is above both the 20% (0.22) and the 15% (0.23) limit of the test. Since it was shown (Fig. 1 of Paper I and Sect. 2 of this paper, resp.) that the absolute disk parameters of both galaxy samples are not affected by selection biases – the distance- and type distribution were found almost indistinguishable – it can be concluded that galactic disks affected by interactions/minor mergers are $`1.5`$ times thicker on average. We therefore infer that the disk thickening is in fact caused by the interaction or the merging process.
However, it is not yet clear whether this thickening effect was evoked only by a locally increased scale height due to a vertically perturbed disk structure or by global disk thickening. This can be clarified only after a detailed analysis of the vertical disk structure. Therefore, the behaviour of the disk scale height with radial distance, i.e. $`z_0(R)`$, and the effects of vertical disk perturbations will be investigated in detail in a forthcoming paper (Paper III, Schwarzkopf & Dettmar 2000c ).
#### 3.1.5 The ratio of radial to vertical parameters “$`h/z_0`$
The ratio of radial to vertical disk scale parameters $`h/z_0`$ – hence a normalized thickness – is very suitable for characterizing and comparing the disk structure of edge-on galaxies independently of their distance. It is thus – unlike the absolute values of scale heights – more reliably for a detection of small changes in the vertical disk structure.
Fig. 6 shows the $`h/z_0`$-distribution for both samples of non-interacting and interacting/merging galaxies. As it can be seen clearly, the ratio $`h/z_0`$ for normal galaxies covers a wide range $`2h/z_018`$ between extreme thick and thin disks. Most of these galaxies are concentrated between $`2h/z_012`$, with a maximum at $`h/z_05.3`$. There is only a slight decrease in the number of flat disks from the maximum of the distribution towards extreme flat disk ratios, i.e. between $`6h/z_012`$.
Unlike this, the most striking feature in the $`h/z_0`$-distribution for interacting/merging galaxies is a very sharp concentration between $`2h/z_06`$, while the flat disks with ratios typically $`h/z_0>7`$ are completely missing. The distribution peaks at $`h/z_03.7`$, and there is no smooth transition towards thicker or thinner disks on both sides of this sharply truncated distribution.
The ratio of the median values for both distributions and hence a lower limit for vertical disk thickening is $`(h/z_0)_{\mathrm{norm}:\mathrm{merg}}=7.1:4.31.7`$. This factor, however, considerably underestimates the differences between both distributions due to the above mentioned lack of (extremely) thin disks ratios. The differences between both distributions are, according to the KS-test, statistically significant with a result of 0.41 even if the strongest test criterium, the 0.1%-level (limit 0.38), is used.
Together with the nearly unchanged scale lengths (Fig. 3) and the differences found between the absolute values of disk scale heights (Fig. 5) it can be concluded that vertical thickening of galactic disks affected by interactions/minor mergers amounts to $`70\%`$. The changes of the disk structure result mainly from an increase in scale height.
Finally, Fig. 7a shows that the ratio $`h/z_0`$ of non-interacting galaxies correlates with the morphological type of galaxies, in the sense that the disks become systematically thinner from early types (S0, $`h/z_02\pm 2`$) to late types (Sc/Sd, $`h/z_08\pm 2`$). Despite the relatively high intrinsic scatter, there is a smooth transition between these two extremes. It should be stressed that – due to the corrections necessary in order to compare disk scale heights derived from different vertical luminosity distributions (exp, sech, and $`\mathrm{sech}^2`$) – the ratio $`h/z_0`$ can be higher than 10 for some of the disks. This value represents the maximum theoretical value allowed for stable disks, derived from the so-called “maximum disk” fits (Bottema bottema1993 (1993)). The results obtained are therefore not unexpected, and comparable ratios were also found in earlier studies (de Grijs grijs1997 (1997); de Grijs & van der Kruit grijs1996 (1996); Schwarzkopf & Dettmar schwarzkopf1997 (1997); Shaw & Gilmore shaw1989 (1989)).
In contrast to this there is no such correlation between galaxy type and ratio $`h/z_0`$ for the sample of interacting/merging galaxies (Fig. 7b). The latter typically possess thickened and disturbed disks, and are therefore concentrated in the lower right part of the panel.
For comparison, the position of our Galaxy disk is also indicated in Fig. 7. The disk of the Milky Way consists presumably of both a thin and a thick component with $`h/z_012`$ and $`h/z_04`$, respectively (according to Gilmore & Reid gilmore1983 (1983)).
### 3.2 The dependence of disk parameters on distance
As briefly mentioned in Sect. 3.1.3 (Table 1), the samples of some previous studies seem to be biased towards nearby objects. In order to check if the distance as a free parameter has any obvious effect on the derived disk parameters of this study, we analyzed 3 sub-samples (no differences were made between interacting and non-interacting galaxies) defined by various distance ranges (Table 2).
To ensure statistically large enough subsamples the total sample (97 galaxies) was split into three subsamples containing about 32 galaxies each. According to the non-uniform distribution of distances (Fig. 1) this leads to different distance intervals. Afterwards, the medians for all disk parameters in each distance-subsample were estimated. For a better comparison, these values are also given in per cent of the median of the total sample. Additionally, the averaged morphological galaxy type is listed for each subsample.
As can be seen in Table 2, there is a clear trend for all absolute disk parameters ($`R_{\mathrm{max}},h,z_0`$) in the sense that their median values are increasing with distance. This behaviour reflects, however, undoubtly a selection effect and is therefore not unusual for all those studies that are using similar selection criteria: due to the limited spatially resolution of the images the selection of suitable, remote galaxies is biased towards large and bright objects, disfavouring the relative number of physically small galaxies. This is confirmed by the strikingly simultaneous trend for all absolute disk parameters listed, which goes, on average, from 70% and 90% to 145% of the median of the total sample (Table 2, values in parenthesis). In spite of this trend the distance-independent ratios $`R_{\mathrm{max}}/h`$ and $`h/z_0`$ both stay nearly constant over the whole range.
Hence, within the given errors there is no correlation between disk parameters of different distance intervals and corresponding galaxy type. Since the distance distribution of interacting and non-interacting galaxies is statistically indistinguishable (Fig. 1b), the observed selection effect applies to both samples and does therefore not influence the comparison of disk parameters derived in this study.
### 3.3 The vertical surface brightness distribution
The disk models applied in this study use a set of 3 different functions $`f(z)`$ in order to describe the luminosity distribution $`L(z)`$ vertically to the disk plane: $`f(z)`$ exp, sech, and $`\mathrm{sech}^2`$. These functions were proposed in some fundamental papers by van der Kruit & Searle (KS I-III), Wainscoat et al. (wainscoat1989 (1989), wainscoat1990 (1990)), and Burkert & Yoshii (burkert1996 (1996)). A detailed description of their properties as well as a comparison between different distributions were given in Sect. 4 of Paper I.
The quantitative results and experiences made in this study after fitting the disk profiles of about 150 highly-inclined/edge-on galaxies in optical and in near infrared (NIR) passbands can be summarized as follows (the complete statistics is listed in Table 3):
* A combination of 3 different luminosity distributions $`f(z)`$ allows a very flexible description of vertical disk profiles of all galaxies observed.
* The fit quality achieved is better than $`\pm \mathrm{\hspace{0.17em}0}.^\mathrm{m}2`$ for nearly all of the non-interacting galaxies investigated both in the optical and in NIR, even at small $`z`$.
* The fit quality of galaxy disks affected by interaction/minor merger is, in principle, comparable to that found for non-interacting galaxies. However, some galaxies in the first sample possess vertical profiles with larger deviations from an ideal disk.
* These deviations are mainly due to tidal perturbations on short scales and/or a warped disk. Such features seem to be characteristic for disks in an intermediate stage of interactions/minor mergers. A detailed study of the vertical disk structure will be given in Paper III.
* The statistics for the best-fitting vertical luminosity distribution $`f(z)`$ – applied to both optical and NIR disk profiles – is as follows (optical : NIR, Table 3, columns 2-4):
(56 : 66)% $``$ sech; (36 : 27)% $``$ exp; (7 : 7)% $`\mathrm{sech}^2`$.
* Thus, the vertical luminosity profiles of nearly all ($`93\%`$) of the galaxies are non-isothermal. In fact the profiles are more sharply peaked and preferentially somewhat closer to a sech- than to an exp-distribution.
* Statistically, the fraction of galaxies with a certain vertical luminosity profile (exp, sech, $`\mathrm{sech}^2`$) is independent of the passband investigated (differences $`<10\%`$).
* In the optical there is no fundamental difference between vertical disk profiles of non-interacting and interacting/merging galaxies.
* Accordingly, almost the same percentage of (non-interacting : interacting) galaxies shows identical vertical distributions in the optical (Table 3):
(56 : 57)% $``$ sech; (36 : 37)% $``$ exp; (8 : 6)% $`\mathrm{sech}^2`$.
* In the NIR, interacting galaxies display preferentially sech-profiles ($`76\%`$ sech; $`19\%`$ exp), while in the optical this distribution is shifted towards the exp-profile ($`57\%`$ sech; $`37\%`$ exp).
* The results obtained are independent of the morphological type of galaxies.
### 3.4 Disk colour gradients
In order to analyze colour gradients derived from measurements of radial and vertical disk parameters in optical and in near infrared passbands ($`R/K`$), the mean ratios of disk cut-off radius $`R_{\mathrm{max}}`$, scale length $`h`$, and scale height $`z_0`$ – obtained for both subsamples – are listed in Table 4.
The radial and vertical disk parameters found for $`K`$ and $`R`$ passbands in the total galaxy sample indicate that the $`R`$-band values are systematically larger, i.e. of the order of $`(R/K)=1.30\pm 0.4`$. In spite of the large intrinsic scatter, comparison with literature data shows that the values (and also the errors) obtained here are consistent with gradients derived by Giovanardi & Hunt (giovanardi1988 (1988)) and de Grijs (grijs1997 (1997)). They found $`(F/K)=1.20\pm 0.42`$ for the $`F`$\- and $`K`$ passbands, and $`(B/K)=1.56\pm 0.45`$; $`(I/K)=1.19\pm 0.17`$ for the $`B`$,- $`I`$\- and $`K`$ passbands, respectively.
While systematically higher optical values have been found for each of the subsamples of interacting/merging $`(R/K)=1.23\pm 0.5`$ and non-interacting galaxies $`(R/K)=1.38\pm 0.5`$, the gradients of the latter sample are, however, systematically higher (Table 4).
Although the results obtained for the two subsamples are difficult to interpret, the systematic differences found for all colour gradients as well as the good agreement with other studies indicate that these gradients are not due to the large intrinsic errors. It should, however, be stressed that the low S/N ratio in the outskirts of disks of a number of faint galaxies obtained in the near infrared largely prevents a precise determination of the cut-off radius $`R_{\mathrm{max}}`$. Longer intergration times than the typical 30–40 min (on source) would be therefore necessary in order to derive more reliable values.
## 4 Discussion
Considering the small mass ratio between merging satellites and disks investigated here – $`M_{\mathrm{sat}}/M_{\mathrm{disk}}0.1`$ – the factor found for vertical disk thickening ($`1.6`$ on average) and thus the efficiency of vertical heating is substantial. This value is, however, significantly lower than the factor of 2–4 obtained in previous studies (Reshetnikov & Combes reshetnikov1996 (1996), reshetnikov1997 (1997); Toth & Ostriker toth1992 (1992)). The differences can be explained by the different mass ratio between strongly interacting systems (galaxies of comparable mass) investigated by Reshetnikov & Combes, their simplyfied disk model (isothermal) applied to all disks and the neglect of precise disk inclination. In contrast to the set of fully self-consistent N-body simulations made recently by Velazquez & White (velazquez1999 (1999)) the analysis of Toth & Ostriker (toth1992 (1992)) ignores the coherent response of the disk and its interaction with the halo. Additionally, their assumption that the orbital energy of the satellite is deposited locally in the disk is clearly unrealistic.
In fact, the increase of disk scale height by a factor of $`1.6`$ found in this study corresponds quite well with the value of 1.5–2 obtained by Velazquez & White (velazquez1999 (1999)). However, as already mentioned by Velazquez & White (velazquez1999 (1999)) and in the introduction of this study, vertical disk thickening due to a minor merger crucially depends on many other factors such as the mass and density profile of the sinking satellite, its orbit (prograde or retrograde), the content of gas deposited in the disk, and – presumably most important – on the morphological type of the galaxy. That, in turn, implies that tidally-triggered disk thickening strongly depends on the B/D ratio and hence on the initial disk thickness of the galaxy, which is characterized by the ratio $`h/z_0`$. This is confirmed by the $`h/z_0`$-statistics obtained in this study (Fig. 6), showing a total lack of thin interacting/merging galaxies with $`h/z_0>7`$. This result has some direct, important consequences on the evolution of disk galaxies on cosmological timescales (Toomre toomre1977 (1977); Weil et al. weil1998 (1998)) and would also constrain the different scenarious discussed for disk heating (Jenkins & Binney jenkins1990 (1990), jenkins1991 (1991); Sanchez-Salcedo sanchez1999 (1999); Valluri valluri1993 (1993)). Therefore a more detailed study of the parameter space using our supplementary N-body simulations, combined with the obtained results, will be given in a forthcoming paper. The mentioned good quantitative agreement between simulation and observation, however, indicates that – in spite of the rather simple approach and the number of open questions on the details of minor merger processes – the changes of the vertical disk structure must be of this order (Kleinschmidt et al. kleinschmidt1999 (1999); Schwarzkopf schwarzkopf1999 (1999); Schwarzkopf & Dettmar schwarzkopf1998 (1998), schwarzkopf1999\_1 (1999)), i.e. somewhat lower than previously argued.
Furthermore, it is still an open question whether the radial disk structure of galaxies suffering interactions/minor mergers within the mass range studied here changes significantly. The differences between both the disk cut-off and scale length statistics obtained (10% and 20% on average, resp.) are just on the level of statistical significance and do therefore not allow for an interpretation. On the other side, there is observational evidence that radial disk shrinking is a typical aftermath of tidal interactions between galaxies with comparable masses (Reshetnikov & Combes reshetnikov1996 (1996), reshetnikov1997 (1997)). If so, this should also apply to smaller interactions and, in particular, to minor mergers (Schwarzkopf schwarzkopf1999 (1999)). For further clarity on this point it is therefore necessary to analyze the radial behaviour of such galactic disks in greater detail and on the basis of an expanded galaxy sample, preferably supported by N-body simulations.
The fact that nearly all galactic disks investigated (93%) possess vertical luminosity profiles which are more sharply peaked than an isothermal distribution reinforces the results of previous observational studies (e.g. de Grijs et al. grijsetal1997 (1997); Schwarzkopf & Dettmar schwarzkopf1997 (1997)) that have ruled out the validity of the $`\mathrm{sech}^2`$-distribution as an adequate quantitative description for most galactic disks, especially close to their plane. This fact, together with the result that almost the same percentage of interacting and non-interacting galaxies shows an identical vertical disk structure (with differences smaller than 2%, see Table 3), indicates that regional damaging effects, asymmetries or perturbations evoked by tidal interactions are non-persistent phenomena with lifetimes significantly shorter than disk thickening. Furthermore, it implies that interactions/minor mergers within the investigated mass range are not capable to destroy the initial vertical disk structure.
## 5 Summary and conclusions
In this work a detailed statistical study is presented in order to investigate the effects of minor mergers and tidal interactions in the range $`M_{\mathrm{sat}}/M_{\mathrm{disk}}0.1`$ on the radial and vertical structure of galactic disks. The fundamental disk parameters of 110 highly-inclined/edge-on disk galaxies are determined in optical and in near infrared passbands. This sample consists of two subsamples of 61 non-interacting and 49 strongly interacting/merging galaxies, respectively. Additionally, 41 of these galaxies were observed in the near infrared. The main conclusions can be summarized as follows:
1. The structural changes of galactic disks affected by interaction/low-mass satellite infall are most noticeable in the direction perpendicular to the disk plane.
2. While the majority of non-interacting galaxies possess a typical exponential disk scale height of $`z_0700`$ pc, disks of minor mergers were found to be systematically thicker with $`z_01.3`$ kpc.
3. On average, galactic disks affected by interactions or minor mergers have $`1.5`$ times larger scale heights and thus vertical velocity dispersions than unperturbed disks.
4. The ratio of radial to vertical scale parameters, i.e. the normalized disk thickness $`h/z_0`$, is $`1.7`$ times smaller for the interacting/merging sample.
5. Ratios $`h/z_0>7`$, as they are typically for flat galaxies, are completely missing in the interacting/merging sample. This implies that vertical disk heating is most efficient for such (extremely) thin disks.
6. The radial disk structure of interacting/merging galaxies, characterized by the cut-off radius $`R_{\mathrm{max}}`$ and scale length $`h`$, shows no statistically significant changes.
7. The vertical luminosity profiles of all galactic disks investigated show the following distribution (independent of the sample and passband):
60% $``$ sech; 33% $``$ exp; 7% $`\mathrm{sech}^2`$.
8. Thus, the majority (93%) of galactic disks possess non-isothermal vertical luminosity profiles and is somewhat closer to a sech- than to an exp distribution.
9. There are no fundamental differences between the vertical luminosity distribution of non-interacting and interacting/merging galaxies. Hence, the intrinsic distribution of disk stars keeps largely retained during and after interactions/minor mergers.
###### Acknowledgements.
We thank the referee of this series of papers, Dr. H. Wozniak, for his useful comments and suggestions. This work was supported by Deutsche Forschungsgemeinschaft, DFG, under grant no. GRK118/2. This research has made use of the NASA/IPAC Extragalactic Database (NED).
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# On the stability of periodic 2𝐷 Euler-𝛼 flows
## 1. Introduction
In Lagrangian mechanics a motion of a natural mechanical system is a geodesic line on a manifold - configuration space in the metric given by the difference of kinetic and potential energy. The configuration space for the fluid motion in a domain $`M`$ is the group $`𝒟_\mu (M)`$ of volume-preserving diffeomorphisms of $`M`$. The corresponding (Lie) algebra is the algebra of divergence-free vector fields on $`M`$ vanishing on the boundary. The standard (Euler) model of an ideal fluid corresponds to the kinetic energy being given by the $`L^2`$ norm of the fluid velocity on $`M`$. That is, the right-invariant metric on $`𝒟_\mu (M)`$ is defined in the following way: its value at the identity of the group on a divergence-free vector field $`v`$ from the algebra is given by $`v,v=v_{l^2}=_M(v,v)𝑑x`$.
Recently, a number of papers (see, e.g., \[HMR, S 98, S 99\]) introduced the so called averaged Euler equations for ideal incompressible flow on a manifold $`M`$. The averaged Euler equations involve a parameter $`\alpha `$; one interpretation is that they are obtained by temporally averaging the Euler equations in Lagrangian representation over rapid fluctuations whose amplitudes are of order $`\alpha `$. The particle flows associated with these equations can be shown to be geodesics on a suitable group of volume-preserving diffeomorphisms but with respect to a right invariant $`H^1`$ metric instead of the $`L^2`$ metric.
The case of area-preserving diffeomorphisms of the two-dimensional torus with a right invariant $`L^2`$ metric was analyzed by Arnold who showed (see, e.g. \[A 66, AK 98\]) that “in many directions the sectional curvature is negative”. In this paper we consider geodesic stability problem for the group $`𝒟_\mu (T^2)`$ with a right invariant $`H^1`$ metric which is related to the average Euler flows.
The instability discussed in this paper is the exponential *Lagrangian* instability of the motion of the fluid, not of its velocity field. A stationary flow can be a Lyapunov stable solution of Euler equations, while the corresponding motion of the fluid is exponentially unstable. The reason is that a small perturbation of the fluid velocity field can induce exponential divergence of fluid particles.
## 2. Instability of the Euler flow on $`T^2`$
Here we review Arnold’s results for the group $`𝒟_\mu (T^2)`$ with a right invariant $`L^2`$ metric closely following \[AK 98\]. Recall some standard notations. Let $`B`$ denote the bilinear form on a Lie algebra $`𝔤`$ defined by the relation $`B(\xi ,\eta ),\zeta =\xi ,[\eta ,\zeta ]`$, where $`\xi ,\eta ,\zeta 𝔤`$ $`[,]`$ is the commutator in $`𝔤`$ and $`,`$ is the inner product in the space $`𝔤`$.
The (Riemannian) *curvature tensor* $`R`$ describes the infinitesimal transformation on a tangent space obtained by parallel translation around an infinitely small parallelogram. For $`u,v,wT_{x_0}M`$, the action of $`R(u,v)`$ on $`w`$ can be expressed in terms of covariant differentiation as follows
$$R(u,v)w=(_{\overline{u}}_{\overline{v}}\overline{w}+_{\overline{u}}_{\overline{v}}\overline{w}++_{\{\overline{u},\overline{u}\}}\overline{w})|_{x=x_0},$$
(2.1)
where $`\overline{u},\overline{v},\overline{w}`$ are any fields whose values at the point $`x_0`$ are $`u,v,w`$.
The *sectional curvature* of $`M`$ in the direction of the two-plane spanned by any two vectors $`u,vT_{x_0}M`$ is the value
$$C_{uv}=\frac{R(u,v)u,v}{u,uv,vu,v^2}.$$
(2.2)
Theorem $`3.2`$ of \[AK 98\] gives explicit formulas for the inner product, commutator, operation $`B`$, connection, and curvature of the right invariant $`L^2`$ metric on the group $`𝒟_\mu (T^2)`$. These formulas allow one to calculate the sectional curvature in any two-dimensional direction.
The divergence-free vector fields that constitute the Lie algebra of the group $`𝒟_\mu (T^2)`$ can be described by their stream (Hamiltonian) functions with zero mean (i.e., $`v={\displaystyle \frac{H}{y}}{\displaystyle \frac{}{x}}+{\displaystyle \frac{H}{x}}{\displaystyle \frac{}{y}}`$). Thus, the Lie algebra can be identified with the space of real functions on the torus having zero average value \[AK 98\]. It is convenient to define such functions by their Fourier coefficients and to carry out all calculations over $``$.
Complexifying the Lie algebra one constructs a basis of this vector space using the functions $`e_k`$ (where $`k`$, called a *wave vector*, is a point of $`^2`$) whose value at a point $`x`$ of our complex plane is equal to $`e^{ı(k,x)}`$. This determines a function on the torus if the inner product $`(k,x)`$ is a multiple of $`2\pi `$ for all $`x\mathrm{\Gamma }`$. All such vectors $`k`$ belong to a lattice $`\mathrm{\Gamma }^{}`$ in $`^2`$, and the functions $`\{e_k|k\mathrm{\Gamma }^{},k0\}`$ form a basis of the complexified Lie algebra.
Consider the parallel sinusoidal steady flow given by the stream function $`\xi =\mathrm{cos}(k,x)`$ and let $`\eta `$ be any other vector of the algebra, i.e. $`\eta =x_le_l`$, where $`x_l=\overline{x}_l`$. Theorem $`3.4`$ of \[AK 98\] states that the curvature of the group $`𝒟_\mu (T^2)`$ in any two-dimensional plane containing the direction $`\xi `$ is *non-positive* and is given by
$$C_{\xi \eta }=\frac{S}{4}\underset{l}{}a_{kl}^2|x_l+x_{l+2k}|^2,$$
(2.3)
where $`a_{kl}={\displaystyle \frac{(k\times l)^2}{|k+l|}}`$, $`k\times l=k_1l_2k_2l_1`$ is the (oriented) area of the parallelogram spanned by $`k`$ and $`l`$, and $`S`$ is the area of the torus. Then, a corollary of this theorem states that the curvature in the plane defined by the stream functions $`\xi =\mathrm{cos}(k,x)`$ and $`\eta =\mathrm{cos}(l,x)`$ is
$$C_{\xi \eta }=(k^2+l^2)\mathrm{sin}^2\beta \mathrm{sin}^2\gamma /4S,$$
(2.4)
where $`\beta `$ is the angle between $`k`$ and $`l`$, and $`\gamma `$ is the angle between $`k+l`$ and $`kl`$.
## 3. Stable directions for the Euler-$`\alpha `$ flow on $`T^2`$
In this section we present new results on the sectional curvature of the group of area-preserving diffeomorphisms of a two-torus with a right invariant $`H^1`$ metric in view of the application to the Lagrangian stability analysis following Arnold \[A 66\]. The foundations for these results were established in \[S 98\] where the continuous differentiability of the geodesic spray of $`H^1`$ metric on $`𝒟_\mu ^s(M)`$ for an arbitrary Riemannian manifold $`M`$ was proved.
We start with an analog of Theorem $`3.2`$ of \[AK 98\]. Define an operator $`A^\alpha :^2_+,kk^2(1+\alpha ^2k^2)`$. It corresponds to the $`H^1`$ norm in the Fourier space and is simply given by $`k^2`$ in the case $`\alpha =0`$ when the $`H^1`$ metric effectively becomes the $`L^2`$ metric.
###### Theorem 3.1.
The explicit formulas for the inner product, commutator, operation $`B`$, and connection of the right invariant $`H^1`$ metric on the group $`𝒟_\mu (T^2)`$ have the following form:
$$e_k,e_l=A^\alpha (k)\delta _{k,l}$$
(3.1)
$$[e_k,e_l]=(k\times l)e_{k+l}$$
(3.2)
$$B(e_k,e_l)=b_{k,l}e_{k+l},\mathrm{where}b_{k,l}=(k\times l)\frac{A^\alpha (k)}{A^\alpha (k+l)}$$
(3.3)
$$_{e_k}e_l=d_{k,k+l}e_{k+l},\mathrm{where}d_{k,k+l}=\frac{k\times l}{s}\left(1\frac{A^\alpha (k)A^\alpha (l)}{A^\alpha (k+l)}\right).$$
(3.4)
Using the definition of the curvature tensor (2.1) we obtain
$$\begin{array}{c}R_{k,l,m,n}R(e_k,e_l)e_m,e_n=(d_{l+m,k+l+m}d_{m,l+m}\hfill \\ \hfill +d_{k+m,k+l+m}d_{m,k+m}+(k\times l)d_{m,k+l+m})A^\alpha (k+l+m)S.\end{array}$$
(3.5)
We do not write here the explicit expression for $`R_{k,l,m,n}`$ as it is rather involved, but we note that it is non-zero only in the case $`k+l+m+n=0`$. We analyze a special case of the curvature in the plane defined by the stream functions $`\xi =\mathrm{cos}(k,x)`$ and $`\eta =\mathrm{cos}(l,x)`$ (notice that the corresponding flow is a solution of the averaged Euler equations). Then the sectional curvature is determined only by two terms (we ignore the scaling factor of the denominator in the definition (2.2)):
$$C_{\xi \eta }^{H^1}=\frac{1}{8}(R_{k,l,k,l}+R_{k,l,k,l})$$
The computation gives an explicit formula
$$\begin{array}{c}C_{\xi \eta }^{H^1}=\frac{S}{36}(k\times l)^2(4A^\alpha (k)+4A^\alpha (l)3A^\alpha (k+l)3A^\alpha (kl)\hfill \\ \hfill +\frac{(A^\alpha (k)A^\alpha (l))^2}{A^\alpha (kl)}+\frac{(A^\alpha (k)A^\alpha (l))^2}{A^\alpha (k+l)})\end{array}$$
(3.6)
which we rewrite in the following form
$$\begin{array}{c}C_{\xi \eta }^{H^1}=\rho ^2\{A^\alpha (k+l)A^\alpha (kl)(4A^\alpha (k)+4A^\alpha (l)3A^\alpha (k+l)3A^\alpha (kl))\hfill \\ \hfill +(A^\alpha (k)A^\alpha (l))^2(A^\alpha (k+l)+A^\alpha (kl))\},\end{array}$$
(3.7)
where $`\rho ^2={\displaystyle \frac{S(k\times l)^2}{36A^\alpha (k+l)A^\alpha (kl)}}`$ is a function of $`k,l,\alpha `$ and is strictly positive. Hence, the sign of the curvature is determined by the expression in the bracket, which is a cubic polynomial in $`\alpha ^2`$:
$$B(\alpha ,k,l)b_0+b_1\alpha ^2+b_2(\alpha ^2)^2+b_3(\alpha ^2)^3,$$
(3.8)
so that $`C_{\xi \eta }^{H^1}=\rho ^2B(\alpha ,k,l)`$.
Numerical analysis of this complicated expression shows that the sectional curvature becomes positive for some values of $`\alpha >\alpha _0`$ when $`kl`$ is small. Fig. (3.1) is representative of a typical behavior of the curvature as a function of $`\alpha `$ for $`l=k+ϵ`$, where $`ϵk`$ is small. Based on this numerical evidence we analyze further analytically the case $`l=k+ϵ`$, where $`ϵk`$ is small. Compute the coefficients $`b_n`$ in (3.8) as power series in $`ϵ`$
$`b_0=64k^4ϵ^2+16k^2(k,ϵ)^2+𝒪(ϵ^4)`$ (3.9)
$`b_1=224k^6ϵ^2+128k^4(k,ϵ)^2+𝒪(ϵ^4)`$ (3.10)
$`b_2=640k^8ϵ^2+320k^6(k,ϵ)^2+𝒪(ϵ^4)`$ (3.11)
$`b_3=256k^8(k,ϵ)^2+𝒪(ϵ^4)`$ (3.12)
Notice that the coefficient of the highest degree is positive while all the rest are negative. Hence, for $`k>1/\alpha `$ it defines the leading term which increases with $`\alpha `$, while the other coefficients are responsible for initial decrease seen in Fig. (3.1). We summarize our result in the following theorem.
###### Theorem 3.2.
Consider the sectional curvature of the group $`𝒟_\mu (T^2)`$ equipped with the right invariant $`H^1`$ metric in the plane defined by the stream functions $`\xi =\mathrm{cos}(k,x)`$ and $`\eta =\mathrm{cos}(l,x)`$, where $`l=k+ϵ`$. Then, for $`|ϵ|`$ sufficiently small, for any $`k`$ there is an $`0<\alpha _0(k)<1`$, such that for all $`\alpha >\alpha _0(k)`$ the corresponding sectional curvature is positive.
## Acknowledgments
The authors would like to thank Jerrold E. Marsden for helpful comments and the Center for Nonlinear Science of Los Alamos National laboratory for providing a valuable setting where much of this work was performed.
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# Tachyon condensation in the D0/D4 system
## 1 Introduction
The presence of a tachyon in the spectrum of a string theory does not necessarily imply that the theory is inconsistent. In recent times there is accumulating evidence that the tachyon found in certain open string theories actually condense to a new stationary point of the potential. For example, the tachyon of the bosonic open string theory represents the instability of the D25-brane to decay to the vacuum. The tachyon of a non-BPS brane of Type II string theories also represents the instability of the non-BPS brane to decay to the vacuum.
Recent works have shown the existence of a stationary point in the tachyon potential. This potential is computed computed in string field theory using level truncation as an approximation method . The value of the tachyon potential at this stationary point agrees with the tension of the unstable branes. This has been shown to a remarkable degree of accuracy both for the bosonic open string field theory and for the superstring field theory on the non-BPS brane in Type II string theories . The physics of the solitons in the tachyon potential also seem to be reproduced in the level truncation approximation scheme . Tachyon condensation has recently been studied in effective field theories on non-commutative spaces in and its relationship to string field theory has been discussed in .
Tachyons are present in other unstable systems. They can be interpreted as the instability of the system to decay to stable systems. For all the cases analyzed so far in string field theory, the system decays to the vacuum. It would be interesting to find systems with tachyons, which decay to BPS states. For example consider the tachyon present for the open strings stretching between Dp-brane and D(p+2)-branes. This tachyonic potential was studied in first quantized string theory in . It represented the instability of the system to for the system to form a bound state of Dp and D(p+2)-branes. Another reason to study other systems with tachyon is that one might be able to find tuning parameter in the tachyon potential. For the systems studied so far, the tachyon potential is universal . Thus there is no tuning parameter which can be used to compare against the level truncation approximation scheme. In this paper we study a system with a tachyon which decays to a BPS state other than the vacuum. There is also a tuning parameter in the tachyon potential.
We consider the system of a single D0 and D4-brane with a Neveu-Schwarz B-field in the spatial directions of the D4-brane. For generic values of the B-field there is a tachyon in the spectrum of the strings joining the D0-brane and the D4-brane . The tachyon in the spectrum of the $`(0,4)`$ strings represents the instability of this system to form a bound state of D0 and D4-brane. The tachyon is in the Neveu-Schwarz sector. We construct the tachyon vertex operator and use the Berkovits formalism of open superstring field theory to evaluate the tachyon potential at the zeroth level. The tachyon vertex contains operators involving matter primaries. Thus the potential falls outside the usual universality class of unstable branes in superstring theory. We show that there is a stationary point in the tachyon potential which approximates the mass defect in the formation of the D0/D4 bound state. The tachyon is almost on shell for small values of the B-field. Then the zeroth level approximation for the tachyon potential contribute to 70% of the mass defect. For large values of the B-field with with the Paffian, Pf$`(2\pi \alpha ^{}B)>0`$ we find that the tachyon potential reduces to that of a D0-brane anti-D0-brane pair. Thus reproducing 62% of the mass defect. For large of the B-field with Pf$`(2\pi \alpha ^{}B)<0`$ we find the zeroth level approximation contributes to 25% of the mass defect. The reason for this might be due to the presence of a large number of low lying states in the spectrum of the $`(0,4)`$ strings for this case .
The unstable D-brane decays to the closed string vacuum. Thus at the tachyon condensate the open string modes are not dynamical. We show that the corresponding statement for the D0/D4 bound state is that the $`(0,4)`$ strings decouple from the dynamics. To do this we look at the small fluctuations around the tachyon condensate. This is facilitated by the description of the D0/D4 bound state as an instanton in the noncommutative $`U(1)`$ gauge theory of the D4-brane. The tachyon instability is the instablity of zero size instanton in the noncommutative $`U(1)`$ gauge theory of the D4-brane. At the tachyon condensate the $`(0,4)`$ strings become massive and decouple from the dynamics. This is seen from the fact that the $`(0,4)`$ strings correspond to the scale of the instanton via the ADHM construction. The noncommutative $`U(1)`$ instanton does not have scale moduli. Thus at the level of the moduli space approximation we can see that the $`(0,4)`$ strings decouple from the dynamics at the tachyon condensate.
The organization of this paper is as follows. In Section 2 we review the D0/D4 system with the B-field in the spatial directions of the D4-brane. In Section 3 we set up the string field theory of the $`(0,4)`$ strings We write down the vertex operators corresponding to the tachyon and calculate the tachyon potential. In Section 4 we calculate the expected mass defect in the formation of the D0/D4 bound state and compare it with the value of the stationary point in the tachyon potential. In Section 5 we identify the tachyon condensate as the instanton of the noncommutative gauge theory on the D4-brane and show the decoupling of the $`(0,4)`$ strings. Section 6 contains our conclusions. Appendix A contains the correlation function of twist operators. Appendix B contains details on the calculation of the tachyon potential.
## 2 The D0/D4 system
In this section we review the D0/D4 system with a Neveu Schwarz B-field along the spatial direction of the D4-brane . We discuss the mode expansions of the strings which join the D0-branes and the D4-branes.
Consider a single D0-brane and a single D4-brane of type IIA string theory in ten dimensions arranged as follows. The D4-brane is extended along the directions $`x^6,x^7,x^8,x^9`$. The open string spectrum consists of excitations of strings joining the D0-branes and the D4-branes among themselves and the excitations of open strings joining the D0-brane and the D4-brane. We denote the open strings joining the D0-brane and the D4-brane as the $`(0,4)`$ strings. Let the string world sheet co-ordinates of the $`(0,4)`$ string be $`X^\mu (\sigma ^0,\sigma ^1)`$. $`\mu `$ runs from $`0,\mathrm{}9`$, and $`\sigma ^1`$ lies between $`0`$ and $`\pi `$. We work with Euclidean world sheet signature. Now turn on a constant B-field along the spatial directions of the D4-brane. We can choose a generic value of the B-field as given below
$$B_{ij}=\frac{1}{2\pi \alpha ^{}}\left(\begin{array}{cccc}0& b& 0& 0\\ b& 0& 0& 0\\ 0& 0& 0& b^{}\\ 0& 0& b^{}& 0\end{array}\right)$$
(1)
Where $`i,j`$ runs from $`6\mathrm{}9`$. We choose the metric $`g_{ij}=\delta _{ij}`$. The boundary conditions of the world sheet co-ordinates with these moduli turned on is given by
$`_{\sigma ^1}X^i+2\pi i\alpha ^{}B_{ij}_{\sigma ^0}X^j|_{\sigma ^1=\pi }=0`$ (2)
$`_{\sigma ^0}X^i|_{\sigma ^1=0}=0`$
$`_{\sigma ^0}X^a|_{\sigma ^1=0,\sigma ^1=\pi }=0\text{ where }a=1,2,3,4,5`$
$`_{\sigma ^1}X^0|_{\sigma ^1=0,\sigma ^1=\pi }=0`$
The non-trivial mode expansions arise for the world sheet co-ordinates $`X^i`$, It is convenient to define co-ordinates
$`X^+=X^6+iX^7X^{}=X^6iX^7`$ (3)
$`X^+=X^6+iX^7X^{}=X^6iX^7`$
The mode expansions of $`X^+`$ and $`X^{}`$ are given by
$`X^+`$ $`=`$ $`i\sqrt{\alpha ^{}}{\displaystyle \underset{n}{}}\left[{\displaystyle \frac{\alpha _{n\nu }^+}{n\nu }}e^{(n\nu )(\sigma ^0+i\sigma ^1)}{\displaystyle \frac{\alpha _{n\nu }^+}{n\nu }}e^{(n\nu )(\sigma ^0i\sigma ^1)}\right]`$ (4)
$`X^{}`$ $`=`$ $`i\sqrt{\alpha ^{}}{\displaystyle \underset{n}{}}\left[{\displaystyle \frac{\alpha _{n+\nu }^{}}{n+\nu }}e^{(n+\nu )(\sigma ^0+i\sigma ^1)}{\displaystyle \frac{\alpha _{n+\nu }^{}}{n+\nu }}e^{(n+\nu )(\sigma ^0i\sigma ^1)}\right]`$
Where
$$e^{2\pi i\nu }=\frac{1+ib}{1ib},\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}\nu <1$$
(5)
The only non zero commutation relations are
$$[\alpha _{n\nu }^+,\alpha _{m+\nu }^{}]=(n\nu )\delta (n+m)$$
(6)
Similar mode expansions arise for the world sheet co-ordinates $`X^+`$ and $`X^{}`$. They are given by
$`X^+`$ $`=`$ $`i\sqrt{\alpha ^{}}{\displaystyle \underset{n}{}}\left[{\displaystyle \frac{\alpha _{n\nu ^{}}^+}{n\nu ^{}}}e^{(n\nu ^{})(\sigma ^0+i\sigma ^1)}{\displaystyle \frac{\alpha _{n\nu ^{}}^+}{n\nu ^{}}}e^{(n\nu ^{})(\sigma ^0i\sigma ^1)}\right]`$ (7)
$`X^{}`$ $`=`$ $`i\sqrt{\alpha ^{}}{\displaystyle \underset{n}{}}\left[{\displaystyle \frac{\alpha _{n+\nu ^{}}^{}}{n+\nu ^{}}}e^{(n+\nu ^{})(\sigma ^0+i\sigma ^1)}{\displaystyle \frac{\alpha _{n+\nu ^{}}^{}}{n+\nu ^{}}}e^{(n+\nu ^{})(\sigma ^0i\sigma ^1)}\right]`$
Where
$$e^{2\pi i\nu ^{}}=\frac{1+ib^{}}{1ib^{}}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}\nu ^{}<1$$
(8)
Similarly the only non zero commutation relations are
$$[\alpha _{n\nu ^{}}^+,\alpha _{m+\nu ^{}}^{}]=(n\nu ^{})\delta (n+m)$$
(9)
The mode expansions of world sheet superpartners of the bosonic fields are fixed by supersymmetry. The mode expansions of $`\psi ^+`$ and $`\overline{\psi }^+`$ of $`X^+`$ is given by
$`\psi ^+`$ $`=`$ $`i\sqrt{\alpha ^{}}{\displaystyle \underset{n}{}}\psi _{n+1/2\nu }^+e^{(n+1/2\nu )(\sigma ^0+i\sigma ^1)}`$ (10)
$`\overline{\psi }^+`$ $`=`$ $`i\sqrt{\alpha ^{}}{\displaystyle \underset{n}{}}\psi _{n+1/2\nu }^+e^{(n+1/2\nu )(\sigma ^0i\sigma ^1)}`$
The mode expansions of the the superpartners of $`X^{}`$ are given by
$`\psi ^{}`$ $`=`$ $`i\sqrt{\alpha ^{}}{\displaystyle \underset{n}{}}\psi _{n+1/2+\nu }^+e^{(n+1/2+\nu )(\sigma ^0+i\sigma ^1)}`$ (11)
$`\overline{\psi }^+`$ $`=`$ $`i\sqrt{\alpha ^{}}{\displaystyle \underset{n}{}}\psi _{n+1/2+\nu }^+e^{(n+1/2+\nu )(\sigma ^0i\sigma ^1)}`$
We have written the mode expansions in the Neveu-Schwarz sector. The only non-zero anti-commutation rules are given by
$$\{\psi _{n+1/2\nu }^+,\psi _{m+1/2+\nu }^{}\}=\delta (m+n)$$
(12)
The mode expansions and the anti-commutation relations for the superpartners of $`X^{}`$ and $`X^+`$ are given by replacing the variables in (10), (11) and (12) by their primed variables.
To show that the $`(0,4)`$ strings have a tachyon in their spectrum we evaluate the zero point energy. The zero point energy for four of the lowest energy states in the Neveu-Schwarz sector are given by
$`E_{}^+={\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}(\nu +\nu ^{})`$ $`E_+^+={\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{2}}(\nu +\nu ^{})`$ (13)
$`E_{}^{}={\displaystyle \frac{\nu \nu ^{}}{2}}`$ $`E_+^{}={\displaystyle \frac{\nu ^{}\nu }{2}}`$
Out of these four states two of them are projected out by the GSO projection. For the case of the D0-brane and the D4-brane the states with zero point energies $`E_\pm ^+`$ are retained . Thus we see unless $`\nu +\nu ^{}=1`$ there is a tachyon in the spectrum of the $`(0,4)`$ strings. For $`\nu +\nu ^{}<1`$ the tachyon (mass)<sup>2</sup> is given by $`E_{}^+/(2\alpha ^{})`$ and for $`\nu +\nu ^{}>1`$ the tachyon (mass)<sup>2</sup> is given by $`E_+^{}/(2\alpha ^{})`$.
## 3 String field theory of the $`(0,4)`$ strings
In this section we will compute the potential for the tachyon in the spectrum of the $`(0,4)`$ strings using string field theory. We will compute the tachyon potential upto the zeroth level in the level truncation approximation. We have seen in section 2 that the tachyon is in the Neveu-Schwarz sector. Therefore we can use the Wess-Zumino-Witten like open superstring field theory formulated by Berkovits to calculate the tachyon potential.
### 3.1 Open Superstring field theory
We will briefly review open superstring field theory. From this section onwards we work in the units $`\alpha ^{}=1`$ The string field theory action is given by
$$S=\frac{1}{g^2}(e^\mathrm{\Phi }Q_Be^\mathrm{\Phi })(e^\mathrm{\Phi }\eta _0e^\mathrm{\Phi })_0^1𝑑t(e^{t\mathrm{\Phi }}_te^{t\mathrm{\Phi }})\{(e^{t\mathrm{\Phi }}Q_Be^{t\mathrm{\Phi }}),(e^{t\mathrm{\Phi }}\eta _0e^{t\mathrm{\Phi }})\},$$
(14)
where $`\{A,B\}AB+BA`$. The open string coupling constant of the $`(0,4)`$ strings, $`g`$ is related to the closed string coupling constant coupling constant $`g_c`$ by $`g^2=g_c`$. The presence of the Neveu-Schwarz B-field does not change the relationship of the open string coupling constant to the closed string coupling constant. This is because the B-field is perpendicular to the common Neumann directions of the $`(0,4)`$ strings. The string field $`\widehat{\mathrm{\Phi }}`$ for the $`(0,4)`$ strings is given by
$$\widehat{\mathrm{\Phi }}=\mathrm{\Phi }_+^{(1)}II+\mathrm{\Phi }_+^{(2)}\sigma _3I+\mathrm{\Phi }_{}^{(3)}\sigma _1\sigma _1+\mathrm{\Phi }_{}^{(4)}\sigma _2\sigma _1$$
(15)
Where $`\mathrm{\Phi }_+^{(1)},\mathrm{\Phi }_+^{(2)}`$ stand for GSO even open string vertex operators and $`\mathrm{\Phi }_{}^{(3)},\mathrm{\Phi }_{}^{(4)}`$ stand for GSO odd open string vertex operators. These operators have ghost number $`0`$ and picture number $`0`$ in the combined conformal field theory of a $`c=15`$ superconformal matter system and the $`b,c,\beta ,\gamma `$ ghost system with $`c=15`$. The bosonized ghost fields $`\xi ,\eta ,\varphi `$ are related to $`\beta ,\gamma `$ by
$$\beta =\xi e^\varphi ,\gamma =\eta e^\varphi $$
(16)
$`\sigma _1,\sigma _2,\sigma _3`$ are $`2\times 2`$ pauli matrices and I is the $`2\times 2`$ identity matrix. The first set of matrices are the external Chan-Paton factors and the second set of matrices are the internal Chan-Paton factors. $`\widehat{Q}_B=Q_BI\sigma _3`$ where $`Q_B`$ is given by
$`Q_B`$ $`=`$ $`{\displaystyle }dzj_B(z)={\displaystyle }dz\{c(T_m+T_{\xi \eta }+T_\varphi +ccb+\eta e^\varphi G_m\eta \eta e^{2\varphi }b\}`$ (17)
$`T_{\xi \eta }`$ $`=`$ $`\xi \eta `$
$`T_\varphi `$ $`=`$ $`{\displaystyle \frac{1}{2}}\varphi \varphi ^2\varphi `$
$`T_m`$ is the mater stress tensor and $`G_m`$ is the matter superconformal generator. $`\widehat{\eta }_0=\eta _0I\sigma _3`$ and $`\eta _0`$ is the zeroth mode of the field $`\eta `$. The products of operators in $`\mathrm{}`$ is defined by
$$A_1\mathrm{}A_n=f_1^{(n)}A_1(0)\mathrm{}f_n^{(n)}A_n(0)$$
(18)
Here $`\mathrm{}`$ denotes the correlation function evaluated on the disc including the trace over the internal and the ordinary (external) Chan-Paton factors. $`fA`$ denotes the conformal transform of $`A`$ by $`f`$. For the case of the disc the functions $`f_k^{(n)}`$ are given by
$$f_k^{(n)}(z)=e^{\frac{2\pi i(k1)}{n}}\left(\frac{1+iz}{1iz}\right)^{2/n}\text{ for }n1$$
(19)
We now expand (14) to compute the tachyon potential to the zeroth level. For this it is sufficient to retain only terms upto the fourth order in the string field in the expansion of (14). This is given by
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2g^2}}{\displaystyle \frac{1}{2}}(\widehat{Q}_B\mathrm{\Phi })(\widehat{\eta }_0\widehat{\mathrm{\Phi }})+{\displaystyle \frac{1}{6}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\mathrm{\Phi }}(\widehat{\eta }_0\mathrm{\Phi })(\widehat{\eta }_0\widehat{\mathrm{\Phi }}\widehat{\mathrm{\Phi }}))`$
$`+`$ $`{\displaystyle \frac{1}{24}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\mathrm{\Phi }}^2(\widehat{\eta }_0\widehat{\mathrm{\Phi }})2\widehat{\mathrm{\Phi }}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})+(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^2)`$
### 3.2 The tachyon vertex operator
We write down the vertex operator corresponding to the tachyon and the first excited state in the spectrum of the $`(0,4)`$ strings. We work in the gauge
$$b_0\widehat{\mathrm{\Phi }}=0,\xi _0\widehat{\mathrm{\Phi }}=0.$$
(21)
Let us take without loss of generality $`\nu +\nu ^{}<1`$. Then the tachyon vertex operator is given by
$$\widehat{T}=\xi ce^\varphi (t_+\mathrm{\Sigma }_\nu \mathrm{\Sigma }_\nu ^{}\sigma _+\sigma _1+t_{}\mathrm{\Sigma }_\nu \mathrm{\Sigma }_\nu ^{}\sigma _{}\sigma _1)$$
(22)
where
$$\mathrm{\Sigma }_\nu \mathrm{\Sigma }_\nu ^{}=\sigma _\nu e^{i\nu H}\sigma _\nu ^{}e^{i\nu ^{}H}\mathrm{\Sigma }_\nu \mathrm{\Sigma }_\nu ^{}=\sigma _\nu e^{i\nu H}\sigma _\nu ^{}e^{i\nu ^{}H}$$
(23)
The twist operators $`\sigma _{\pm \nu },\sigma _{\pm \nu ^{}}`$ and $`H,H^{}`$ are defined in the appendix. $`\sigma _+`$ and $`\sigma _{}`$ are given by
$$\sigma _+=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)\sigma _{}=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)$$
(24)
The first set of matrices in (22) stand for the external Chan-Paton indices, the second set stand for the internal Chan-Paton indices. The twist operator $`\mathrm{\Sigma }_\nu \mathrm{\Sigma }_\nu ^{}`$ and the anti-twist operator $`\mathrm{\Sigma }_\nu \mathrm{\Sigma }_\nu ^{}`$ comes with external Chan-Paton factors $`\sigma _+`$ and $`\sigma _{}`$ respectively. This is because the insertion of $`\mathrm{\Sigma }_\nu \mathrm{\Sigma }_\nu `$ changes the boundary conditions to that of the string joining the D0-brane to the D4-brane. The mode expansion of these strings are given in (4), (7), (10), (11). While the insertion of $`\mathrm{\Sigma }_\nu \mathrm{\Sigma }_\nu ^{}`$ changes the boundary conditions to that of the string joining the D4-brane to the D0-brane, that is a string of opposite orientation. Note that the conformal dimension of the tachyon vertex agrees with the zero point energy $`E_{}^+`$.
The tachyon vertex operator in (22) reduces to the tachyon vertex of a D-brane anti-D-brane pair when $`\nu =0`$ and $`\nu ^{}=0`$. When $`\nu =0`$ and $`\nu ^{}=0`$ we see from the mode expansions in (4), (7), (10), (11) the boundary condition is Dirichlet at both ends of the $`(0,4)`$ strings. From (5) and (8) we see that $`\nu =0`$ and $`\nu ^{}=0`$ can occur only if $`b\mathrm{}`$ and $`b^{}\mathrm{}`$ or $`b\mathrm{}`$ and $`b^{}\mathrm{}`$. The D0-brane charge induced on the D4-brane is proportional to $`\text{Pf}(2\pi B)=bb^{}`$. The D4-brane effectively becomes an anti-D0-brane. Thus the tachyon vertex in (22) reduces to the tachyon vertex of a D-brane anti-D-brane pair.
For completeness we include the vertex operator of the first excited state. The first excited state becomes the tachyon and the tachyon in (22) becomes the first excited state when $`\nu +\nu ^{}>1`$. The vertex operator for the first excited state is given by
$$\widehat{E}=\xi ce^\varphi (e_+\mathrm{\Sigma }_{1\nu }\mathrm{\Sigma }_{1\nu ^{}}\sigma _+\sigma _1+e_{}\mathrm{\Sigma }_{(1\nu )}\mathrm{\Sigma }_{(1\nu ^{})}\sigma _{}\sigma _1)$$
(25)
where
$`\mathrm{\Sigma }_{1\nu }\mathrm{\Sigma }_{1\nu ^{}}`$ $`=`$ $`\sigma _{1\nu }e^{i(1\nu )H}\sigma _{1\nu ^{}}e^{i(1\nu ^{})H^{}}`$ (26)
$`\mathrm{\Sigma }_{(1\nu )}\mathrm{\Sigma }_{(1\nu ^{})}`$ $`=`$ $`\sigma _{(1\nu )}e^{i(1\nu )H}\sigma _{(1\nu ^{})}e^{i(1\nu ^{})H^{}}`$
To compute the tachyon potential to the zeroth level we take the string field to be
$$\widehat{\mathrm{\Phi }}=\widehat{T}+\widehat{E}$$
(27)
### 3.3 The tachyon potential
In the case of the unstable brane or the D-brane anti-D-brane system the tachyon potential is universal . The tachyon potential in this case is independent of the background conformal field theory except for an overall multiplicative factor. <sup>1</sup><sup>1</sup>1The overall factor is the open string coupling constant. The relationship of the open string coupling constant to the closed string coupling constant can depend on the background conformal field theory. The computation of the tachyon potential for the case involves correlations functions involving the ghosts fields and the matter energy momentum tensor with central charge $`15`$. It is clear from the vertex operator for the tachyon in (22) that it involves matter primaries. Thus the tachyon potential for the $`(0,4)`$ strings does not belong to the same universality class as that of the unstable D-branes. In fact it depends on the background B-field.
The details on the computation of the tachyon potential is given in appendix B. We state the result. The tachyon potential of the $`(0,4)`$ strings at the zeroth level is given by
$`V`$ $`=`$ $`S={\displaystyle \frac{1}{g^2}}({\displaystyle \frac{1}{2}}(1(\nu +\nu ^{})t_{}t_+++{\displaystyle \frac{1}{2}}(1(\nu +\nu ^{})e_{}e_+`$
$`+`$ $`{\displaystyle \frac{1}{F(\nu ,1\nu ,1;1/2)F(\nu ^{},1\nu ^{},1;1/2)}}t_{}^2t_+^2`$
$`+`$ $`{\displaystyle \frac{e_{}^2e_+^2}{F(1\nu ,\nu ,1;1/2)F(1\nu ^{},\nu ^{},1;1/2)}}+{\displaystyle \frac{2t_+t_{}e_+e_{}}{F(\nu ,1\nu ,1;1)F(\nu ^{},1\nu ^{},1;1)}})`$
where the hypergeometric function $`F(\nu ,1\nu ,1;x)`$ is defined in (64). We have the relation
$$F(a,1a,1;1/2)=\sqrt{\pi }\frac{1}{\mathrm{\Gamma }(\frac{1}{2}+\frac{a}{2})\mathrm{\Gamma }(1\frac{a}{2})}$$
(29)
Using this, the tachyon potential becomes
$`V`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}({\displaystyle \frac{1}{2}}(1(\nu +\nu ^{}))t_{}t_++{\displaystyle \frac{1}{2}}(1(\nu +\nu ^{}))e_{}e_+`$
$`+`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{\Gamma }({\displaystyle \frac{1}{2}}+{\displaystyle \frac{\nu }{2}})\mathrm{\Gamma }(1{\displaystyle \frac{\nu }{2}})\mathrm{\Gamma }({\displaystyle \frac{1}{2}}+{\displaystyle \frac{\nu ^{}}{2}})\mathrm{\Gamma }(1{\displaystyle \frac{\nu ^{}}{2}})(t_{}^2t_+^2+e_{}^2e_+^2)`$
$`+`$ $`{\displaystyle \frac{2t_+t_{}e_+e_{}}{F(\nu ,1\nu ,1;1)F(\nu ^{},1\nu ^{},1;1)}})`$
## 4 Analysis of the tachyon potential
The mass defect for D-brane anti-D-brane annihilation is given by $`2\times `$(mass of the D-brane). There is a minimum of the tachyon potential computed by the level truncation approximation scheme which approximates the mass defect. For the D0/D4 system with the Neveu-Schwarz B-field in the spatial direction we expect the tachyon condensate to be the bound state of the D0-brane within the D4-brane. Thus we expect the minimum of the tachyon potential of the $`(0,4)`$ strings to approximate the mass defect of the formation of the D0/D4 bound state from the D0-brane and the D4-brane. In this section we will compare this mass defect with the value of the minimum of the tachyon potential.
### 4.1 The calculation of the mass defect
The mass of the D0-brane in the conventions of is given by
$$M_{D0}=\frac{1}{2\pi ^2g^2}$$
(31)
This formula is true for any Dp-brane. In calculating this mass the longitudinal directions of the Dp-brane is assumed to be compact. The mass of the D4-brane with the B-field in the spatial directions is given by
$$M_{D4}=\frac{1}{2\pi ^2g^2}\sqrt{(1+b^2)(1+b^2)}$$
(32)
We can understand this expression from the Dirac-Born-Infeld action.
$$S=\frac{1}{2\pi ^2g^2}\sqrt{\text{Det}(G+2\pi B)}$$
(33)
where $`G`$ is the induced metric. Substituting the value of the B-field in (1) and the metric $`g_{ij}=\delta _{ij}`$ and expanding the Dirac-Born Infeld action in the static gauge we obtain the mass for the D4-brane with the B-field as given in (32). It is instructive to understand the formula in (32) by another method. The BPS formula for the D0/D4 system with moduli is given by
$$M^2=(Q_0+(1\text{Pf}(2\pi B))Q_4)^2+(2\pi ^2B_{ik}g^{ij}g^{jl}B_{kl}+2\text{Pf}(2\pi B))Q_4^2$$
(34)
Here $`Q_0=N_0/(2\pi ^2g^2)`$ and $`Q_4=N_4/(2\pi ^2g^2)`$. $`N_0`$ and $`N_4`$ are the number of D0-branes and D4-branes respectively. Substituting the values of the B-field from (1) we get
$$M^2=(Q_0+(1bb^{})Q_4)^2+(b+b^{})^2Q_4^2$$
(35)
Substituting $`Q_0=0`$ and $`Q_4=1/(2\pi ^2g^2)`$ in (35) we get (32).
Now it is easy to write down the mass formula for the D0/D4 bound state. Substitute $`Q_0=1/(2\pi ^2g^2)`$ and $`Q_4=1/(2\pi ^2g^2)`$ in (35) we get the following mass formula for the D0/D4 bound state.
$`M_{D0/D4}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2g^2}}\sqrt{(1+(1bb^{}))^2+(b+b^{})^2}`$
$`=`$ $`{\displaystyle \frac{1}{2\pi ^2g^2}}\sqrt{4+b^2+b^22bb^{}+b^2b^2}`$
Thus the mass defect is given by
$$\mathrm{\Delta }M=M_{D0/D4}(M_{D0}+M_{D4})$$
(37)
We analyze the mass defect $`\mathrm{\Delta }M`$ in three limiting cases.
Case 1. $`b1,b^{}1`$
The leading order terms in the mass defect is given by
$$\mathrm{\Delta }M_1=\frac{1}{2\pi ^2g^2}(\frac{b+b^{}}{2})^2$$
(38)
Case 2. $`b\mathrm{},b^{}\mathrm{}`$ or $`b\mathrm{},b^{}\mathrm{}`$; Pf$`(2\pi B)>0`$
The leading order terms in the mass defect is given by
$$\mathrm{\Delta }M_2=2\frac{1}{2\pi ^2g^2}$$
(39)
Notice that in this case the mass defect reduces to that of the D-brane anti-D-brane system. This is expected from the discussion in Section 3.2
Case 3. $`b\mathrm{},b^{}\mathrm{}`$ or $`b\mathrm{},b^{}\mathrm{}`$; Pf$`(2\pi B)<0`$
For this case the leading order terms in the mass defect is given by
$$\mathrm{\Delta }M_3=\frac{1}{4\pi ^2g^2}(\frac{1}{b}+\frac{1}{b^{}})^2$$
(40)
### 4.2 The minimum of the tachyon potential and the mass defect
In this section we compute the minimum of the tachyon potential and compare it with the mass defect obtained in section 4.1. The tachyon potential in (3.3) has an extrema at
$`t_{}t_+`$ $`=`$ $`{\displaystyle \frac{\pi }{4}}{\displaystyle \frac{1(\nu +\nu ^{})}{\mathrm{\Gamma }(1/2+\nu /2)\mathrm{\Gamma }(1\nu /2)\mathrm{\Gamma }(1/2+\nu ^{}/2)\mathrm{\Gamma }(1\nu ^{}/2)}}`$ (41)
$`e_{}=0`$ $`e_+=0`$
The minimum of the tachyon potential at this extrema is given by
$$V_{\mathrm{min}}=\frac{\pi }{16g^2}\frac{(1(\nu +\nu ^{}))^2}{\mathrm{\Gamma }(1/2+\nu /2)\mathrm{\Gamma }(1\nu /2)\mathrm{\Gamma }(1/2+\nu ^{}/2)\mathrm{\Gamma }(1\nu ^{}/2)}$$
(42)
We now compare this to the mass defect obtained in the three cases in section 4.1
Case 1. $`b1,b^{}1`$
To the leading order the value of $`\nu `$ and $`\nu ^{}`$ are given by
$$\nu =\frac{1}{2}\frac{b}{\pi }\nu ^{}=\frac{1}{2}\frac{b^{}}{\pi }$$
(43)
We substitute this values of $`\nu `$ and $`\nu ^{}`$ in (42). We find the minimum of the tachyon to the leading order is given by
$$V_{\mathrm{min}}=\frac{1}{2g^2\pi ^2}\left(\frac{b+b^{}}{2}\right)^2\frac{\pi }{2(\mathrm{\Gamma }(3/4))^4}$$
(44)
On comparison with the mass defect for this case from (38) we obtain $`70\%`$ of the expected result.
$$\frac{V_{\mathrm{min}}}{\mathrm{\Delta }M_1}=\frac{\pi }{2(\mathrm{\Gamma }(3/4)^4}=.70$$
(45)
Case 2. $`b\mathrm{},b^{}\mathrm{}`$ or $`b\mathrm{},b^{}\mathrm{}`$; Pf$`(2\pi B)>0`$
As we have seen in section 3.2 this case is expected to reduce that of the D0-brane anti-D0-brane system. Without loss of generality we take $`b\mathrm{},b^{}\mathrm{}`$. The value of $`\nu `$ and $`\nu ^{}`$ to the leading order are given by
$$\nu =\frac{1}{\pi b}\nu ^{}=\frac{1}{\pi b^{}}$$
(46)
It is easy to see the tachyon potential in (3.3) reduces to that of the D0-brane anti-D0-brane system. The minimum of the tachyon potential is given by
$$V_{\mathrm{min}}=\frac{1}{16g^2}$$
(47)
This is the value of the minimum of the tachyon potential computed in for the D-brane anti-D-brane system. Comparison with the mass defect in (39) gives 62% of the expected result.
$$\frac{V_{\mathrm{min}}}{\mathrm{\Delta }M_2}=\frac{\pi ^2}{16}=.62$$
(48)
Case 3. $`b\mathrm{},b^{}\mathrm{}`$ or $`b\mathrm{},b^{}\mathrm{}`$; Pf$`(2\pi B)<0`$
Without loss of generality we can consider the case $`b\mathrm{},b^{}\mathrm{}`$. For this case the values of $`\nu `$ and $`\nu ^{}`$ to the leading order are given by
$$\nu =1\frac{1}{\pi b}\nu ^{}=\frac{1}{\pi b^{}}$$
(49)
Substituting this in (42) we obtain
$$V_{\mathrm{min}}=\frac{1}{4\pi ^2g^2}\left(\frac{1}{b}+\frac{1}{b^{}}\right)^2\frac{1}{4}$$
(50)
Comparing the value of the minimum of the tachyon potential to the mass defect in (4.1) yields only 25% of the expected result.
$$\frac{V_{\mathrm{min}}}{\mathrm{\Delta }M_2}=\frac{1}{4}=.25$$
(51)
Perhaps the low contribution from the zeroth level approximation to the tachyon potential might be due to the large number of low lying states found for this case .
It is nice to see that we obtain the same dependence of the mass defect on the B-field in Case 1 and Case 2 from two different functions. Graphically it can be easily seen that $`V_{\mathrm{min}}\mathrm{\Delta }M`$ for all values of $`\nu `$ and $`\nu ^{}`$ between $`0`$ and $`1`$. Equality holds only when $`\nu +\nu ^{}=1`$, that is when the tachyon is massless and there is no mass defect. The best approximation to the mass defect over all the allowed values of $`\nu `$ and $`\nu ^{}`$ is 70%.
## 5 Excitations around the tachyon condensate
In this section we examine the excitations of the $`(0,4)`$ strings around the D0/D4 bound state. This is facilitated by the description of the D0/D4 bound state from the gauge theory of the D4-brane. The existence of the bound state was shown from the gauge theory in . The bound state of the D0-brane within the D4-brane can be identified with a single instanton of the D4-brane gauge theory. We will briefly review the arguments here.
The gauge theory of the D4-brane with the B-field is a noncommutative $`U(1)`$ gauge theory with $`N=4`$ supersymmetry. This theory admits non-singular instanton solutions . This noncommutative $`U(1)`$ gauge theory with a single instanton has the right properties to be identified as the bound state the D0-brane within the D4-brane. The arguments for this is the same as the case with no B-field. The instanton in the D4-brane gauge theory carries the Ramond-Ramond charge of the D0-brane. It preserves the same supercharges as the D0-brane . The instanton of the D4-brane gauge theory has the right action as the tension of the D0-brane. The ADHM construction of the instanton of noncommutative D4-brane gauge theory reproduces the moduli space of the D0-brane gauge theory .
### 5.1 The decoupling of the $`(0,4)`$ strings
In this section we show that at the tachyon condensate the $`(0,4)`$ strings become massive and decouple from the dynamics.
We examine the small fluctuations around the tachyon condensate. The tachyon condensate is identified with the instanton of the non-commutative gauge theory. The small fluctuations around the instanton can be determined by examining the moduli space of the instanton of the noncommutative $`U(1)`$ gauge theory. This is given by the ADHM construction of the instanton. The moduli space is determined by solving the following equations
$`AB`$ $`=`$ $`0`$ (52)
$`|A|^2|B|^2`$ $`=`$ $`r`$
These are the D-terms of a $`(4,4)`$ supersymmetric $`U(1)`$ gauge theory with one hypermultiplet, that is the gauge theory of the D0-brane. $`r`$ is the Fayet-Illiopoulous term. It has been shown that a bound state exists for $`r>0`$ . For Case 1 and Case 2 $`r`$ has been identified to be proportional to the square of the self dual part of the B-field . $`A`$ and $`B`$ are the bosonic components of the hypermultiplet. $`A`$ is is a complex field which carries a charge of $`+1`$ under the $`U(1)`$ gauge symmetry and $`B`$ carries a charge of $`1`$ under the $`U(1)`$ gauge symmetry. To determine the moduli space we must solve the equations in (52) and identify gauge equivalent configurations.
The solution of (52) is given by $`B=0`$ with $`|A|=\sqrt{r}`$. Using the $`U(1)`$ gauge degree of freedom one can choose the solution $`B=0`$ and $`A=\sqrt{r}`$. Thus we see there is no modulus for the non-commutative instanton other than its position.
Small fluctuations of the $`(0,4)`$ strings can be identified with the small fluctuations of the fields $`A`$ and $`B`$. This can be seen from the fact that the fields corresponding to the $`(0,4)`$ strings are hypermultiplets in the gauge theory of the D4-brane and the D0-brane. Therefore at the tachyon condensate, the $`(0,4)`$ strings are massive and decouple from the low energy dynamics. The decoupling of the $`(0,4)`$ strings is the analog of the decoupling of the strings joining the D-brane with the anti-D-brane. It is the analog of the decoupling of the relative $`U(1)`$ for the D-brane, anti-D-brane case.
## 6 Conclusions
We have studied tachyon condensation in the D0/D4 system with B-field in the spatial directions of the D4-brane using string field theory. The tachyon in the spectrum of the $`(0,4)`$ strings signals the instability of the D0/D4 system to form the D0/D4 bound state. Here after tachyon condensation we are left with a BPS state. We computed the tachyon potential in the zeroth level approximation. This tachyon potential is outside the universality class analyzed for the D-brane anti-D-brane systems. It is a function of the background B-field. We compared the minimum of the tachyon potential to the mass defect for the formation of the D0/D4 bound state in three cases. When the tachyon is almost on shell for small values of the B-field the zeroth level approximation for the tachyon potential contributes to 70% of the expected mass defect. For large values of the B-field with Pf$`(2\pi B)`$ the tachyon potential reduces to that of the D-brane anti-D-brane pair. For large values of the B-field with Pf$`(2\pi B)`$ the zeroth level approximation contributes only to 25% of the mass defect. Note that in this last case the level truncation method is not very successful at the zeroth level. It would be interesting to know whether the inclusion of higher levels for this case converges to the expected answer.
We see that for the case of large values of the B-field with Pf$`(2\pi B)>0`$ the D0/D4 system reduces to that of the D-brane anti-D-brane system. One is tempted to think that the decoupling of the open string modes in tachyon condensation of the D-brane anti-D-brane system can be understood from the description of the D0/D4 bound state as an instanton in the noncommutative gauge theory of the D4-brane for large values of the B-field.
###### Acknowledgments.
The author thanks Shinji Hirano, Nissan Itzhaki and especially Shiraz Minwalla for discussions. He is grateful for encouragement from Joe Polchinski and Ashoke Sen. The work of the author is supported by NSF grant PHY97-2202.
## Appendix A The twist operator and its correlation function
This appendix discusses the definition of the bosonic and fermionic twist operators and calculates their correlations functions. The twist operators are located at the boundary of the world sheet. Thus they are on the real axis when the world sheet is on the upper half plane. The bosonic twist operator $`\sigma _\nu `$ inserted at the origin changes the boundary conditions of the open string as shown in the Figure 1. The insertion of $`\sigma _\nu `$ changes the boundary conditions of the string to that of a string joining the D0-brane to the D4-brane. The anti-twist operator $`\sigma _\nu `$ changes the boundary conditions to that of a string joining the D4-brane to the D0-brane. From the mode expansions in (4) and (7) the bosonic twist operators $`\sigma _\nu `$ and $`\sigma _\nu `$ have the following OPE with the world sheet bosons
$`X^+(z)\sigma _\nu (w)={\displaystyle \frac{1}{(zw)^{(1\nu )}}}\tau _\nu (w)`$ $`X^{}(z)\sigma _\nu (w)={\displaystyle \frac{1}{(zw)^\nu }}\tau _\nu ^{}(w)`$ (53)
$`X^+(z)\sigma _\nu (w)={\displaystyle \frac{1}{(zw)^\nu }}\tau _\nu (w)`$ $`X^{}(z)\sigma _\nu (w)={\displaystyle \frac{1}{(zw)^{(1\nu )}}}\tau _\nu ^{}(w)`$
$`\overline{}X^+(\overline{z})\sigma _\nu (w)={\displaystyle \frac{1}{(\overline{z}w)^{1\nu }}}\tau _\nu (w)`$ $`\overline{}X^{}(\overline{z})\sigma _\nu (w)={\displaystyle \frac{1}{(\overline{z}w)^\nu }}\tau _\nu ^{}(w)`$
$`\overline{}X^+(\overline{z})\sigma _\nu (w)={\displaystyle \frac{1}{(\overline{z}w)^\nu }}\tau _\nu (w)`$ $`\overline{}X^{}(\overline{z})\sigma _\nu (w)={\displaystyle \frac{1}{(\overline{z}w)^{(1\nu )}}}\tau _\nu ^{}(w)`$
$`w`$ is a point on the real axis. $`\tau `$’s are the excited twist operators. We are interested in obtaining the following four point function of the twist fields
$`Z_1(z_1,z_2,z_3,z_4)`$ $`=`$ $`\sigma _\nu (z_1)\sigma _\nu (z_2)\sigma _\nu (z_3)\sigma _\nu (z_4)`$ (54)
$`Z_2(z_1,z_2,z_3,z_4)`$ $`=`$ $`\sigma _{1\nu }(z_1)\sigma _{(1\nu }(z_2)\sigma _\nu (z_3)\sigma _\nu (z_4)`$
where $`z_1,z_2,z_3,z_4`$ are four points on the real axis. We will demonstrate the calculation for $`Z_1`$, the calculation of $`Z_2`$ follows along similar lines. We follow the method in and as developed for open strings in . Consider the auxiliary correlators
$$g(z,w)=\frac{\frac{1}{2}X^+(z)X^{}(w)\sigma _\nu (z_1)\sigma _\nu (z_2)\sigma _\nu (z_3)\sigma _\nu (z_4)}{\sigma _\nu (z_1)\sigma _\nu (z_2)\sigma _\nu (z_3)\sigma _\nu (z_4)}$$
(55)
and
$$h(\overline{z},w)=\frac{\frac{1}{2}\overline{}X^+(\overline{z})X^{}(w)\sigma _\nu (z_1)\sigma _\nu (z_2)\sigma _{\nu (z_3)}\sigma _\nu (z_4)}{\sigma _\nu (z_1)\sigma _\nu (z_2)\sigma _{\nu (z_3)}\sigma _\nu (z_4)}$$
(56)
Define
$$\omega _\nu (z)=\frac{1}{[(zz_1)(zz_3)]^\nu }\frac{1}{[(zz_2)(zz_4)]^{1\nu }}$$
(57)
Now $`g(z,w)`$ is given by
$`g(z,w)`$ $`=`$ $`\omega _\nu (z)\omega _{1\nu }(w)(\nu {\displaystyle \frac{(zz_1)(zz_3)(wz_2)(wz_4)}{(zw)^2}}`$
$`+`$ $`(1\nu ){\displaystyle \frac{(zz_2)(zz_4)(wz_1)(wz_3)}{(zw)^2}}+A(z_1,z_2,z_3,z_4))`$
The form for $`g(z,w)`$ given above has the required singularity structure so that the (53) is obeyed. When $`zw`$, then
$$g(z,w)=\frac{1}{(zw)^2}$$
(59)
$`h(\overline{z},w)`$ is given by
$`h(\overline{z},w)`$ $`=`$ $`\omega _\nu (\overline{z})\omega _{1\nu }(w)(\nu {\displaystyle \frac{(\overline{z}z_1)(\overline{z}z_3)(wz_2)(wz_4)}{(zw)^2}}`$
$`+`$ $`(1\nu ){\displaystyle \frac{(\overline{z}z_2)(\overline{z}z_4)(wz_1)(wz_3)}{(zw)^2}}+A(z_1,z_2,z_3,z_4))`$
This form for $`h(\overline{z},w)`$ also has the singularity structure so that the (53) is obeyed. It satisfies the condition
$$h(\overline{z},w)=g(\overline{z},w)$$
(61)
The origin of this condition can be seen from the mode expansions (4) and (7). The anti-holomorphic components $`\overline{}X^\pm (z)`$ can be obtained from the holomorphic components by replacing $`z`$ by $`\overline{z}`$ along with a negative sign. We now determine $`A`$ from the monodromy conditions. We have the monodromy condition <sup>2</sup><sup>2</sup>2We are neglecting instanton sectors as the D4-brane is extended in the spatial directions.
$$_Cg(z,w)𝑑z+_Ch(\overline{z},w)𝑑\overline{z}=0$$
(62)
Where the contour $`C`$ is shown in Figure 2. This monodromy condition arises because integration along the contour $`C`$ gives the change in $`X^+`$ for a string ending on the D0-brane at the two ends. There are two equivalent nontrivial contours $`C`$ and $`C^{}`$.
To solve for $`A`$ first divide (62) by $`\omega _{1\nu }(w)`$ and take $`w\mathrm{}`$. We then use the $`SL(2,R)`$ invariance to choose $`z_1=0,z_2=x,z_3=1`$ and $`z_4=\mathrm{}`$, where $`0<x<1`$. Then $`A`$ is given by
$$A=z_4x(1x)\frac{d}{dx}\mathrm{ln}F(\nu ,1\nu ,1;x)$$
(63)
Here $`z_4`$ stands for $`\mathrm{}`$. $`F(\nu ,1\nu ;x)`$ is the hypergeometric function whose integral representation is given by
$$F(\nu ,1\nu ,1;x)=\frac{1}{\pi }\mathrm{sin}(\pi \nu )_0^1𝑑y\frac{1}{y^\nu (1y)^{1\nu }(1xy)^\nu }$$
(64)
Now we can calculate the correlation function $`Z_1(z_1,z_2,z_3,z_4)`$. Consider the following correlation function
$`{\displaystyle \frac{T(z)\sigma _\nu (z_1)\sigma _\nu (z_2)\sigma _\nu (z_3)\sigma _\nu (z_4)}{\sigma _\nu (z_1)\sigma _\nu (z_2)\sigma _\nu (z_3)\sigma _\nu (z_4)}}`$ $`=`$ $`\underset{wz}{lim}\left[g(z,w){\displaystyle \frac{1}{(zw)^2}}\right]`$
$`=`$ $`{\displaystyle \frac{1}{2}}\nu (1\nu )({\displaystyle \frac{1}{(zz_1)^2}}+{\displaystyle \frac{1}{(zz_3)^2}}`$
$``$ $`{\displaystyle \frac{1}{(zz_3)^2}}{\displaystyle \frac{1}{(zz_4)^2}})`$
$`+`$ $`{\displaystyle \frac{A}{(zz_1)(zz_2)(zz_3)(zz_4)}}`$
where $`T`$ is the stress energy tensor. From the operator product
$$T(z)\sigma _\nu (z_2)=\frac{1}{2}\frac{\nu (1\nu )T(z_2)}{(zz_2)^2}+\frac{_{z_2}\sigma _\nu (z_2)}{zz_2}$$
(66)
we can obtain the following differential equation for $`Z_1(z_1,z_2,z_3,z_4)`$
$`_{z_2}\mathrm{ln}Z_1(z_1,z_2,z_3,z_4)`$ $`=`$ $`\nu (1\nu )\left({\displaystyle \frac{1}{(z_2z_1)}}+{\displaystyle \frac{1}{(z_2z_3)}}{\displaystyle \frac{1}{(z_3z_4)}}\right)`$
$`+`$ $`{\displaystyle \frac{A}{(z_2z_1)(z_2z_3)(z_2z_4)}}`$
Substituting the values of $`z_1,z_2,z_3,z_4`$ in the above equation we obtain
$$_xZ_1(x)=\nu (1\nu )\left(\frac{1}{x}\frac{1}{1x}\right)\frac{d}{dx}\mathrm{ln}F(\nu ,1\nu ,1;x)$$
(68)
This equation can easily be integrated to obtain
$$Z_1(x)=\frac{1}{[x(1x)]^{\nu (1\nu )}F(\nu ,1\nu ,1,x)}$$
(69)
We now use the $`SL(2,R)`$ invariance to obtain $`Z_1(z_1,z_2,z_3,z_4)`$. We get
$$Z_1(z_1,z_2,z_3,z_4)=z_{21}^{2h}z_{31}^{2h}z_{41}^{2h}z_{32}^{2h}z_{42}^{2h}z_{43}^{2h}\frac{1}{F(\nu ,1\nu ,1;x)}$$
(70)
where $`h=\nu (1\nu )/2`$, $`z_{ij}=z_iz_j`$ and
$$x=\frac{z_{21}z_{43}}{z_{31}z_{42}}$$
(71)
To construct the fermionic twist operators we first bosonize the fermions by
$$\psi ^+(w)=e^{iH(w)}\psi ^{}(w)=e^{iH(w)}$$
(72)
where $`H`$ is a free boson. The fermionic twist operator are is given by $`e^{i\nu H}`$ and the anti-twist operator is given by $`e^{i\nu H}`$. The correlation function for the fermionic twists are easy to evaluate and they are given by
$$e^{i\nu H(z_1)}e^{i\nu H(z_2)}e^{i\nu H(z_3)}e^{i\nu H(z_4)}=z_{21}^{\nu ^2}z_{31}^{\nu ^2}z_{41}^{\nu ^2}z_{32}^{\nu ^2}z_{42}^{\nu ^2}z_{43}^{\nu ^2}$$
(73)
Putting the bosonic and the fermionic twists together one obtains
$$\mathrm{\Sigma }_\nu (z_1)\mathrm{\Sigma }_\nu (z_2)\mathrm{\Sigma }_\nu (z_3)\mathrm{\Sigma }_\nu (z_4)=z_{21}^\nu z_{31}^\nu z_{41}^\nu z_{32}^\nu z_{42}^\nu z_{43}^\nu \frac{1}{F(\nu ,1\nu ,1,x)}$$
(74)
Using similar methods we obtain the correlation function $`Z_2(z_1,z_2,z_3,z_4)`$. It is given by
$`Z_2(z_1,z_2,z_3,z_4)=z_{21}^{2h}z_{31}^{2h}z_{41}^{2h}z_{32}^{2h}z_{42}^{2h}z_{43}^{2h}{\displaystyle \frac{1}{F(\nu ,1\nu ,1;y)}}`$ (75)
Here $`y=x/(x1)`$. Using $`Z_2`$ we can find the following correlation functions
$`\mathrm{\Sigma }_{1\nu }(z_1)\mathrm{\Sigma }_{(1\nu )}(z_2)\mathrm{\Sigma }_\nu (z_3)\mathrm{\Sigma }_\nu (z_4)`$ $`=`$ $`{\displaystyle \frac{1}{z_{21}^{1\nu }z_{43}^\nu }}{\displaystyle \frac{1}{F(\nu ,1\nu ,1;y)}}`$ (76)
$`\mathrm{\Sigma }_\nu (z_1)\mathrm{\Sigma }_\nu (z_2)\mathrm{\Sigma }_{(1\nu )}(z_3)\mathrm{\Sigma }_{1\nu }(z_4)`$ $`=`$ $`{\displaystyle \frac{1}{z_{21}^\nu z_{43}^{1\nu }}}{\displaystyle \frac{1}{F(\nu ,1\nu ,1;y)}}`$
For completeness we write down the two point function of the twist operators. This is fixed by conformal invariance.
$$\mathrm{\Sigma }_\nu (z_1)\mathrm{\Sigma }_\nu (z_2)=\frac{1}{(z_2z_1)^\nu }$$
(77)
We have chosen a normalization for the two point function which differs from the conventional one to make sure all coefficients are real in the tachyon potential. The four point function of the twist operators in (74) and (76) are consistent with this normalization of the two point function <sup>3</sup><sup>3</sup>3If we had worked with the conventional normalization we can ensure that the coefficients of the tachyon potential are real by a redefinition of the tachyon field..
## Appendix B Details on the calculation of the tachyon potential
The vertex operators $`\widehat{T}`$ and $`\widehat{E}`$ are all primary, therefore they transform under a conformal transformation $`f`$ as
$$f𝒪(0)=(f^{}(0))^h𝒪(0)$$
(78)
Here $`𝒪`$ is a primary operator of dimension $`h`$.
### B.1 The quadratic term
We now focus on the evaluation of the quadratic term in the string field theory action given by
$$S^{(2)}=\frac{1}{4g^2}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\eta }_0\widehat{\mathrm{\Phi }})$$
(79)
Substituting $`\widehat{\mathrm{\Phi }}=\widehat{T}+\widehat{E}`$, the terms which contribute in $`S^{(2)}`$ are
$$S^2=\frac{1}{4g^2}\left((\widehat{Q}_B\widehat{T})(\widehat{\eta }_0\widehat{T})+(\widehat{Q}_B\widehat{E})(\widehat{\eta }_0\widehat{E})\right)$$
(80)
The cross terms do not contribute due to twist conservation. $`\widehat{T}`$ is an operator of dimension $`(1(\nu +\nu ^{}))/2`$ and $`\widehat{E}`$ is an operator of dimention $`(1(\nu +\nu ^{}))/2`$. We calculate the correlations functions on the upper half plane. Using (77) and evaluating the traces over the Chan-Paton factors and grouping the terms we obtain
$`S^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}(+{\displaystyle \frac{1}{2}}(1(\nu +\nu ^{}))z_{21}^{1(\nu +\nu ^{})}(g_1^{(2)}(0)g_2^{(2)}(0))^{(1(\nu +\nu ^{})/2}t_{}t_+`$
$``$ $`{\displaystyle \frac{1}{2}}(1(\nu +\nu ^{}))z_{21}^{1+(\nu +\nu ^{})}(g_1^{(2)}g_2^{(2)})^{(1(\nu +\nu ^{})/2}e_{}e_+)`$
where $`z_1=g_1^{(2)}(0)`$ and $`z_2=g_2^{(2)}(0)`$. The $`g`$’s are defined in . We write them down here for completeness
$`g_1^{(2)}=\mathrm{tan}({\displaystyle \frac{\pi }{4}})`$ $`g_1^{(2)}=2`$ (82)
$`g_2^{(2)}=\mathrm{tan}({\displaystyle \frac{\pi }{4}})`$ $`g_1^{(2)}=2`$
Substituting the values of each of the terms in $`S^{(2)}`$ we obtain
$$S^{(2)}=\frac{1}{g^2}\left(\frac{1}{2}(1(\nu +\nu ^{}))t_{}t_+\frac{1}{2}(1(\nu +\nu ^{}))e_{}e_+\right)$$
(83)
### B.2 The cubic term
The cubic term in the superstring field theory action is given by
$$S^{(3)}=\frac{1}{12g^2}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\mathrm{\Phi }}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})(\widehat{\eta }_0\widehat{\mathrm{\Phi }}\widehat{\mathrm{\Phi })}$$
(84)
Substituting $`\widehat{\mathrm{\Phi }}=\widehat{T}+\widehat{E}`$ in $`S^{(3)}`$ we see that it vanishes. This can be easily seen by taking the traces of the external Chan-Paton factors.
### B.3 The quartic term
The quartic term in the string field theory action is given by
$$S^{(4)}=\frac{1}{48g^2}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\mathrm{\Phi }}^2(\widehat{\eta }_0\widehat{\mathrm{\Phi }})2\widehat{\mathrm{\Phi }}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})+(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^2)$$
(85)
From the fact that there should be only three $`c`$ ghosts for the correlation functions to contribute we can see that after the substitution $`\widehat{\mathrm{\Phi }}=\widehat{T}+\widehat{E}`$ there is a factorization of the action given by
$$S^{(4)}=𝒮\times 𝒯$$
(86)
where
$`𝒮`$ $`=`$ $`{\displaystyle \frac{1}{48g^2}}(Q_B\sigma _3O\sigma _1)((O\sigma _1)^2(\eta _0\sigma _3O\sigma _1)`$
$``$ $`2O\sigma _1(\eta _0\sigma _3O\sigma _1)+(\eta _0\sigma _3O\sigma _1)(O\sigma _1)^2)`$
$`𝒯`$ $`=`$ $`PPPP`$
Here $`O=\xi ce^\varphi `$ and
$`P`$ $`=`$ $`t_+\mathrm{\Sigma }_\nu \mathrm{\Sigma }_\nu ^{}\sigma _++t_{}\mathrm{\Sigma }_\nu \mathrm{\Sigma }_\nu ^{}\sigma _{}`$
$`+`$ $`e_+\mathrm{\Sigma }_{1\nu }\mathrm{\Sigma }_{1\nu ^{}}\sigma _++e_{}\mathrm{\Sigma }_{(1\nu )}\mathrm{\Sigma }_{(1\nu ^{})}\sigma _{}`$
The value of $`𝒮`$ has been calculated in . $`𝒮=1/(2g^2)`$. We have to use the correlation functions found in Appendix A to evaluate $`𝒯`$. As an example we write down the coefficient of $`(t_{}t_+)^2`$
$`2\left(g_1^{(4)}(0)g_2^{(4)}(0)g_3^{(4)}(0)g_3^{(4)}(0)\right)^{(\nu +\nu ^{})/2}`$ (89)
$`\mathrm{\Sigma }_\nu (z_1)\mathrm{\Sigma }_\nu (z_2)\mathrm{\Sigma }_\nu (z_3)\mathrm{\Sigma }_\nu (z_4)\mathrm{\Sigma }_\nu ^{}(z_1)\mathrm{\Sigma }_\nu ^{}(z_2)\mathrm{\Sigma }_\nu ^{}(z_3)\mathrm{\Sigma }_\nu ^{}(z_4)`$
where $`z_1=g_1^{(4)}(0)z_2=g_2^{(4)}(0)z_3=g_3^{(4)}(0)z_4=g_4^{(4)}(0)`$. and the $`g`$’s are defined in They are given by
$`g_1^{(4)}(z)`$ $`=`$ $`4+6z9z^2+\mathrm{}`$ (90)
$`g_2^{(4)}(z)`$ $`=`$ $`1+{\displaystyle \frac{3}{4}}z{\displaystyle \frac{3}{16}}z^2+\mathrm{}`$
$`g_3^{(4)}(z)`$ $`=`$ $`0+{\displaystyle \frac{2}{3}}z+{\displaystyle \frac{1}{9}}z^2+\mathrm{}`$
$`g_4^{(4)}(z)`$ $`=`$ $`2+3z+3z^2+\mathrm{}`$
Substituting these values in (89) and using the correlation function in (74) we find that the coefficient is that given in (3.3). Note that the value of argument $`x`$ of the hypergeometric function is $`1/2`$ and $`y=1`$.
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# Force fluctuation in a driven elastic chain
## I Introduction
The dissipative motion of an elastic line in a random potential is an interesting example of a nonequilbrium interacting system and is relevant for several phenomena in condensed matter physics. Examples are numerous and include the motion of magnetic interfaces in ferromagnetic materials , solid friction , wetting , charge density waves , fluids in porous media , vortex dynamics in high temperature superconductors , cracks and dislocations. These systems are characterized by a dynamic phase transition ruled by the interplay between quenched disorder and elastic interactions.
Due to the effect of the disorder, an elastic chain at zero temperature is pinned when the applied force is below a critical value $`F_c`$: after a sufficiently long time, independently of the initial conditions, the chain reaches a configuration where no movement is possible. For $`F>F_c`$ the chain can escape from any pinning configuration and moves with constant average velocity. When $`F`$ is close to $`F_c`$ the motion is dominated by collective effects and the depinning of a single bead produces a large reorganization of the chain. In other words, for $`F=F_c`$ the system is critical and the motion of the beads is highly correlated.
An elastic chain moving in a disordered potential is a useful model to understand some general features of sliding friction in particular of the experiment reported in Ref. , done using two artificial surfaces with controlled roughness and elasticity. Beads of diameter 2mm were randomly put inside an elastic matrix, with a maximum roughness of 0.5mm. The two surfaces were then displaced against each other at constant velocity and the friction force was measured, varying the elasticity of the matrix and the driving velocity. The distribution of the amplitude of the slip events is generally found to decay as a power law at small velocities, suggesting the presence of an underlying critical point. The exponents characterizing the power law distribution are found to decrease with the applied velocity, in analogy with other driven systems such as domain walls in ferromagnet .
Several variant of the chain model can be studied in order to reproduce the experiments: periodic or disordered arrangements of beads, on a rigid or elastic substrate. Simulations of a one dimensional, periodic or disordered, chain over a rigid disordered potential have been performed by Cule and Hwa . The measurement of the velocity and the roughness exponents seem to indicate that periodic and disordered chains are described by two different universality classes. These and other simulations are performed considering a constant applied force, while experiments are performed driving the system at constant velocity.
Here we consider explicitly the second case and measure the force fluctuations as a function of the applied velocity. The distribution of the slips events is power law distributed, and characterized by an exponent $`\tau `$, which appears to decrease with the applied velocity, in agreement with the experiments. In the limit of low velocity the exponent $`\tau `$ can be related to the critical exponent obtained tuning the applied force as first discussed in Ref (see also ). We measure the fluctuation in the position of the beads in the constant force and constant velocity cases and obtain the same roughness exponent $`\zeta `$. Next we evaluate the force velocity diagram, using finite size scaling to locate the critical force and compute the exponent $`\beta `$. The values of the exponents are consistent with scaling relation and in good numerical agreement with the exponents of interface depinning, but disagree with previous simulations for a periodic chain in a random potential . In order to confirm this conclusion, we simulate the motion of a disordered chain and see no evidence for the existence of two different universality classes for periodic and disordered chains, in disagreement with the conclusions of Ref. .
The paper is organized as follows in Sec. II we introduce the model, and in Sec III we define the critical exponents and discuss some scaling relations. In Sec. IV A we present the numerical results obtained at constant velocity, in Sec. IV B we discuss the constant force case studying the scaling behavior close to the depinning transition. In Sec. V we summarize the main results of the paper.
## II Model
We consider the overdamped dynamics of a one dimensional line of elastically coupled beads, driven on a disordered substrate. The disordered potential is a succession of identical Gaussian potentials, randomly distributed in space. The beads can be driven directly, applying a constant force, or indirectly coupling them to an intermediate spring which is pulled at constant velocity. The equation of motion is,
$$\eta \frac{r_i(t)}{t}=D(r_{i+1}2r_i+r_{i1})+f(r_i(t))+F,$$
(1)
where $`r_i(t)`$ is the position of the bead $`i`$ at time $`t`$, $`\eta `$ is the coefficient of viscosity, $`D`$ is the stiffness of the elastic line, $`F`$ is the driving force, and $`f(x)`$ is a random force, due to the contribution of the ensemble of pinning centers. We model the random force by the sum of $`N`$ derivatives of Gaussian potentials located on the pinning site
$$f(x)=C\underset{i=1}{\overset{N}{}}(xx_i^p)\mathrm{exp}\left[\frac{1}{2}\frac{(xx_i^p)^2}{\sigma ^2}\right],$$
(2)
where $`C`$ represents the strength of the disorder, $`\sigma `$ quantifies the width of the wells, and $`x_i^p`$ is the location of the pinning site $`i`$, which we chose to be Poisson distributed. The equation of motion is integrated numerically using a fourth-order Runge-Kutta method.
The interplay between disorder and elastic interactions in our model can be understood computing the Larkin length $`l_L`$ . For distances smaller than $`l_L`$, the beads are interacting strongly and the chain moves coherently, while for distances larger than $`l_L`$ the random forces become dominant and the chain deforms considerably. The Larkin length can be estimated considering both the effect of the rigidity of the line and the strength of the disorder and for our model it is given by
$$l_L[\frac{Da\sigma ^2}{\rho ^{1/2}C}]^{2/3},$$
(3)
where $`a`$ is the distance between the beads and $`\rho `$ is the density of pinning centers. In order to analyze the critical properties of the system, we have to consider the limit where the Larkin length is larger than the mean distance between the beads and the dynamics is governed by the collective motion of the beads. To this end we carefully choose the parameters of the model so that $`l_La`$.
## III Scaling relations
Depending on the method used to drive the chain the measured quantities change, but the corresponding critical exponents can be related by scaling relations. Here we summarize the scaling properties of the depinning transition in the case where the beads are driven by springs of stiffness $`K`$ pulled at constant velocity $`V`$ and in the case where they are submitted to a constant force $`F`$.
### A Constant velocity driving
Friction experiments are generally performed under a constant velocity driving. In some cases, a traction machine turning at constant velocity is coupled by a spring to the sliding system. In other cases, an effective spring coupling is due to the elastic deformation of the material driven imposing a constant strain rate far from the sliding interface. To simulate constant velocity driving , we attach each bead to a spring of stiffness $`K`$ , so that the force is given by
$$F=K(r_i(0)+Vtr_i(t))$$
(4)
where $`V`$ is the applied velocity. When the velocity $`V`$ and the stiffness $`K`$ are small, the motion displays large fluctuations. In particular, in the limit $`V0`$ and $`K0`$ the system reaches the depinning transition and the force fluctuates around $`F_c`$. This feature is common to other non equilibrium critical phenomena, such as absorbing state phase transitions and self-organized criticality .
The distribution of the friction force fluctuations is directly related to the size distribution of slip events $`\mathrm{\Delta }x`$ of the chain, since $`\mathrm{\Delta }F=K\mathrm{\Delta }x`$. Here $`x_ir_i`$ and the slip is defined by a drop in the measured friction force (See Fig. 1). Close to the depinning transition, we expect that the distribution of $`\mathrm{\Delta }x`$ decays as
$$P(\mathrm{\Delta }x)s^\tau h(\mathrm{\Delta }x/\mathrm{\Delta }x_0),$$
(5)
where $`h(x)`$ is a scaling function and $`\mathrm{\Delta }x_0`$ the cutoff value. The value of the cutoff depends on various parameters, such as the system size employed in the simulation. Clearly the avalanche size cannot be greater than the total length of the line. Furthermore, we expect that the stiffness of the springs $`K`$ and the driving velocity $`V`$ will in general change the value of the cutoff. In the low velocity limit and for large enough system sizes, $`K`$ becomes the dominant parameter that determines the value of the cutoff .
The scaling of the cutoff with $`K`$ can be used to evaluate the roughness exponent $`\zeta `$, using the relation
$$\mathrm{\Delta }x_0K^{\zeta /2}.$$
(6)
This relation can be obtained noting that the chain can not jump over a distance larger than $`\xi _0`$, which represents the length for which the elastic interaction term is overcame by the restoring force due to the springs: $`D\xi _0^2K`$. Then, using the scaling relation for $`\mathrm{\Delta }x_0\xi _0^\zeta `$, we obtain Eq (6) . This result implies that $`P(\mathrm{\Delta }x,K)`$ should satisfy the scaling form
$$P(\mathrm{\Delta }x,K)K^{\tau \zeta /2}=H(t),t\mathrm{\Delta }xK^{\zeta /2}.$$
(7)
Eq. (5) is used to compute $`\tau `$ and Eq. (6) and Eq. (7) are used to compute the exponent $`\zeta `$.
The exponent $`\zeta `$ can also be evaluated directly using the scaling of the fluctuations of the displacements with $`K`$. Defining the relative displacements of the beads as $`u_i(t)r_i(t)(Vt+r_i(0))`$, the fluctuations can be quantified by
$$W=\underset{i=0}{\overset{L}{}}(u_i(t)m_i(t))^2/L,$$
(8)
where $`m_i(t)_{i=0}^Lu_i(t)/L`$, and $`L`$ is the number of beads. The roughness $`W`$ scales with the correlation length $`\xi _0`$ as $`W\xi _0^\zeta `$, and since $`\xi _0K^{1/2}`$
$$WK^{\zeta /2}.$$
(9)
Eq. (6) and Eq. (9) have the same origin and can independently be used to estimate $`\zeta `$.
### B Constant force driving
When the system is driven at constant force we expect a depinning transition as a function of $`F`$. For $`F>F_c`$, the chain moves with constant average velocity $`v`$ defining an exponent $`\beta `$
$$v(FF_c)^\beta .$$
(10)
Close to the depinning transition the motion is very irregular, and large regions of the chain move collectively. The correlation length diverges at the transition as
$$\xi (FF_c)^\nu .$$
(11)
In order to estimate the critical exponents $`\beta `$ and $`\nu `$ we employ a particular finite size scaling method , in analogy with absorbing state phase transitions .
We first compute the critical force analyzing the decay of the average velocity with time for different system sizes. For finite systems, the average velocity reaches a quasi steady state $`v(F,L)`$. When $`F>F_c`$ we expect that as the size $`L\mathrm{}`$, $`v(F,L)`$ approaches a non vanishing value given by Eq. (10), while it decays to zero for $`F<F_c`$. At the depinning transition we expect that
$$v(F_c,L)L^{\beta /\nu }.$$
(12)
Once $`F_c`$ is known with good precision, we can measure directly the exponent $`\beta `$ from Eq. (10). As in the constant velocity case, we can evaluate the roughness exponent, measuring the width at $`F_c`$ as a function of $`L`$, which should scale as
$$W(L)L^\zeta .$$
(13)
For a periodic chain Ref. reported $`\beta 0.4`$ and $`\zeta 1.5`$, while for a disordered $`\beta 0.25`$ and $`\zeta 1.2`$. These last value are consistent with interface depinning that in $`d=1`$ yields $`\beta 0.25`$ and $`\zeta 1.25`$ .
## IV Numerical results
### A Constant velocity.
The primary interest of this study is to compute the exponent $`\tau `$ which characterizes the collective motion of the particles at $`FF_c`$, and in particular its dependence on the driving velocity. We note that the avalanche exponent was found to decrease with the driving velocity in the Barkhausen effect, due to the motion of domain walls in a ferromagnet . The same effect was observed in the friction experiments reported in Ref .
In order to reach the scaling regime, we progressively decrease $`V`$ and $`K`$ and compute the friction force. Fig. 1 shows two typical plot for the friction force as a function of the position of the line, in Fig. 1a the driving velocity is $`V=0.05`$ and in Fig. 1b $`V=5`$. As the driving velocity increases the friction force becomes smoother, and in the limit $`V1`$ we obtain a viscous behavior ($`F\eta V`$) with small relative fluctuations. On the contrary, for small velocities the dynamics is jerky: the force increases with time until the beads are sufficiently stressed so that the chain depins decreasing the force. In this case, the friction force displays a characteristic stick-slip pattern.
We thus measure the friction force drops $`\mathrm{\Delta }F`$, or the slip sizes $`\mathrm{\Delta }X=\mathrm{\Delta }F/K`$ and analyze their distribution. The distributions are averaged over 20 realizations of the disorder for $`V=0.05`$ and 100 realizations for $`V=10`$; in all cases the system was composed of $`L=1000`$ beads, and the disorder was produced by $`N=20000`$ pining sites, Poissonian distributed. The value of the exponent $`\tau `$ is obtained by a direct fit of the linear part of the distribution plotted in a log-log graph. The results are shown in Fig. 2, the main graph presents the log-log plot of the probability distribution function of the jumps for various $`V`$ and in the inset we report the value of $`\tau `$ as function of the velocity. We see a slow decrease of the exponent when the velocity increases. This result is in good qualitative agreement with the experiments reported in Ref. . The value of $`\tau `$ for $`V0`$ agrees well with the exponent obtained in elastic line depinning under quasistatic conditions (see table I).
Fig. 2 shows that the cutoff $`\mathrm{\Delta }X_0`$ clearly depends on $`V`$ when $`K`$ is hold fixed. We can quantify this variation and the result is reported in Fig. 3. For small velocities ($`V<V^{}`$) the cutoff is a constant (i.e it does not depend on $`V`$) which in principle depends on $`K`$, while for high velocities the cutoff decreases with $`V`$, roughly as a power law. Next, we study the behavior of the cutoff when $`V<V^{}`$ as $`K`$ is varied. In Fig. 4 we shows the distribution of slip sizes for various $`K`$ for very small driving velocity. We see that the cut off increases as $`K`$ is decreased. Fig. 5 shows the collapse of the curves after the rescaling with $`K`$, in accordance with Eq. (7).
We also measure the roughness exponent of the system following Eq. (9) and the result is reported in Fig. 6. We obtain with both methods $`\zeta =1.26`$, which is consistent with the numerical value found in interface depinning , for a disordered chain , and with a recent two loop renormalization group calculation . This value nevertheless disagrees with previous results on a periodic chain .
The simulations of Ref. for periodic and disordered chain, suggest the presence of two different universality classes. In order to test this result we study the force fluctuations of a disordered chain. The equilibrium length of the springs connecting the beads is chosen randomly (Poisson). The chain is then driven at constant velocity and the distribution of the slip sizes is calculated. The result is shown in Fig. 7 where we also report the distribution obtained with a periodic chain using similar parameters. The two distributions are clearly indistinguishable, casting some doubt on the relevance of disorder in the spring lengths. From this study, one would conclude that the two dynamics are in the same universality class.
### B Constant force.
For constant force driving, we employ system sizes varying from $`L=20`$ to $`L=200`$ and the density of pinning sites was chosen equal to unity $`(L=N)`$. In order to determine the exponent $`\beta `$, we need an accurate estimate of the critical force, since an error in $`F_c`$ can strongly bias the fit.
Fig. 8 shows the value of the average velocity of the interface as a function of $`L`$. For $`F`$ smaller than $`F_c`$ in the limit of a large system $`v`$ tends to zero, and for $`F`$ greater than $`F_c`$ and for the same limit ($`L1`$) $`v`$ should tend a non vanishing value. In this way we can locate the critical force, which results to be $`F_c=2.195\pm 0.005`$. This result appears clearly from Fig. 8, from the log-log plot of $`v`$ as a function of $`L`$. We see that $`F=2.195`$ is compatible with a power-law behavior, whereas for $`F=2.200`$ in the large $`L`$ limit the mean velocity tends to a nonzero constant. For $`F=2.190`$ the velocity tends to zero faster than a power law in the limit of large system sizes. The numerical results are averaged over a number of disorder configurations which varies from $`4000`$ for $`L=20`$ to $`500`$ for $`L=200`$. The curve plotted in Fig. 8 allows also to estimate $`\beta /\nu =0.16\pm 0.02`$ (see Eq. (12)).
Next, we calculate the exponent $`\beta `$ directly, plotting $`v`$ vs $`(FF_c)`$. The fit in Fig. 9 yields $`\beta =0.22\pm 0.02`$. We restrict the fit to the the six smaller values since a crossover to linear behavior is expected at high forces (i.e. $`v=F/\eta `$) and this can bias the numerical estimate of the exponent. The simulations are made with a system of $`L=540`$ particles, and the results are averaged over $`100`$ configurations of the disorder. In this way, we can obtain $`\beta `$ and $`\nu `$ as summarized in table I.
To further test the consistency of our results, we calculate the exponent $`\zeta `$ measuring the scaling of $`W`$ with $`L`$ at $`F=F_c`$ (see Eq. 13) The results shown in Fig. 10 give $`\zeta =1.28\pm 0.03`$ in agreement with the result obtained in Sec. IV A.
## V Conclusion
In conclusion, we have investigated the dynamics of an elastic chain sliding on a disordered substrate. We have analyzed numerically the scaling close to the depinning transition focusing on the effect of different driving modes. Usually the problem is analyzed under constant force driving, while friction experiments are usually performed controlling the velocity. The two problems are closely related as discussed in Refs . We compute the critical exponents characterizing the transition and analyze the effect of the driving velocity and the loading spring stiffness. Our results are in qualitative agreement with friction experiments performed with macroscopic asperities coupled by an elastic matrix. In general friction experiments our model will not apply since in many instances the Larkin length is extremely large . In addition, inertial effect are present in most cases and could lead to different force fluctuations.
The exponents we measure agree well with the values expected for the depinning of elastic interfaces in quenched disordered media and with the renormalization group calculation of Ref. . This is not surprising, since it is possible to show using the method discussed in Ref. that the continuum limit of the model we study is described by
$$\frac{h(x,t)}{t}=Dh+F+f_p(x,h),$$
(14)
where $`h(x,t)`$ is a coarse grained version of $`u_i(t)`$ and $`f_p`$ is a coarse grained random pinning force. It is thus expected that the simulations performed in Ref. for a periodic chain agree with this result. The only difference between those simulations and ours lies in the way disorder is implemented: in Ref. the pinning point are arranged in a periodic structure and have random strength, while we use constant strength and random positions. We also note that the roughness exponent measured in Ref. is very close to $`\zeta =3/2`$, which is expected below the Larkin length, although the parameter employed do not seem to be consistent with that regime.
S. Z. acknowledges financial support from EC TMR Research Network under contract ERBFMRXCT960062.
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# Ground States in Non-relativistic Quantum Electrodynamics
## 1 INTRODUCTION
An established picture of an atom or molecule is that even in the presence of a quantized radiation field there is a ground state. The excited states that exist in the absence of coupling to the field are expected to melt into resonances, which means that they eventually decay with time into the ground state plus free photons. This picture has been established by Bach, Fröhlich and Sigal in for sufficiently small values of the various parameters that define the theory. Here we show that a ground state exists for all values of the parameters (including a variable $`g`$-factor) in the one particle case and, under a physically appropriate assumption, in the many-particle case.
We know that the Hamiltonian for the system is bounded below, so a ground state energy always exists in the sense of being the infimum of the spectrum, but the existence of a genuine normalizable solution to the eigenvalue equation is a more delicate matter that has received a great deal of attention, especially in recent years. A physically simple example where no ground state exists (as far as we believe now) is the free particle (i.e., particle plus field). In the presence of an external potential, however, like the Coulomb potential of a nucleus, a ground state should exist.
The difficulty in establishing this ground state comes from the fact that the bottom of the spectrum lies in the continuum (i.e., essential spectrum), not below it, as is the case for the usual Schrödinger equation. We denote the bottom of spectrum of the free-particle Hamiltonian for $`N`$ particles with appropriate statistics by $`E^0(N)`$. The ‘free-particle’ Hamiltonian includes the interparticle interaction (e.g., the Coulomb repulsion of electrons) but it does not include the interaction with a fixed external potential, e.g., the interaction with nuclei. When the latter is included we denote the bottom of the spectrum by $`E^V(N)`$. It is not hard to see in many cases that $`E^V(N)<E^0(N)`$, but despite this inequality $`E^V(N)`$ is, nevertheless, the bottom of the essential spectrum. The reason is that we can always add arbitrarily many, arbitrarily ‘soft’ photons that add arbitrarily little energy. It is the soft photon problem that is our primary concern here.
The main point of this paper is to show how to overcome this infrared problem and to show, quite generally for a one-particle system, that a ground state exists for all values of the particle mass, the coupling to the field (fine-structure constant $`\alpha =e^2/\mathrm{}c`$), the magnetic $`g`$-factor, and the ultraviolet cutoff $`\mathrm{\Lambda }`$ of the electromagnetic field frequencies, provided a bound state exists when the field is turned off. This result implies, in particular, that for a fixed ultraviolet cutoff renormalization of the various physical quantities will not affect the existence of a ground state. Of course, nothing can be said about the limit as the cutoff tends to infinity. We also include a large class of interactions much more general than the usual Coulomb interaction.
The model we discuss has been used quite frequently in field theory. In its classical version it was investigated by Kramers who seems to have been the first to point out the possibility of renormalization. The quantized version was investigated by Pauli and Fierz in connection with scattering theory. Most importantly, it was used by Bethe to obtain a suprisingly good value for the Lamb shift.
Various restricted versions of the problem have been attacked successfully. In the early seventies Fröhlich investigated the infrared problem in translation invariant models of scalar electrons coupled to scalar bosons . In he proved that for an electron coupled to a massive field a unique ground state exists for fixed total momentum. External potentials were not considered in these papers.
The first rigorous result on the bound state problem, to our knowledge, is due to Arai and concerns one particle confined by an $`x^2`$-potential, the interaction with the photons being subject to the dipole-approximation. For this model, which is explicitly solvable, Arai proved existence and uniqueness of the ground state . Later, Spohn showed by perturbation theory that this result extends to perturbed $`x^2`$\- potentials . No bounds on the parameters were needed to obtain these results but the methods oviously do not admit extensions to more realistic models.
Bach , Fröhlich and Sigal initiated the study of the full nonrelativistic QED model (the same model as considered in the present paper) under various simplifying assumptions. In the existence of a ground state in this model for particles subject to an external binding potential was proved for $`\alpha \mathrm{\Lambda }`$ small enough. The main achievement of this paper is the elimination of the earlier simplifying assumptions, especially the infrared regularization. This is the first, and up to now the only paper where a ‘first principles’ QED model was successfully analyzed, but with a restricted parameter range. A weaker result for the same system but with simplifications such as infrared regularization were independently obtained by Hiroshima by entirely different methods; he also showed uniqueness of the ground state, assuming its existence .
In a parallel development Arai and Hirokawa , Spohn , and Gerard investigated the ground-state problem for systems similar to the one of Bach et al in various degrees of mathematical generality. Arai and Hirokawa proved existence in what they called a ‘generalized spin-boson model’. If specialized to the case of a non-relativistic $`N`$-particle system interacting with bosons, their result proves the existence of a ground state for confined particles and small $`\alpha `$. This result was extended in a recent preprint to account for non-confined particles and systems with the true infrared singularity of QED. Concrete results in the infrared-singular case concern only special models, however, such as the Wigner-Weisskopf Hamiltonian which describes a two-level atom. Hirokawa continued this work in a recent preprint . Spohn and Gerard both proved existence of a ground state for a confined particle and arbitrary coupling constant, the result of Gerard being somewhat more general .
For the existence of a ground state in the case of massive bosons, which is a typical intermediate result in the cited works, a short and elegant proof was given by Derezinski and Gerard for the case of linear coupling, in which the $`A^2`$ term is omitted, and a confining potential. Some of their ideas are used in our paper.
The Hamiltonian for $`N`$-particles has four parts which are described precisely in the next section,
$$H^V=T+V+I+H_f.$$
(1)
The dependence of $`H^V`$ on $`N`$ is not noted explicitly. The first term, $`T=_{j=1}^NT_j`$, is the kinetic energy with ‘minimal coupling’ in the Coulomb gauge (i.e., $`p`$ is replaced by $`p+\sqrt{\alpha }A`$, where $`A`$ is the magnetic vector potential satisfying $`\mathrm{div}A=0`$), $`V(X)`$, with $`X=(x_1,x_2,\mathrm{},x_N)`$ is the external potential, typically a Coulomb attraction to one or more nuclei. In any case, we assume that $`V`$ is a sum of one-body potentials, i.e.,
$$V(X)=\underset{j=1}{\overset{N}{}}v(x_j).$$
The particle-particle interaction, $`I`$, has the important feature that it is translation invariant and, of course, symmetric in the particle labels. Both $`I`$ and $`V`$ could be spin dependent, but we shall not burden the notation with this latter possibility. Typically $`I`$ is a Coulomb repulsion, but we do not have to assume that $`I`$ is merely a sum of two-body potentials. The only requirements are: (1) the negative part of $`v`$ vanishes at infinity; (2) the negative part of $`I`$ satisfies clustering, i.e., the intercluster part of $`I_{}`$ tends to 0 when the spacing between any two clusters tends to infinity; (3) $`V_{}+I_{}`$ are dominated by the kinetic energy as in (5). Another assumption we make (in the $`N`$-particle case) is that there is binding, as described below.
The natural choice for the kinetic energy is the ‘Pauli operator’ $`T=(p+\sqrt{\alpha }A(x))^2+\sqrt{\alpha }\sigma B`$, but we can generalize this to include the case of the usual kinetic energy $`T=(p+\sqrt{\alpha }A(x))^2`$ by introducing a ‘g-factor’, $`g`$. Thus, we take
$$T=(p+\sqrt{\alpha }A(x))^2+\frac{g}{2}\sqrt{\alpha }\sigma B.$$
(2)
Note that $`T`$ is a positive operator if $`0g2`$ and may not be otherwise. Nevertheless, the Hamiltonian is always bounded below because of the ultraviolet cutoff we shall impose on the $`A`$ field, which implies that $`(g/2)\sqrt{\alpha }\sigma B+H_f`$ is always bounded below.
We believe that the ‘relativistic’ operator $`T=|p+\sqrt{\alpha }A(x)|`$, presents no real difficulty either, but we do not want to overburden this paper with a lengthy proof. This problem is currently under investigation.
A model that is frequently discussed is the ‘Pauli-Fierz’ model, but it is not entirely clear how this is defined since several variants appear in . One version uses $`T=(p+\sqrt{\alpha }A(x))^2+\sqrt{\alpha }\sigma B`$, which is one of the models under consideration. Another variant uses a linearized version of this operator, $`T=p^2+2\sqrt{\alpha }pA(x)+\sqrt{\alpha }\sigma B`$ or $`T=p^2+2\sqrt{\alpha }pA(x)`$. These variants are not gauge invariant and, therefore, depend on the choice of gauge for $`A`$. Our method is applicable to these linearized models in some gauges, but not in others. We omit further discussion of this point since these variants are not the most relevant ones for quantum electrodynamics.
There is one important point that as far as we know, has not been mentioned in the QED context. This is the binding condition. Our proof of the existence of a ground state uses, as input, the assumption that
$$E^V(N)<E^V(N^{})+E^0(NN^{})\mathrm{for}\mathrm{all}N^{}<N.$$
(3)
Our work can be generalized (but we shall not do so here) to the case in which the external potential is that of attractive nuclei and these nuclei are also dynamical particles. Then, of course, it is necessary to work in the center of mass system and then (3) must be replaced by the condition that $`E^V(N)`$ is less than the lowest two-cluster threshold. While this condition, or (3) in the static case, are physically necessary for the existence of a ground state, the validity of these conditions cannot be taken for granted.
We prove inequality (3) for one particle ($`N=1`$) quite generally, using only the assumption that the ordinary Schrödinger operator $`p^2+v`$ has a negative energy ground state. This certainly holds for the Coulomb potential. Indeed, one could expect, on physical grounds, that there could be binding even if $`p^2+v`$ is has no negative energy bound state, because the interaction with the field increases the effective mass of the particle — and hence the binding energy. The same argument shows that when there are $`N`$ particles at least one them is necessarily bound, i.e., $`E^V(N)<E^0(N)`$. When we consider more than one particle, we are not able to show (3) for all $`N^{}>1`$, even if $`p_j^2+I+V`$ has a ground state.
In the Coulomb case, it is possible to show (but we shall not do so here) that condition (3) is satisfied if the nuclear charge $`Z`$ is large enough. The basic idea is that if breakup into two groups, one of them with $`N^{}`$ particles close to the nucleus and the second consisting of $`NN^{}`$ particles far away occurs, then there will be an attractive Coulomb tail acting on the separated particles at a distance $`R`$ away with net attractive potential $`(ZN^{})/R`$. However, to localize one of these particles within a distance $`R`$ of the nucleus will require a field energy localization error of the order of $`C/R`$, by dimensional analysis arguments. If $`(ZN^{})>C`$ then the energy can be decreased by bringing one of the unbound particles close to the nucleus.
Section 2 introduces the precise definition of our problem and the main result Theorem 2.1.
In Section 3, we show how to prove that $`E^V(1)<E^0(1)`$. More generally, $`E^V(N)<E^0(N)`$ if $`v0`$.
Our strategy to establish a ground state is the usual one of showing that a minimizing sequence of trial vectors for the energy actually has a weak limit that, in fact, is a minimizer. The problem here is that one can easily construct minimizing sequences that converge weakly to zero by choosing vectors with too many soft photons. To avoid this we take a special sequence.
To define this sequence we first consider an artificial model in which the photons have a mass, i.e., $`\omega (k)=\sqrt{k^2+m^2}`$. Here there is no soft photon problem and we show in Section 4 that this model has a ground state $`\mathrm{\Phi }_m`$.
In Section 5 we show that as $`m0`$ the $`\mathrm{\Phi }_m`$ sequence is minimizing. Then in Section 6 we use the Schrödinger equation for $`\mathrm{\Phi }_m`$ to deduce certain properties of $`\mathrm{\Phi }_m`$ which we call infrared bounds. One of these was proved in but we need one more, which is new.
With these bounds we can show in Section 7 that $`\mathrm{\Phi }_m`$ has a strong limit as $`m0`$, which is a minimizer for $`H^V`$.
Acknowledgment: We thank Professor Fumio Hiroshima for a useful correspondence concerning equation (53) and for sending us his preprint .
## 2 DEFINITIONS AND MAIN THEOREM
The Hamiltonian for N particles interacting with the quantized radiation field and with a given external potential $`V(X)`$, with $`X=(x_1,x_2,\mathrm{},x_N)`$ and $`x_j^3`$, is
$$H^V=\underset{j=1}{\overset{N}{}}\left\{\left(p_j+\sqrt{\alpha }A(x_j)\right)^2+\frac{g}{2}\sqrt{\alpha }\sigma _jB(x_j)\right\}+V(X)+I(X)+H_f.$$
(4)
The unit of energy is $`Mc^2`$/2, where $`M`$ is the particle mass, the unit of length is $`\mathrm{}_c=2\mathrm{}/Mc`$, twice the Compton wavelength of the particle, and $`\alpha =e^2/\mathrm{}c`$ is the dimensionless “fine structure constant” (=1/137 in nature). The electric charge of the particle is $`e`$. The unit of time is the time it takes a light wave to travel a Compton wavelength, i. e., the speed of light is $`c=1`$.
The operator $`p=i`$, while $`A`$ is the (ultraviolet cutoff) magnetic vector potential (we use the Coulomb, or radiation gauge). The unit of $`A^2(x)`$ is $`Mc^2/2\mathrm{}_c`$. The magnetic field is $`B`$ = curl$`A`$ and the unit of B is $`\alpha ^{3/2}`$ times the quantity $`M^2e^3c/4\mathrm{}^3`$, which is the value of $`B`$ for which the magnetic length ($`\mathrm{}c/eB)^{1/2}`$ equals the Bohr radius $`2\mathrm{}^2/Me^2`$.
The reader might wonder why we use these units, which seem to be more appropriate for a relativistic theory than for the nonrelativistic theory we are considering. Why not use the Bohr radius as the unit of length, for example? Our reason is that we want to isolate the electric charge, which is the quantity that defines the interaction of matter with the electromagnetic field, in precisely one place, namely $`\alpha `$. ‘Atomic units’ have the charge built into the length, etc. and we find this difficult to disentangle.
The Hilbert space is an appropriate subspace of
$$=^NL^2(^3;^2),$$
where $``$ is the Fock space for the photon field. We have in mind Fermi statistics (the antisymmetric subspace of $`^NL^2(^3;^2)`$) and the $`^2`$ is to accomodate the electron spin. We can also deal with ”Boltzmann statistics”,in which case we would set $`g=0`$ and use $`^NL^2(^3)`$, or bose statistics, in which case we would set $`g=0`$ and use the symmetric subspace of $`^NL^2(^3)`$. These generalizations are mathematically trivial and we do not discuss them further.
For our purposes we assume that for every $`\epsilon >0`$ there exists a constant $`a(\epsilon )`$ such that the negative part of the potentials, $`V_{}(X)`$ and $`I_{}(X)`$, satisfy
$$V_{}+I_{}\epsilon \underset{j=1}{\overset{N}{}}p_j^2+a(\epsilon )$$
(5)
as quadratic forms on $``$.
The vector potential is
$$A(x)=\underset{\lambda =1,2}{}_{|k|<\mathrm{\Lambda }}\frac{1}{\sqrt{|k|}}\left[\epsilon _\lambda (k)a_\lambda ^{}(k)+\epsilon _\lambda (k)a_\lambda ^{}(k)\right]e^{ikx}d^3k$$
(6)
where the operators $`a_\lambda ^{},a_\lambda ^{}`$ satisfy the usual commutation relations
$$[a_\lambda ^{}(k),a_\nu ^{}(q)]=\delta (kq)\delta _{\lambda ,\nu },[a_\lambda ^{}(k),a_\nu (q)]=0,\mathrm{etc}$$
(7)
and the vectors $`\epsilon _\lambda (k)`$ are the two possible orthonormal polarization vectors perpendicular to $`k`$. They are chosen for convenience in (59,60).
The number $`\mathrm{\Lambda }`$ is the ultraviolet cutoff on the wavenumbers $`k`$. Our results hold for all finite $`\mathrm{\Lambda }`$. The details of the cutoff in (6) are quite unimportant, except for the requirement that rotation symmetry in $`k`$-space is maintained. E.g., a gaussian cutoff can be used instead of our sharp cutoff. We avoid unnecessary generalisations.
The field energy $`H_f`$, sometimes called $`d\mathrm{\Gamma }(\omega )`$ is given by
$$H_f=\underset{\lambda =1,2}{}_^3\omega (k)a_\lambda ^{}(k)a_\lambda ^{}(k)d^3k$$
(8)
The energy of a photon is $`\omega (k)`$ and the physical value of interest to us is
$$\omega (k)=|k|,$$
(9)
in our units. Indeed, any continuous function that is bounded below by $`\mathrm{const}.|k|`$ for small $`|k|`$ is acceptable. In the process of proving the existence of a ground state for $`H`$ we will first study the unphysical ‘massive photon’ case, in which
$$\omega _m(k)=\sqrt{k^2+m^2}$$
(10)
for some $`m>0,`$ called the ‘photon mass’.
In the remainder of this paper, unless otherwise stated, we shall always assume that there is no restriction on $`\alpha ,\mathrm{\Lambda }`$ and $`g`$ and that $`\omega (k)`$ can be either as in (9) or as in (10).
By Lemma A.5 we see easily that $`H^V`$ is bounded below for all values of the parameters, including $`m=0`$. Thus, $`H^V`$ defines a closable quadratic form and hence a selfadjoint operator, the Friedrich’s extension. We denote this extension again by $`H^V`$.
Our main theorem is
###### 2.1 THEOREM (Existence of a ground state).
Assume that the binding condition (3) and the condition (5) hold. Then there is a vector $`\mathrm{\Phi }`$ in the $`N`$-particle Hilbert space $``$ such that
$$H^V\mathrm{\Phi }=E^V(N)\mathrm{\Phi }.$$
(11)
## 3 UPPER BOUND
We shall prove the binding condition (3) for one particle and, with an additional assumption, for the $`N`$ particle case as well. This is that if $`N`$-particles are present then at least one of them binds.
As we mentioned before, at least in the single particle case, that our requirement that the system without the radiation field has a bound state is somewhat unnaturally restrictive, since one expects that the radiation field enhances binding; this has been shown to be true in the ‘dipole’, or Kramers approximation . We are able to show in the one-particle case, that the photon field cannot decrease the binding energy. It is quite possible that there could be binding even when the operator $`p^2+v`$ does not have a negative energy state, but we cannot shed any light on that question.
For the one-particle case the situation is less delicate than the $`N`$-particle case.
###### 3.1 THEOREM (Binding of at least one particle).
Assume that the one-particle Hamiltonian $`p^2+v(x)`$ has a negative energy bound state with eigenfunction $`\varphi (x)`$ and energy $`e_0`$. Then,
$$E^V(1)E^0(1)e_0,$$
(12)
i.e., binding continues to exist when the field is turned on.
For the $`N`$-particle case we make the additional assumption that $`v(x)0`$ for all $`x`$. Then,
$$E^V(N)E^0(N)e_0,$$
(13)
i.e., at least one particle is bound.
###### Proof.
It suffices to prove that $`E^V(N)E^0(N)+\epsilon e_0`$ for all $`\epsilon >0`$. There is a normalized vector $`F`$ such that $`(F,H^0F)<E^0(N)+\epsilon `$. ($`F`$ is antisymmetric according to the Pauli principle.) We use the notation $`,`$ to denote the inner product in Fock space and spin space. Then we can write $`(F,H^0F)=G(X)d^{3N}X`$ with
$`G(X)=`$
$`{\displaystyle \underset{j=1}{\overset{N}{}}}\left\{(i_j+\sqrt{\alpha }A(x_j))F,(i_j+\sqrt{\alpha }A(x_j))F(X)+\sqrt{\alpha }(g/2)F,\sigma _jB(x_j)F(X)\right\}`$
$`+F,(I+H_f)F(X).`$ (14)
As a (unnormalized) variational trial vector we take the vector $`\psi =\left[_{j=1}^N\varphi (x_j)^2\right]^{1/2}F`$. Recall that $`\varphi (x)0`$ since $`\varphi `$ is the ground state of $`p^2+v`$. We also recall the Schwarz inequality
$$\left|\frac{_{j=1}^N\varphi (x_j)\varphi (x_j)}{\left[_{j=1}^N\varphi (x_j)^2\right]^{1/2}}\right|^2\underset{j=1}{\overset{N}{}}\left|\varphi (x_j)\right|^2.$$
(15)
Using (15), integration by parts, and the fact that $`\varphi `$ satisfies the Schrödinger equation $`(p^2+v)\varphi =e_0\varphi `$, we easily find that
$`(\psi ,\left[H^V(E^0(N)+\epsilon e_0)\right]\psi )`$
$`{\displaystyle \left\{G(X)(E^0(N)+\epsilon )F,F(X)\right\}\underset{j=1}{\overset{N}{}}\varphi (x_j)^2d^{3N}X}`$
$`+{\displaystyle \underset{jk}{}v(x_k)\varphi (x_j)^2F,F(X)d^{3N}X}.`$ (16)
When $`N=1`$ the last term in (3) is not present so no assumption about the potential $`v`$ is needed. When $`N>1`$ we can omit the last term because it is negative by assumption.
Now, by the $`^3`$-translation invariance of $`H_0`$, for every $`y^3`$ there is a ‘translated’ vector $`F_y`$ so that $`G(X)G(X+(y,\mathrm{},y))`$ and $`F_y,F_y(X)=F,F(X+(y,\mathrm{},y))`$. (This is accomplished by the unitary operator on $``$ that takes $`x_jx_j+y`$ for every $`j`$ and $`a_\lambda ^{}(k)exp(iky)a_\lambda ^{}(k)`$.) Thus, if we denote the quantity in $`\{\}`$ in (3) by $`W(X)`$, and if we define $`\psi _y`$ by replacing $`F`$ by $`F_y`$ in the definition of $`\psi `$, we have
$`\mathrm{\Omega }(y)=(\psi _y,H^V(E_0+\epsilon e_0)\psi _y)`$
$`{\displaystyle W(X+(y,\mathrm{},y))\underset{j=1}{\overset{N}{}}\varphi (x_j)^2d^{3N}X}={\displaystyle W(X)\underset{j=1}{\overset{N}{}}\varphi (x_jy)^2d^{3N}X}.`$ (17)
Note that $`\mathrm{\Omega }(y)𝑑yNW(X)d^{3N}X`$. But $`W(X)d^{3N}X=(F,(H^0E^0(N)\epsilon )F)`$ and this is strictly negative by assumption. Hence, for some $`y^3`$ we have that $`\mathrm{\Omega }(y)<0`$ and thus $`\psi _y0`$, which proves the theorem. ∎
Remark: \[Alternative theorem\]
It may be useful to note, briefly, a different proof of Theorem 3.1, for long range potentials $`v(x)`$, such as the attractive Coulomb potential $`Z/|x|`$, which shows that the bottom of the spectrum of $`H^V`$ lies strictly below $`E^0`$. Unfortunately, this proof does not show that the difference is at least $`e_0`$. We sketch it for the one-body case. Using the notation of the proof above, the first step is to replace $`F`$ by $`F_R=u(x_1/R)F`$ where $`u`$ is a smooth function with support in a ball of radius $`1`$. One easily finds that $`(F_R,H_0F_R)/(F_R,F_R)=E^0+\epsilon +c/R^2`$, where $`c`$ is a constant that depends only on $`u`$ and not on $`\epsilon `$ and $`R`$. On the other hand $`(F_R,VF_R)/(F_R,F_R)Z/R`$, to use the Coulomb potential as an example. To complete the argument, choose $`R=2c/Z`$ and then choose $`\epsilon =c/R^2`$. What we have used here is the fact that localization ‘costs’ a kinetic energy $`R^2`$, while the potential energy falls off slower than this, e.g., $`R^1`$.
## 4 GROUND STATE WITH MASSIVE PHOTONS
As we emphasized in the introduction, not every minimizing sequence converges to the minimizer for our $`m=0`$ problem, i.e., with $`\omega (k)=|k|`$. The situation is much easier for the massive case (10). The Hamiltonian in this case is given by (4) and $`H_f`$ is given by (8) with (10). To emphasize the dependence on $`m`$ we denote this Hamiltonian and field energy by $`H_m^V`$ and $`H_f(m)`$, respectively. Likewise, $`E^V(m,N)`$ and $`E^0(m,N)`$ denote the mass dependent energies, as defined before.
We emphasize that the vector potential is still given by (6), but we could, if we wished, easily replace $`|k|^{1/2}`$ in (6) by $`(k^2+m^2)^{1/4}`$.
It will be shown in this section that $`H_m^V`$ has a ground state. More precisely we prove
###### 4.1 THEOREM (Existence of ground state).
Assume that for some fixed value of the ultraviolet cutoff $`\mathrm{\Lambda }`$ there is binding for the Hamiltonian $`H_m^V`$, i.e., $`E^V(m,N)<\mathrm{\Sigma }^V(m,N)`$ where $`\mathrm{\Sigma }^V(m,N)=min\{E^V(m,N^{})+E^0(m,NN^{}):\mathrm{all}N^{}<N\}`$ is the ‘lowest two-cluster threshold’. Then $`E(m,N)`$ is an eigenvalue, i.e., there exists a state $`\mathrm{\Phi }_m`$ in $``$ such that $`H_m^V\mathrm{\Phi }_m=E(m,N)\mathrm{\Phi }_m`$.
###### Proof.
Let us first show that it suffices to prove that for any normalized sequence $`\mathrm{\Psi }^j`$, $`j=1,2,\mathrm{}`$, (not necessarily minimizing) tending weakly to zero
$$\underset{j\mathrm{}}{lim\; inf}(\mathrm{\Psi }^j,H_m^V\mathrm{\Psi }^j)>E^V(m,N).$$
(18)
To prove this let $`\mathrm{\Phi }^j`$ be some minimizing sequence, i.e., assume that
$$\mathrm{\Phi }^j=1,$$
(19)
and that
$$(\mathrm{\Phi }^j,H_m^V\mathrm{\Phi }^j)E^V(m,N).$$
(20)
By the Banach Alaoglu Theorem we can assume that this sequence, as well as the sequence $`H_m^V\mathrm{\Phi }^j`$ converge weakly in the sense that for any $`\mathrm{\Psi }`$ with $`(\mathrm{\Psi },H_m^V\mathrm{\Psi })<\mathrm{}`$ we have that
$$(\mathrm{\Psi },H_m^V\mathrm{\Phi }^j)(\mathrm{\Psi },H_m^V\mathrm{\Phi }_m),$$
(21)
where $`\mathrm{\Phi }_m`$ is the weak limit of $`\mathrm{\Phi }^j`$. Our goal is to show that $`(\mathrm{\Phi }_m,H_m^V\mathrm{\Phi }_m)=E^V(m,N)`$ and that $`\mathrm{\Phi }_m=1`$.
Write $`\mathrm{\Phi }^j=\mathrm{\Phi }_m+\mathrm{\Psi }^j`$. Obviously $`\mathrm{\Psi }^j`$ as well as $`H_m^V\mathrm{\Psi }^j`$ go weakly to zero. Thus
$`0`$ $`=\underset{j\mathrm{}}{lim}(\mathrm{\Phi }^j,(H_m^VE^V(m,N))\mathrm{\Phi }^j)`$
$`=\underset{j\mathrm{}}{lim}((\mathrm{\Phi }_m+\mathrm{\Psi }^j),(H_m^VE^V(m,N))(\mathrm{\Phi }_m+\mathrm{\Psi }^j))`$
$`=\underset{j\mathrm{}}{lim}(\mathrm{\Psi }^j,(H_m^VE^V(m,N))\mathrm{\Psi }^j)+(\mathrm{\Phi }_m,(H_m^VE^V(m,N))\mathrm{\Phi }_m)`$
where we used that the cross terms vanish. Since $`H_m^VE^V(m,N)0`$ this shows that $`\mathrm{\Phi }_m`$ minimizes the energy, and, furthermore, that
$$0\underset{j\mathrm{}}{lim}(\mathrm{\Psi }^j,(H_m^VE^V(m,N))\mathrm{\Psi }^j)\delta \underset{j\mathrm{}}{lim\; inf}\mathrm{\Psi }^j^2$$
for some positive constant $`\delta `$. The second inequality is trivial if $`lim\; inf_j\mathrm{}\mathrm{\Psi }^j^2=0`$ and otherwise follows from our assumption (18). This proves that $`\mathrm{\Psi }^j`$ converges strongly to zero along a subsequence, which implies that $`\mathrm{\Phi }_m=1`$. Hence $`\mathrm{\Phi }_m`$ is a normalized ground state. Thus, it suffices to prove (18).
The steps that lead to a proof of (18) are quite standard. The only difficulty is that one has to localize in Fock space, which we describe first. We follow with some necessary modifications and some simplifications.
Recall that, when the $`a_\lambda ^\mathrm{\#}`$ operators are viewed in $`x`$-space
$$a_\lambda ^{}(f):,a_\lambda ^{}(g):,$$
(22)
they obey the commutation relations
$$[a_\lambda ^{}(f),a_\lambda ^{}(g)]=_^3\overline{f}(x)g(x)d^3x=:(f,g).$$
(23)
Consider now two smooth localization functions $`j_1`$ and $`j_2`$ that satisfy $`j_1^2+j_2^2=1`$ and $`j_1`$ is supported in a ball of radius $`P`$. The first derivatives of $`j_1`$ and $`j_2`$ are of order $`1/P`$.
The operators
$$c_\lambda (f)=a_\lambda ^{}(j_1f)+a_\lambda ^{}(j_2f),c_\lambda ^{}(g)=a_\lambda ^{}(j_1g)+a_\lambda ^{}(j_2g)$$
(24)
act both on the space $``$. Note that
$$[c_\lambda (f),c_\lambda ^{}(g)]=(f,g).$$
(25)
Thus, these new creation and anihilation operators create another Fock space $`^l`$ that is a subspace of $``$ and is isomorphic to the old Fock space $``$. Hence, there exists a map
$$U:^l$$
(26)
that is an invertible isometry between Fock spaces. It is uniquely specified by the properties
$$a^\mathrm{\#}=U^{}c^\mathrm{\#}U,$$
(27)
and the vacuum in $``$ is mapped to the vacuum in $``$.
The map $`U^{}`$ is defined on $`^l`$ only, but we can extend it to all of $``$ by setting $`U^{}F=0`$ whenever $`F`$ is perpendicular to $`^l`$. In other words $`U^{}`$ is a partial isometry between Fockspaces where $`U^{}U=`$ on $``$, and where $`UU^{}`$ is the orthogonal projection onto $`^l`$. We continue to denote the extended map by $`U^{}`$.
Let $`\varphi `$ and $`\overline{\varphi }`$ be smooth nonnegative functions, with $`\varphi ^2+\overline{\varphi }^2=1`$, $`\varphi `$ identically one on the unit ball, and vanishing outside the ball of radius $`2`$. Set $`\varphi _R(X)=\varphi (X/R)`$. It is a standard calculation to show that for any $`\mathrm{\Psi }`$ with finite energy
$$(\mathrm{\Psi },H_m^V\mathrm{\Psi })=(\varphi _R\mathrm{\Psi },H_m^V\varphi _R\mathrm{\Psi })+(\overline{\varphi _R}\mathrm{\Psi },H_m^V\overline{\varphi _R}\mathrm{\Psi })(\mathrm{\Psi },(\varphi _R)^2\mathrm{\Psi })(\mathrm{\Psi },(\overline{\varphi _R})^2\mathrm{\Psi }).$$
(28)
The last two terms in (28) are bounded by $`const./R^2`$.
One goal will be to show that for any $`\mathrm{\Psi }`$ with finite energy
$`(\mathrm{\Psi },\varphi _RH_m^V\varphi _R\mathrm{\Psi })=`$
$`(\mathrm{\Psi },\varphi _RU^{}\left\{H_m^V+H_f\right\}U\varphi _R\mathrm{\Psi })+o(1).`$ (29)
The error term $`o(1)`$ vanishes as both $`R`$ and $`P`$ go to infinity and depends otherwise only on the energy of $`\mathrm{\Psi }`$. Notice that the invertible map $`U`$ depends on the cutoff parameter $`P`$ as well. (29) will be proved in Lemma A.1 in the Appendix. The intuition behind the estimate (29) is that localized electrons interact only weakly with far away photons. Those photons are described solely by their own field energy.
An immediate consequence of (29) is the estimate
$$(\mathrm{\Psi },\varphi _RH_m^V\varphi _R\mathrm{\Psi })(E^V(m,N)+m)\varphi _R\mathrm{\Psi }^2m(\varphi _R\mathrm{\Psi },U^{}P_2U\varphi _R\mathrm{\Psi })+o(1),$$
(30)
which is obtained by noting that the field energy in the second factor can be estimated from below by
$$H_fmmP_2,$$
(31)
where $`P_2`$ is the projection onto the vacuum of the second factor of $``$.
In a further step we prove in Lemma A.3 that the sequence
$$(\varphi _R\mathrm{\Psi }^j,U^{}P_2U\varphi _R\mathrm{\Psi }^j)0$$
(32)
as $`j\mathrm{}`$.
Returning to (28), using Corollary A.2 we have that
$$(\overline{\varphi _R}\mathrm{\Psi },H_m^V\overline{\varphi _R}\mathrm{\Psi })\mathrm{\Sigma }^V(m,N)\overline{\varphi _R}\mathrm{\Psi }^2o(1),$$
(33)
with $`o(1)`$ going to zero as $`R\mathrm{}`$. Roughly speaking $`\overline{\varphi }`$ forces some of the particles to be far away from the origin. Any such particle configuration can be described by two clusters with no interaction between them. In particular the interaction between these clusters via the radiation field is turned off. This means that each cluster carries its own field energy. To prove this the localization in Fock space is used. Moreover the cluster that is far away from the origin does not interact with the external potential although the repulsion among its particles is still present.
To summarize, by combining (30), (32) and (33) we have proved that
$$lim\; inf(\mathrm{\Psi }^j,H_m^V\mathrm{\Psi }^j)(E^V(m,N)+\delta )+o(1)$$
(34)
where
$$\delta =\mathrm{min}\{m,\mathrm{\Sigma }^V(m,N)E^V(m,N)\},$$
(35)
and $`o(1)`$ tends to zero as $`R\mathrm{}`$ and $`P\mathrm{}`$. ∎
## 5 A MINIMIZING SEQUENCE
We consider the Hamiltonian $`H_m^V`$ defined in (4) with field energy $`H_f(m)`$ defined using $`\omega _m(k)=\sqrt{k^2+m^2}`$. Our main goal here is Theorem 5.3, which shows that the ground states of the $`m>0`$ problem form a minimizing sequence for the $`m=0`$ problem.
###### 5.1 THEOREM ($`E^V(m,N)`$ converges to $`E^V(0,N)`$).
As $`m0`$,
$$E^V(m,N)E^V(0,N)\mathrm{and}E^0(m,N)E^0(0,N)$$
(36)
###### Proof.
First, note that $`H_m^V>H_m^{}^V>H_0^V`$ if $`m>m^{}>0`$, because $`\omega _m`$ has this same monotonicity property. Therefore, for any sequence of $`m0`$, $`E^V(m,N)`$ is monotonically decreasing and has a sequence-independent, finite limit, which we call $`E^{}`$, and we note that $`E^{}E^V(0,N)`$. To prove the opposite, namely $`E^{}E^V(0,N)`$, we shall prove that $`E^{}E^V(0,N)+2\epsilon `$ for every $`\epsilon >0`$.
Let $`\mathrm{\Phi }`$ be normalized and such that $`(\mathrm{\Phi },H_0^V\mathrm{\Phi })<E^V(0,N)+\epsilon `$.
We note that $`H_m^V<H_0^V+m𝒩`$, where $`𝒩`$ is the number operator
$$𝒩=\underset{\lambda =1,2}{}_{|k|<\mathrm{\Lambda }}a_\lambda ^{}(k)a_\lambda ^{}(k)d^3k.$$
(37)
Thus, if we use $`\mathrm{\Phi }`$ as a variational function for $`H_m^V`$ we have $`E^V(m,N)<E^V(0,N)+\epsilon +m(\mathrm{\Phi },𝒩\mathrm{\Phi })`$, and our goal is accomplished provided that $`\mathrm{\Phi }`$ can be chosen so that $`(\mathrm{\Phi },𝒩\mathrm{\Phi })<\mathrm{}`$, in addition to $`(\mathrm{\Phi },H_0^V\mathrm{\Phi })<E^V(0,N)+\epsilon `$. If a way can be found to modify $`\mathrm{\Phi }`$ to another vector $`\stackrel{~}{\mathrm{\Phi }}`$ so that $`(\stackrel{~}{\mathrm{\Phi }},𝒩\stackrel{~}{\mathrm{\Phi }})<\mathrm{}`$, in addition to $`(\stackrel{~}{\mathrm{\Phi }},H_0^V\stackrel{~}{\mathrm{\Phi }})<E^V(0,N)+2\epsilon `$ the proof will be complete.
A suitable choice is $`\stackrel{~}{\mathrm{\Phi }}=\mathrm{\Pi }_n\mathrm{\Phi }`$, with $`\mathrm{\Pi }_n`$ being the projector onto the subspace of $``$, with $`n`$ or fewer photons, i.e., $`_n=^NL^2(^3;^2)_n`$, for an appropriately large $`n`$. It is easy to see, with the help of Lemma A.5 that $`(\mathrm{\Pi }_n\mathrm{\Phi },[V+I]\mathrm{\Pi }_n\mathrm{\Phi })(\mathrm{\Phi },[V+I]\mathrm{\Phi })`$, that $`(\mathrm{\Pi }_n\mathrm{\Phi },H_f\mathrm{\Pi }_n\mathrm{\Phi })(\mathrm{\Phi },H_f\mathrm{\Phi })`$, and $`(\mathrm{\Pi }_n\mathrm{\Phi },p_j^2\mathrm{\Pi }_n\mathrm{\Phi })(\mathrm{\Phi },p_j^2\mathrm{\Phi })`$ as $`n\mathrm{}`$. The following Lemma 5.2, shows that the other terms converge as well. The same proof works for $`E^0(m,N)`$. ∎
###### 5.2 LEMMA (Finite photon number approximation).
Let $`\mathrm{\Phi }`$ be such that $`(\mathrm{\Phi },H_f\mathrm{\Phi })<\mathrm{}`$. Denote by $`\mathrm{\Pi }_n`$ the projection in $``$ onto states with photon number less than or equal to $`n`$, i.e., $`\mathrm{\Pi }_n`$ is the projection onto the subspace $`_n=^NL^2(^3;^2)_n`$. Let $`\mathrm{\Phi }_n=\mathrm{\Pi }_n\mathrm{\Phi }`$. Then we have the following strong convergence as $`n\mathrm{}`$ (in addition to $`\mathrm{\Phi }_n\mathrm{\Phi }`$ )
$`A(x_j)\mathrm{\Phi }_n`$ $``$ $`A(x_j)\mathrm{\Phi }`$
$`B(x_j)\mathrm{\Phi }_n`$ $``$ $`B(x_j)\mathrm{\Phi }.`$ (38)
###### Proof.
Write $`A=D+D^{}`$ where $`D`$ contains the annihilation operators and $`D^{}`$ contains the creation operators. We omit $`(x_j)`$ and we omit the vector index of $`A`$ for simplicity. Since $`(\mathrm{\Phi },H_f\mathrm{\Phi })<\mathrm{}`$ we learn from Lemma A.4 that $`D\mathrm{\Phi }`$ and $`D^{}\mathrm{\Phi }`$ are in $``$. Since $`\mathrm{\Pi }_{n1}`$ strongly,
$$D\mathrm{\Phi }_n=\mathrm{\Pi }_{n1}D\mathrm{\Phi }D\mathrm{\Phi }.$$
The same holds if $`D`$ is replaced by $`D^{}`$ and $`n1`$ by $`n+1`$. This proves the statement for $`A`$. The statement for $`B`$ is proved in the same way. ∎
The following are two corollaries of Theorem 5.1. The first is the input for Section 6. The second is important for showing that it is only necessary to state the ‘binding’ condition (3) once (for $`m=0`$). It is a trivial consequence of Theorem 5.1.
###### 5.3 THEOREM (Minimizing sequence).
Suppose that $`m_1>m_2>\mathrm{}>0`$ is a sequence tending to zero and suppose that $`\mathrm{\Phi }_j`$, for $`j=1,2,\mathrm{}`$ is an approximate minimizer for $`H_{m_j}^V`$ in the sense that
$$\delta _j(\mathrm{\Phi }_j,H_{m_j}^V\mathrm{\Phi }_j)E^V(m_j,N)0\mathrm{as}j\mathrm{}.$$
(39)
Then $`\mathrm{\Phi }_1,\mathrm{\Phi }_2,\mathrm{}`$ is a minimizing sequence for $`H_0^V`$.
###### Proof.
$`E^V(m_j,N)+\delta _j=(\mathrm{\Phi }_j,H_{m_j}^V\mathrm{\Phi }_j)>(\mathrm{\Phi }_j,H_0^V\mathrm{\Phi }_j)E^V(0,N).`$
###### 5.4 LEMMA (Binding without mass implies uniform binding with mass).
Assume (3), i.e., assume that
$$E^V(0,N)<\mathrm{min}\{E^V(0,N^{})+E^0(0,NN^{}):N^{}<N\}2\epsilon .$$
Then, for all sufficiently small $`m`$,
$$E^V(m,N)<\mathrm{min}\{E^V(m,N^{})+E^0(m,NN^{}):N^{}<N\}\epsilon .$$
###### Proof.
Each of the various energies converges as $`m0`$ by Theorem 5.1. ∎
## 6 TWO INFRARED BOUNDS
We have seen that any sequence of (approximate) minimizers for the $`H_m^V`$ problem ($`m>0`$) is a minimizing sequence for the $`H_0^V`$ problem. A crucial point was the possibility of finding an approximate minimizer for the $`H_0^V`$ problem that has a finite expectation value for the total photon number operator $`𝒩`$.
We also know by Corollary 5.4 and Theorem 4.1 that the $`H_m^V`$ problem has a ground state, $`\mathrm{\Phi }_m`$, and in this section we shall prove two theorems about the soft photon behavior of $`\mathrm{\Phi }_m`$. The first, Theorem 6.1, is about the photon number $`a_\lambda ^{}(k)a_\lambda ^{}(k)`$ and is based on a method of . The second, Theorem 6.3, which has no antecedent we are aware of, is about the derivative of $`a_\lambda ^{}(k)a_\lambda ^{}(k)`$ with respect to $`k`$.
###### 6.1 THEOREM (Photon number bound).
Assume that there is binding, i.e., $`E^V(m,N)<\mathrm{\Sigma }^V(m,N)`$. Assume that $`\mathrm{\Phi }_m`$ is a normalized ground state for the many-body Hamiltonian $`H_m^V`$, $`m0`$. Then
$$(\mathrm{\Phi }_m,a_\lambda ^{}(k)a_\lambda ^{}(k)\mathrm{\Phi }_m)<\frac{P\alpha }{|k|}(1+g^2)\chi _\mathrm{\Lambda }(k),$$
where $`P`$ is a finite constant independent of $`\mathrm{\Phi }_m`$, $`g`$, $`\alpha `$, and depends on $`m`$ only via the binding energy $`\mathrm{\Sigma }^V(m,N)E^V(m,N)>0`$ and of course on $`\mathrm{\Lambda }`$ . The function $`\chi _\mathrm{\Lambda }(k)`$ is the characteristic function of the ball of radius $`\mathrm{\Lambda }`$.
Remark: We have proved in Corollary 5.4 that the binding energy is a uniformly positive function of $`m`$ for small $`m`$, if the binding energy at $`m=0`$ is not zero. Therefore, Theorem 6.1 implies that the number operator $`𝒩`$ in (37) is uniformly bounded for small $`m`$.
The proof of Theorem 6.1 will be based on the following lemma about exponential decay of eigenfunctions, concerning which there is a vast literature (see ). We do not strive at all to get the best exponential decay constants. For us the only really relevant goal is a decay estimate that depends only on $`\mathrm{\Sigma }^V(m,N)E^V(m,N)>0`$, but not otherwise on $`m`$.
###### 6.2 LEMMA (Exponential decay).
Let $`H_m^V`$ be the N-body Hamiltonian in (4) and let $`\mathrm{\Phi }_m`$ be a groundstate wave function, which necessarily satisfies the Schrödinger equation
$$H_m^V\mathrm{\Phi }_m=E^V(m,N)\mathrm{\Phi }_m.$$
(40)
We assume that $`\mathrm{\Sigma }^V(m,N)E^V(m,N)>0`$ and choose $`\beta >0`$ with $`\beta ^2<\mathrm{\Sigma }^V(m,N)E^V(m,N)`$. Then
$$\mathrm{exp}(\beta |X|)\mathrm{\Phi }_m^2C(1+\frac{1}{\mathrm{\Sigma }^V(m,N)E^V(m,N)\beta ^2})\mathrm{\Phi }_m^2$$
(41)
where the constant $`C`$ does not depend on $`m`$.
The strategy of the following proof is probably due to Agmon . We learned it from .
###### Proof.
Let $`G(X)`$ any smooth, bounded function on $`^{3N}`$ with bounded first derivative. We easily compute
$$[[H_m^VE^V(m,N),G],G]=2|G|^2.$$
(42)
We use the Schrödinger equation (40) to compute
$$(G\mathrm{\Phi }_m,[H_m^VE^V(m,N)]G\mathrm{\Phi }_m)=\frac{1}{2}(\mathrm{\Phi }_m,[[H_m^VE(m),G],G]\mathrm{\Phi }_m)=(\mathrm{\Phi }_m,|G|^2\mathrm{\Phi }_m).$$
(43)
Now we choose $`G`$ to be
$`G(X)`$ $`=`$ $`\chi (X/R)\mathrm{exp}[f(X)],\mathrm{where}`$
$`f(X)`$ $`=`$ $`\left[{\displaystyle \frac{\beta |X|}{1+\epsilon |X|}}\right],`$ (44)
and where $`0\chi 1`$ is a smooth cutoff function that is identically equal to 1 outside the ball of radius 2, and identically zero inside the ball of radius 1. We let $`\epsilon 0`$ at the end.
Next, we calculate
$$|G|^2=|\chi |^2e^{2f}+2\chi fe^fG+|f|^2G^2,$$
and note that the first and second terms are compactly supported in $`^{3N}`$ and each is bounded by a constant $`C`$ that depends on $`\beta `$ and $`R`$.
Returning to (43), we obtain, after rearranging terms,
$$(G\mathrm{\Phi }_m,(H_m^VE^V(m,N)|f|^2))G\mathrm{\Phi }_m)C\mathrm{\Phi }_m^2.$$
(45)
Since $`|f|\beta `$ we know by Corollary A.2 that
$`\chi \left(H_m^VE^V(m,N)|f|^2\right)\chi `$ (46)
$``$ $`\left(\mathrm{\Sigma }^V(m,N)E^V(m,N)\beta ^2o(1)\right)\chi ^2(R\mathrm{}).`$
In conjunction with (45) this shows that
$$G\mathrm{\Phi }_m^2\frac{2C}{\mathrm{\Sigma }^V(m,N)E^V(m,N)\beta ^2}\mathrm{\Phi }_m^2$$
for $`R`$ large enough. After letting $`\epsilon 0`$ by monotone convergence a similar bound with $`G`$ replaced by $`\chi \mathrm{exp}(\beta |X|)`$ is obtained. ∎
###### Proof of Theorem 6.1.
This proof is a slight modification of the one in . The basic idea is to show that there is effectively no interaction between localized particles and low momentum photons. To make this idea explicit we write our Hamiltonian in a gauge different from the usual Coulomb gauge.
To be precise, define
$$\stackrel{~}{A}(x)=A(x)A(0),$$
(47)
which is well defined owing to the ultraviolet cutoff. The unitary operator that accomplishes this is $`U=\mathrm{exp}[i_{j=1}^N\sqrt{\alpha }x_jA(0)]`$. This is an ‘operator-valued gauge transformation’. It commutes with $`A(x)`$ for all $`x`$, but not with $`a_\lambda ^{}(k)`$ or with $`H_f`$.
Define
$$b_\lambda (k,X)=Ua_\lambda ^{}(k)U^{}=a_\lambda ^{}(k)iw_\lambda (k,X),$$
(48)
with $`w_\lambda (k,X)=\chi _\mathrm{\Lambda }(k)|k|^{1/2}\epsilon _\lambda (k)_{j=1}^Nx_j`$. The transformed Hamiltonian $`\stackrel{~}{H}_m`$ is
$$\stackrel{~}{H}_m=UH_m^VU^{}=\underset{j=1}{\overset{N}{}}\left\{\left(p_j+\sqrt{\alpha }\stackrel{~}{A}(x_j)\right)^2+\frac{g}{2}\sqrt{\alpha }\sigma _jB(x_j)\right\}+V+\stackrel{~}{H}_f(m)$$
(49)
$$\stackrel{~}{H}_f(m)=\underset{\lambda =1,2}{}_{}^3\omega _m(k)b_\lambda ^{}(k,X)b_\lambda (k,X).$$
(50)
To estimate $`a_\lambda ^{}(k)\mathrm{\Phi }_m`$, write
$$a_\lambda ^{}(k)\mathrm{\Phi }_m=U^{}a_\lambda ^{}(k)\stackrel{~}{\mathrm{\Phi }}_miw_\lambda (k,X)\mathrm{\Phi }_m$$
(51)
where $`\stackrel{~}{\mathrm{\Phi }}_m=U\mathrm{\Phi }_m`$, and note that
$$w_\lambda (k,X)\mathrm{\Phi }_m\sqrt{N}\frac{\chi _\mathrm{\Lambda }(k)}{|k|^{1/2}}|X|\mathrm{\Phi }_m.$$
(52)
It remains to estimate $`a_\lambda ^{}(k)\stackrel{~}{\mathrm{\Phi }}_m`$. By the Schrödinger equation for $`\stackrel{~}{\mathrm{\Phi }}_m`$
$`\left(\stackrel{~}{H}_mE^V(m,N)\right)a_\lambda ^{}(k)\stackrel{~}{\mathrm{\Phi }}_m`$ $`=`$ $`[\stackrel{~}{H}_m,a_\lambda ^{}(k)]\stackrel{~}{\mathrm{\Phi }}_m`$ (53)
$`=`$ $`2\sqrt{\alpha }|k|^{1/2}\epsilon _\lambda (k){\displaystyle \underset{j=1}{\overset{N}{}}}(p_j+\sqrt{\alpha }\stackrel{~}{A}(x_j))(1e^{ikx_j})\stackrel{~}{\mathrm{\Phi }}_m`$
$`+i{\displaystyle \frac{g}{2}}\sqrt{\alpha }{\displaystyle \frac{k\epsilon _\lambda (k)}{\sqrt{|k|}}}{\displaystyle \underset{j=1}{\overset{N}{}}}\sigma _je^{ikx_j}\stackrel{~}{\mathrm{\Phi }}_m`$
$`\omega _m(k)b_\lambda (k,X)\stackrel{~}{\mathrm{\Phi }}_m.`$
While this equation is correct, the derivation is somewhat formal. This is rigorously justified in Appendix B.
Now add $`\omega _m(k)a_\lambda ^{}(k)\stackrel{~}{\mathrm{\Phi }}_m`$ on both sides. Since $`E^V(m,N)`$ is the ground state energy and $`\omega _m(k)>0`$ the operator $`\stackrel{~}{H}_mE^V(m,N)+\omega _m(k)`$ has a bounded inverse $`R(\omega _m(k))`$ and hence
$$\begin{array}{cc}\hfill a_\lambda ^{}(k)\stackrel{~}{\mathrm{\Phi }}_m=& 2\sqrt{\alpha }R(\omega _m(k))\frac{\chi _\mathrm{\Lambda }(k)}{|k|^{1/2}}\epsilon _\lambda (k)\underset{j=1}{\overset{N}{}}(p_j+\sqrt{\alpha }\stackrel{~}{A}(x_j))(1e^{ikx_j})\stackrel{~}{\mathrm{\Phi }}_m\hfill \\ & +iR(\omega _m(k))\chi _\mathrm{\Lambda }(k)\frac{g}{2}\sqrt{\alpha }\frac{k\epsilon _\lambda (k)}{\sqrt{|k|}}\underset{j=1}{\overset{N}{}}\sigma _je^{ikx_j}\stackrel{~}{\mathrm{\Phi }}_m\hfill \\ & R(\omega _m(k))\omega _m(k)iw_\lambda (k,X)\stackrel{~}{\mathrm{\Phi }}_m.\hfill \end{array}$$
(54)
For consistency, note that for $`|k|>\mathrm{\Lambda }`$ $`w_\lambda (k,X)=0`$ and hence $`a_\lambda ^{}(k)\stackrel{~}{\mathrm{\Phi }}_m=0`$, i.e, for these modes $`\stackrel{~}{\mathrm{\Phi }}_m`$ is the vacuum as it should be for a minimizer. Since $`R(\omega _m(k))\omega _m(k)^1`$ the norm of the last term is bounded by
$$\sqrt{\alpha }\sqrt{N}\frac{\chi _\mathrm{\Lambda }(k)}{|k|^{1/2}}|X|\stackrel{~}{\mathrm{\Phi }}_m.$$
(55)
To bound the norm of the first term we need to estimate
$`{\displaystyle \underset{j=1}{\overset{N}{}}}R(\omega _m(k))\epsilon _\lambda (k)(p_j+\sqrt{\alpha }\stackrel{~}{A}(x_j))(1e^{ikx_j})\stackrel{~}{\mathrm{\Phi }}_m`$ (56)
$`=`$ $`\underset{\eta 1}{sup}\left|{\displaystyle \underset{j=1}{\overset{N}{}}}(\epsilon _\lambda (k)(p_j+\sqrt{\alpha }\stackrel{~}{A}(x_j))R(\omega _m(k))\eta ,(1e^{ikx_j})\stackrel{~}{\mathrm{\Phi }}_m)\right|`$
$``$ $`\underset{\eta 1}{sup}\left[{\displaystyle \underset{j=1}{\overset{N}{}}}(p_j+\sqrt{\alpha }\stackrel{~}{A}(x_j))R(\omega _m(k))\eta ^2\right]^{1/2}\left[{\displaystyle \underset{j=1}{\overset{N}{}}}(1e^{ikx_j})\stackrel{~}{\mathrm{\Phi }}_m^2\right]^{1/2}`$
Next estimate the square of the first factor to get
$`(\eta ,R(\omega _m(k))\left[{\displaystyle \underset{j=1}{\overset{N}{}}}(p_j+\sqrt{\alpha }\stackrel{~}{A}(x_j))^2\right]R(\omega _m(k))\eta )`$ (57)
$``$ $`a(\eta ,R(\omega _m(k))H_m^VR(\omega _m(k))\eta )+b`$
$``$ $`a(\eta ,R(\omega _m(k))\eta )+(aE^V(m,N)+b)(\eta ,R(\omega _m(k))^2\eta )`$
$``$ $`C{\displaystyle \frac{(\mathrm{\Lambda }+1)}{|k|^2}}\text{for}|k|\mathrm{\Lambda }`$
where $`a`$ and $`b`$ are independent of $`m`$. Since $`sup_{m<1}E^V(m,N)<\mathrm{}`$ the constant $`C`$ is also independent of $`m`$. Finally the second factor in (56) is bounded by $`|k||X|\stackrel{~}{\mathrm{\Phi }}_m`$. The term containing the Pauli matrices in (54) is estimated similarly. In conjuction with (54), (52)(56) and (57) this shows that
$$a_\lambda ^{}(k)\mathrm{\Phi }_mC\sqrt{\alpha }(\mathrm{\Lambda }+1)^{1/2}\frac{\chi _\mathrm{\Lambda }(k)}{|k|^{1/2}}|X|\stackrel{~}{\mathrm{\Phi }}_m$$
(58)
This, together with Lemma 6.2, proves the theorem. ∎
Next we differentiate (54) with respect to $`k`$. There is a slight problem with this calculation since the polarization vectors cannot be defined in a smooth fashion globally. We make the following choice for the polarization vectors.
$$\epsilon _1(k)=\frac{(k_2,k_1,0)}{\sqrt{k_1^2+k_2^2}}$$
(59)
and
$$\epsilon _2(k)=\frac{k}{|k|}\epsilon _1(k).$$
(60)
###### 6.3 THEOREM (Photon derivative bound).
Assume that there is binding, i.e., $`\mathrm{\Sigma }^V(m,N)E^V(m,N)>0`$. Assume that $`\mathrm{\Phi }_m`$ is a normalized ground state for the many-body Hamiltonian $`H_m^V`$, $`m0`$. Then for $`|k|<\mathrm{\Lambda }`$ and $`(k_1,k_2)(0,0)`$
$$_ka_\lambda ^{}(k)\mathrm{\Phi }_m<\frac{Q\sqrt{\alpha }(1+|g|)}{|k|^{1/2}\sqrt{k_1^2+k_2^2}},$$
(61)
where $`Q`$ is a finite constant independent of $`\mathrm{\Phi }_m`$, $`g`$, $`\alpha `$, $`\mathrm{\Lambda }`$, and depends on $`m`$ only through the binding energy $`\mathrm{\Sigma }^V(m,N)E^V(m,N)>0`$ .
###### Proof.
We differentiate 54 with respect to $`k`$ and obtain
$`_k(a_\lambda ^{}(k))\stackrel{~}{\mathrm{\Phi }}_m=`$
$`2\sqrt{\alpha }R(\omega _m(k))^2{\displaystyle \frac{k}{|k|}}\epsilon _\lambda (k){\displaystyle \underset{j=1}{\overset{N}{}}}(p_j+\sqrt{\alpha }\stackrel{~}{A}(x_j)){\displaystyle \frac{(1e^{ikx_j})}{|k|^{1/2}}}\stackrel{~}{\mathrm{\Phi }}_m+`$
$`2\sqrt{\alpha }R(\omega _m(k))_k(\epsilon _\lambda (k)){\displaystyle \underset{j=1}{\overset{N}{}}}(p_j+\sqrt{\alpha }\stackrel{~}{A}(x_j)){\displaystyle \frac{(1e^{ikx_j})}{|k|^{1/2}}}\stackrel{~}{\mathrm{\Phi }}_m+`$
$`2\sqrt{\alpha }R(\omega _m(k))\epsilon _\lambda (k){\displaystyle \underset{j=1}{\overset{N}{}}}(p_j+\sqrt{\alpha }\stackrel{~}{A}(x_j))_k\left({\displaystyle \frac{(1e^{ikx_j})}{|k|^{1/2}}}\right)\stackrel{~}{\mathrm{\Phi }}_m`$
$`+{\displaystyle \frac{g}{2}}\sqrt{\alpha }_k\left(iR(\omega _m(k)){\displaystyle \frac{k\epsilon _\lambda (k)}{\sqrt{|k|}}}{\displaystyle \underset{j=1}{\overset{N}{}}}\sigma _je^{ikx_j}\stackrel{~}{\mathrm{\Phi }}_m\right)`$
$`_k\left(R(\omega _m(k))\omega _m(k)iw_\lambda (k,X)\stackrel{~}{\mathrm{\Phi }}_m\right).`$
The norms of the first and third terms can estimated precisely the same way as in (56) and (57), and yields a bound of the form
$$\frac{C\sqrt{\alpha }}{|k|^{3/2}}\left(1+|X|\right)\stackrel{~}{\mathrm{\Phi }}_m.$$
(62)
For the second term, a straightforward calculation shows that
$$|_k\epsilon _i(k)|\frac{\mathrm{const}.}{\sqrt{k_1^2+k_2^2}}\mathrm{for}i=1,2.$$
(63)
The last term is dealt with in a similar fashion as the previous ones. Using the steps in (56) and (57), this leads to the bound
$$_k(a_\lambda ^{}(k))\stackrel{~}{\mathrm{\Phi }}_m\frac{C\sqrt{\alpha }}{|k|^{1/2}\sqrt{k_1^2+k_2^2}}\left(1+|X|\right)\stackrel{~}{\mathrm{\Phi }}_m.$$
(64)
The fourth term can be estimated in the same fashion to yield a similar result.
Differentiating (51) leads to the same estimate with $`\stackrel{~}{\mathrm{\Phi }}_m`$ replaced by $`\mathrm{\Phi }_m`$. This, together with Lemma 6.2, proves the theorem. ∎
As for the proof of Theorem 6.1, our somewhat formal calculations above are rigorously justified in Appendix B.
## 7 PROOF OF THEOREM 2.1
The proof will be done in two steps.
###### Proof.
Step 1. The Hamiltonian $`H_m^V`$ has a normalized ground state $`\mathrm{\Phi }_m`$, by Theorem 4.1. Pick a sequence $`m_1>m_2>\mathrm{}`$ tending to zero and denote the corresponding eigenvectors by $`\mathrm{\Phi }_j`$. This sequence is a minimizing sequence for $`H_0^V`$ by Theorem 5.3. Since $`\mathrm{\Phi }_j`$ is bounded there is a subsequence (call it again $`\mathrm{\Phi }_j`$) which has a weak limit $`\mathrm{\Phi }`$. Since $`H_0^VE^V(0,N)0`$ and by the lower semi-continuity of non-negative quadratic forms (in our case, $`H_0^VE^V(0,N)`$)
$$0(\mathrm{\Phi },(H_0^VE^V(0,N))\mathrm{\Phi })\underset{j\mathrm{}}{lim\; inf}(\mathrm{\Phi }_j,(H_0^VE^V(0,N))\mathrm{\Phi }_j)=0.$$
Hence $`\mathrm{\Phi }`$ will be a (normalized) ground state if we show that $`\mathrm{\Phi }=1`$ (i.e. $`\mathrm{\Phi }_j\mathrm{\Phi }`$ strongly). It is important to note, however that if we write $`\mathrm{\Phi }_j=\{\mathrm{\Phi }_{j,0},\mathrm{\Phi }_{j,1},\mathrm{},\mathrm{\Phi }_{j,n},\mathrm{}\}`$, where $`\mathrm{\Phi }_{j,n}`$ is the $`n`$-photon component of $`\mathrm{\Phi }_j`$ then it suffices to prove the $`L^2`$ norm-convergence of each $`\mathrm{\Phi }_{j,n}`$. The reason is the uniform bound on the total average photon number; see the remark after Theorem 6.1 which implies
$$\underset{nN}{}\mathrm{\Phi }_{j,n}^2\text{const}N^1.$$
Likewise, it suffices to prove the strong $`L^2`$ convergence in the bounded domain in which $`|X|<R`$ for each finite $`R`$. The reason for this is the exponential decay given in Lemma 6.2, which is uniform by Lemma 5.4. Finally, by Theorem 6.1 $`\mathrm{\Phi }_{j,n}(X,k_1,\mathrm{},k_n)`$ vanishes if $`|k_i|>\mathrm{\Lambda }`$ for some $`i`$. So it suffices to show $`L^2`$ convergence for $`\mathrm{\Phi }_{j,n}`$ restricted to
$$\mathrm{\Omega }=\{(X,k_1,\mathrm{},k_n):|X|<R;|k_i|<\mathrm{\Lambda },i=1,\mathrm{},n\}^{3(N+n)}$$
for each $`R>0`$.
Step 2. For each $`p<2`$ and $`R>0`$ we show that $`\mathrm{\Phi }_{j,n}`$ restricted to $`\mathrm{\Omega }`$ is a bounded sequence in $`W^{1,p}(\mathrm{\Omega })`$. The key to this bound is (61) and
$$(a_\lambda ^{}(k)\mathrm{\Phi }_j)_{n1}(X,k_1,\mathrm{},k_{n1})=\sqrt{n}\mathrm{\Phi }_{j,n}(X,k,k_1,\mathrm{},k_{n1})$$
(65)
where the arguments $`\lambda ,\lambda _1,\mathrm{},\lambda _{n1}`$ and the spin indices have been supressed. By the symmetry of $`\mathrm{\Phi }_{j,n}`$, (65), Hölder’s inequality and (61)
$$\begin{array}{c}_{B_R}𝑑X_{|k_1|,\mathrm{},|k_n|<\mathrm{\Lambda }}𝑑k_1\mathrm{}𝑑k_n\underset{i=1}{\overset{n}{}}|_{k_i}\mathrm{\Phi }_{j,n}(X,k_1,\mathrm{},k_n)|^p\hfill \\ \hfill =n^{1p/2}_{B_R}dX_{|k_1|,\mathrm{},|k_n|<\mathrm{\Lambda }}dk_1\mathrm{}dk_n|_{k_1}(a_\lambda ^{}(k_1)\mathrm{\Phi }_j)_{n1}(X,k_2,\mathrm{},k_n)|^p\\ \hfill C_{|k_1|<\mathrm{\Lambda }}dk_1(_{B_R}dX_{|k_2|,\mathrm{},|k_n|<\mathrm{\Lambda }}dk_2\mathrm{}dk_n|_{k_1}(a_\lambda ^{}(k_1)\mathrm{\Phi }_j)_{n1}(X,k_2,\mathrm{},k_n)|^2)^{p/2}\\ \hfill C_{|k_1|<\mathrm{\Lambda }}𝑑k_1_{k_1}a_\lambda ^{}(k_1)\mathrm{\Phi }_j^p\text{const}\end{array}$$
(66)
independent of $`j`$. The constant $`C`$ depends on all the parameters, but is finite because $`|k_i|\mathrm{\Lambda }`$ in the integration. Similarly, by Hölder’s inequality
$`\chi (|X|<R)_X\mathrm{\Phi }_{j,n}_p^p`$ $``$ $`C\chi (|X|<R)_X\mathrm{\Phi }_{j,n}_2^p`$
$``$ $`C(\mathrm{\Phi }_j,{\displaystyle \underset{i=1}{\overset{N}{}}}p_i^2\mathrm{\Phi }_j)^{p/2}`$
which is uniformly bounded by Lemma (A.5).
Since the classical derivative of $`a_\lambda ^{}(k)\mathrm{\Phi }_m`$ is not defined in all of $`\mathrm{\Omega }`$ one has to check that the weak derivative coincides with the classical derivative a.e.. Because of our definitions (59) and (60), the classical derivative is not defined along the $`3`$-axis.
One has to show that
$$_\mathrm{\Omega }_i\psi \mathrm{\Phi }_{j,n}=\underset{\epsilon 0}{lim}_{\mathrm{\Omega }_\epsilon }_i\psi \mathrm{\Phi }_{j,n}=\underset{\epsilon 0}{lim}_{\mathrm{\Omega }_\epsilon }\psi _i\mathrm{\Phi }_{j,n}i=1,\mathrm{},3(N+n).$$
for any test function $`\psi C_c^{\mathrm{}}(\mathrm{\Omega })`$. Here $`\mathrm{\Omega }_\epsilon `$ is $`\mathrm{\Omega }`$ with an $`\epsilon `$ cylinder around the $`3`$-axis removed in each $`k`$-ball. The first equality is trivial; it is the second equality that has to be checked. This amounts to showing that the boundary term, coming from the integration by parts vanishes in the limit as $`\epsilon `$ tends to zero. But this follows immediately from Theorem 6.1.
This shows that $`\mathrm{\Phi }_{j,n}`$ as a function of all its $`3(N+n)`$ variables, is in the Sobolev space $`W^{1,p}(\mathrm{\Omega })`$ and that $`sup_j\mathrm{\Phi }_{j,n}_{W^{1,p}(\mathrm{\Omega })}<\mathrm{}`$. Since $`\mathrm{\Phi }_{n,j}`$ converges weakly in $`L^2(\mathrm{\Omega })`$ it converges weakly in $`L^p(\mathrm{\Omega })`$ and since the sequence is bounded in $`W^{1,p}(\mathrm{\Omega })`$, $`\mathrm{\Phi }_{n,j}`$ converges weakly to $`\mathrm{\Phi }_n`$.
The Rellich-Kondrachov theorem (see Theorem 8.9) states that such a sequence converges strongly in $`L^q(\mathrm{\Omega })`$ if $`1q<\left[\mathrm{\hspace{0.17em}3}p(N+n)/3(N+n)p\right]`$. The boundedness of $`\mathrm{\Omega }`$ is crucial here. For our purposes we need $`q=2`$, and hence we have to pick $`p`$ such that
$$2>p>\frac{23(N+n)}{2+3(N+n)}$$
(67)
which is possible for each $`N`$ and $`n`$. We conclude that $`\mathrm{\Phi }_{j,n}\mathrm{\Phi }_n`$ strongly in $`L^2(\mathrm{\Omega })`$ as $`j\mathrm{}`$, for each $`n`$ and $`R`$. This proves the theorem. ∎
Remark: Theorem 6.3 essentially says that the derivative is almost, but not quite in $`L^2`$. For high dimensions, $`3(N+n)`$, the required $`p`$ is as close as we please to 2 if we require $`q=2`$, but $`p=2`$ is not allowed. The way out of the difficulty was to prove a uniform bound on the number operator and use this to say that that it suffices to prove strong convergence for each $`n`$ separately. With $`n`$ then fixed, it is possible to find a $`p<2`$ that yields $`q=2`$. It is, therefore, crucial to have the derivative in every $`L^p`$ space with $`p<2`$. The resolution of the problem of the infrared singularity is thus seen to be a delicate matter.
## Appendix A Appendix: LOCALIZATION ESTIMATES
In this appendix we collect a few facts which we use several times in this paper. Generally we worry about localizations of Hamiltonians in configuration space. While this is standard for Schrödinger operators it is somewhat more complicated in the presence of the radiation field. This is chiefly due to the problem of localization of photons.
We begin by stating a few well known facts about partitions of unity. Let $`\beta `$ denote one of the $`2^N`$ subsets of the set of integers $`1,2,\mathrm{},N`$. Its complement is denoted by $`\beta ^c`$. As shown in there exists a family of smooth functions $`j_\beta `$ having the following four properties.
(i)
$$\underset{\beta }{}j_\beta ^2=1.$$
(68)
(ii) For $`\beta \{1,\mathrm{},N\}`$ (including the empty set) the $`j_\beta `$’s are homogeneous of degree $`0`$ and live outside the ball of radius $`R`$ centered at the origin and
$$suppj_\beta \{X:\underset{i\beta ,j\beta ^c}{\mathrm{min}}(|x_ix_j|,|x_j|)c|X|\},$$
(69)
where $`C`$ is some positive constant.
(iii) In the case where $`\beta =\{1,\mathrm{},N\}`$, $`j_\beta `$ is compactly supported.
Corresponding to these electron localizations we define photon localizations.
For given $`\beta \{1,\mathrm{},N\}`$ consider the function
$$g_1(y;\beta ,X)=\mathrm{\Pi }_{j\beta ^c}\left(1\chi (\frac{yx_j}{P})\right)$$
(70)
where $`\chi `$ is a smoothed characteristic function of the unit ball. Define $`g_2(y;\beta ,X)=1g_1(y,\beta ,X)`$. In the variable $`y`$, the function $`g_1`$ is supported away from the particles in $`\beta ^c`$ while $`g_2`$ lives close to the particles in $`\beta ^c`$. Next define, for $`i=1,2`$,
$$j_i(y;\beta ,X)=\frac{g_i(y;\beta ,X)}{\sqrt{g_1(y;\beta ,X)^2+g_2(y;\beta ,X)^2}}.$$
(71)
Certainly $`j_1^2+j_2^2=1`$ and a simple computation shows that
$$|j_i|\frac{\mathrm{const}.}{P}.$$
(72)
In the case where $`\beta =\{1,\mathrm{},N\}`$ the construction of $`j_1`$ and $`j_2`$ is similar to the above one except that the function $`g_1`$ depends on $`y`$, is equal to one in a neighborhood of the origin and is compactly supported.
With the help of $`j_1`$ and $`j_2`$ the photons can now be localized as was done in Section 4. Let $`U_\beta (X)`$ be the corresponding isometric transformation, i.e., the one that is defined via the relation
$$U_\beta (X)a^\mathrm{\#}(h)U_\beta ^{}(X)=a^\mathrm{\#}(j_1h)+a^\mathrm{\#}(j_2h).$$
(73)
The tensor product indicated is a tensor product between Fock spaces.
We denote by $`H_\beta `$ the Hamiltonian of the form (4) with photon mass, but only for the particles in the set $`\beta `$. More precisely this operator acts on $`L^2(^{3|\beta |})`$. By $`H^{\beta ^c}`$ we denote the Hamiltonian of the form (4) with photon mass, but only for the particles in the set $`\beta ^c`$ where the interaction with the nuclei has been dropped. This operator acts on $`L^2(^{3|\beta ^c|})`$ In particular we keep the interaction among those particles. In the case where $`\beta =\{1,\mathrm{},N\}`$ the Hamiltonian $`H^{\beta ^c}=H_f(m)`$.
###### A.1 LEMMA (Localization of Hamiltonian).
For every $`\beta `$
$$j_\beta Hj_\beta =U_\beta ^{}(X)j_\beta \left[H_\beta +H^{\beta ^c}\right]j_\beta U_\beta (X)+o(1).$$
(74)
For $`\beta \{1,\mathrm{},N\}`$, $`o(1)0`$ as first $`R\mathrm{}`$ and then $`P\mathrm{}`$. If $`\beta =\{1,\mathrm{},N\}`$ then $`o(1)0`$ as $`P\mathrm{}`$ for every fixed $`R>0`$.
###### Proof.
Our immediate aim is to compare the field energy $`H_f`$ with the localized field energy $`U_\beta ^{}(X)[H_f+H_f]U_\beta (X)`$. For simplicity the various indices are supressed and $`U_\beta (X)`$ is replaced by $`U_\beta `$. The variable $`X`$ plays no role here. Pick an orthonormal basis $`\{g_j\}_{j=1}^{\mathrm{}}`$ of $`L^2(^3)`$ in $`H^{1/2}(^3)`$. States of the form
$$\zeta =const.a_{\lambda _{i_1}}^{}(g_{i_1})\mathrm{}a_{\lambda _{i_k}}^{}(g_{i_k})|0>$$
(75)
where $`k`$ is finite, form an orthonormal basis in the Fock space. The field energy acts on such states as
$$H_f\zeta =\underset{j=1}{\overset{k}{}}a_{\lambda _{i_1}}^{}(g_{i_1})\mathrm{}a_{\lambda _{i_j}}^{}(\omega g_{i_j})\mathrm{}a_{\lambda _{i_k}}^{}(g_{i_k})|0>.$$
(76)
Thus, we have that
$$H_f\zeta =U_\beta ^{}\underset{j=1}{\overset{k}{}}c_{\lambda _{i_1}}^{}(g_{i_1})\mathrm{}c_{\lambda _{i_j}}^{}(\omega g_{i_j})\mathrm{}c_{\lambda _{i_k}}^{}(g_{i_k})U_\beta |0>$$
(77)
and
$$H_f=U_\beta ^{}\left[H_f+H_f\right]U_\beta +E_f$$
(78)
where the error $`E_f`$ is given by
$$\begin{array}{cc}\hfill E_f\zeta =U_\beta ^{}\underset{j=1}{\overset{k}{}}c_{\lambda _{i_1}}^{}(g_{i_1})\mathrm{}& (a_{\lambda _{i_j}}^{}([j_1,\omega ]g_{i_j})+a_{\lambda _{i_j}}^{}([j_2,\omega ]g_{i_j}))\hfill \\ & \mathrm{}c_{\lambda _{i_k}}^{}(g_{i_k})U_\beta |0>.\hfill \end{array}$$
(79)
Thus $`E_f`$ is given by the operator (the $`\lambda `$’s are omitted)
$`E_f=`$ $`U_\beta ^{}{\displaystyle \underset{k}{}}\left[a^{}([j_1,\omega ]g_k)+a^{}([j_2,\omega ]g_k)\right]\left[a^{}(j_1g_k)+a^{}(j_2g_k)\right]U_\beta .`$ (80)
The expression for the operator $`E_f`$ does not look hermitian but it is, remembering that $`U_\beta ^{}`$ is a partial isometry. Standard estimates lead to
$$|(\mathrm{\Psi },E_f\mathrm{\Psi })|\left([\omega ,j_1]+[\omega ,j_2]\right)(\mathrm{\Psi },\left[𝒩+1\right]\mathrm{\Psi }).$$
(81)
where $`𝒩`$ is the number operator.
Here $`[j_1,\omega ]`$ denotes the operator norm associated with the kernel $`[j_1,\omega ]`$. This norm can be estimated using the formula
$$[j_1,\omega ]=[j_1,\omega ^2]\frac{1}{\omega }+\omega ^2[j_1,\frac{1}{\omega }].$$
(82)
Recalling the definition of $`j_1`$, the operator norm of the first term is easily seen to be bounded by a $`const./P`$. Likewise, the second term, using the formula
$$\frac{1}{\sqrt{p^2+m^2}}=\frac{1}{\pi }_0^{\mathrm{}}\frac{1}{t+p^2+m^2}\frac{dt}{\sqrt{t}},$$
(83)
can be estimated by $`const./P`$. The term $`[j_2,\omega ]`$ is estimated in a similar fashion. The estimate (81) immediately shows that for a general state $`\mathrm{\Phi }`$ we have that
$$|(\mathrm{\Phi },\left[H_fU_\beta ^{}\left[H_f+H_f\right]U_\beta \right]\mathrm{\Phi })|\frac{const.}{P}(\mathrm{\Phi },𝒩\mathrm{\Phi }).$$
(84)
Since the photons have a mass we can estimate the number operator in terms of the field energy. The field energy is relatively bounded with respect to the Hamiltonian, i.e., $`H_faH_m^V+b`$ for some positive constants $`a`$ and $`b`$, and thus we obtain
$$|(\mathrm{\Phi },\left[H_fU_\beta ^{}\left[H_f+H_f\right]U_\beta \right]\mathrm{\Phi })|\frac{const.}{Pm}(\mathrm{\Phi },[aH_m^V+b]\mathrm{\Phi }).$$
(85)
Note that this estimate had nothing to do with the electron, in particular the $`x`$–space cutoff is not present in the calculation.
Next we have to compare $`_{j=1}^N(p_j+\sqrt{\alpha }A(x_j))^2`$ with
$$U_\beta ^{}(X)j_\beta \left[\underset{i\beta }{}(p_i+\sqrt{\alpha }A(x_i))^2+\underset{j\beta ^c}{}(p_j+\sqrt{\alpha }A(x_j))^2\right]j_\beta U_\beta (X).$$
This time the $`X`$–space cutoff is important. We would like to estimate the difference
$`j_\beta {\displaystyle \underset{i\beta }{}}\left[(p_i+\sqrt{\alpha }A(x_i))^2U_\beta ^{}(X)(p_i+\sqrt{\alpha }A(x_i))^2U_\beta (X)\right]j_\beta `$ (86)
$`+`$ $`j_\beta {\displaystyle \underset{i\beta ^c}{}}\left[(p_i+\sqrt{\alpha }A(x_i))^2U_\beta ^{}(X)(p_i+\sqrt{\alpha }A(x_i))^2U_\beta (X)\right]j_\beta .`$
It suffices to treat the first term, the other is similar. It can be easily expressed as
$$j_\beta \left[\underset{i\beta }{}(p_i+\sqrt{\alpha }A(x_i))Q_i+Q_i(p_i+\sqrt{\alpha }A(x_i))Q_i^2\right]j_\beta $$
(87)
where
$$Q_i=p_i+\sqrt{\alpha }A(x_i)U_\beta ^{}(X)(p_i+\sqrt{\alpha }A(x_i))U_\beta (X).$$
(88)
Using the form boundedness of the kinetic energy with respect to the full Hamiltonian, we have
$$(\mathrm{\Psi },\underset{j=1}{\overset{N}{}}(p_j+\sqrt{\alpha }A(x_j))^2\mathrm{\Psi })a(\mathrm{\Psi },H_m^V\mathrm{\Psi })+b(\mathrm{\Psi },\mathrm{\Psi })$$
(89)
for positive constants $`a`$ and $`b`$.. Thus, using Schwarz’ inequality it suffices to show that
$$Q_ij_\beta \mathrm{\Psi }=o(1)\mathrm{for}i\beta ,$$
(90)
as $`R`$ (the localization radius for the electrons) tends to infinity. Denote by
$$h_{i,x}^\lambda (y)=(2\pi )^{3/2}_{|k|<\mathrm{\Lambda }}\frac{1}{\omega (k)}\epsilon _i^\lambda (k)e^{ik(yx)}d^3k.$$
(91)
Explicitly, $`Q_i`$ is given by
$`p_iU_\beta ^{}(X)p_iU_\beta (X)`$ (92)
$`+`$ $`U_\beta ^{}(X)\left[{\displaystyle \underset{\lambda }{}}a_\lambda ^{}([j_11]h_x^\lambda )+a_\lambda ^{}(j_2h_x^\lambda )\right]U_\beta (X)`$
$`+`$ $`U_\beta ^{}(X)\left[{\displaystyle \underset{\lambda }{}}a_\lambda ^{}([j_11]h_x^\lambda )+a_\lambda ^{}(j_2h_x^\lambda )\right]U_\beta (X),`$
and it suffices to estimate each of these terms separately. Each of the last two terms can be brought into the form
$$U_\beta ^{}(X)a^\mathrm{\#}(f)U_\beta (X)j_\beta \mathrm{\Psi }$$
(93)
where $`f`$ is one of the functions
$$[j_1(y,\beta ,X)1]h_{1,x_j}^\lambda (y)\mathrm{or}j_2(y,\beta ,X)h_{1,x_j}^\lambda (y)j\beta .$$
(94)
The terms (93) are estimated by
$$\underset{X}{sup}\{j_\beta (X)[j_11]h_{i,x_j}^\lambda _2\}\sqrt{(\mathrm{\Psi },(𝒩+1)\mathrm{\Psi })}$$
(95)
respectively
$$\underset{X}{sup}\{j_\beta (X)j_2h_{1,x_j}^\lambda _2\}\sqrt{(\mathrm{\Psi },(𝒩+1)\mathrm{\Psi })}.$$
(96)
In both formulas the index $`j`$ is in $`\beta `$. The function $`j_\beta `$ lives in the region where $`|x_ix_j|cR`$ for $`i\beta `$ and $`j\beta ^c`$. The function $`j_11`$ (and likewise $`j_2`$) is not zero only if $`|yx_j|P`$ for some $`j\beta ^c`$. Thus, $`j_\beta (X)(j_11)(y)`$ and $`j_\beta (X)j_2(y)`$ are nonzero only if $`|yx_i|cRP`$. As $`cRP`$ gets large only the tail of the function $`h^\lambda `$ contributes to the integral which can be made as small as we please. The number operator is bounded by the field energy times $`1/m`$ which in turn is bounded by the full energy.
To estimate the first term in (92) we calculate
$$\begin{array}{cc}& p_iU_\beta ^{}(X)p_iU_\beta (X)=\hfill \\ & U_\beta ^{}(X)\underset{k}{}[a^{}([p_i,j_1]g_k)+a^{}([p_i,j_2]g_k)]\times \hfill \\ & \left[a^{}(j_1g_k)+a^{}(j_2g_k)\right]U_\beta (X).\hfill \end{array}$$
(97)
Note that the tensor product in the first line is different from the second. In the first the identity acts on $`L^2(^{3|\beta ^c|})`$ while in the second $``$ indicates the tensor product of the Fock spaces only. The functions $`g_k`$ indicates a basis of $`L^2(^3)`$. The operators $`U_\beta ^{}(X)`$ and $`U_\beta (X)`$ have unit norm. Thus
$$\begin{array}{cc}& U_\beta ^{}(X)\underset{k}{}[a^{}([p_i,j_1]g_k)+a^{}([p_i,j_2]g_k)]\times \hfill \\ & [a^{}(j_1g_k)+a^{}(j_2g_k)]U_\beta (X)\mathrm{\Psi }\hfill \\ & \left([p_i,j_1]+[p_i,j_2]\right)\sqrt{𝒩+1}\mathrm{\Psi },\hfill \end{array}$$
(98)
where $``$ indicates that the operator norm has been taken. The norms of the commutators are of the order $`1/P`$ and hence vanish as $`P\mathrm{}`$. Since the photons have a mass we can estimate the number operator in terms of the field energy.
Similar consideration apply to the $`\beta ^c`$ term in (86). The only difference is that instead of (95) and (96) we have
$$\underset{X}{sup}\{j_\beta (X)j_1h_{x_j}^\lambda _2\}\sqrt{(\mathrm{\Psi },(𝒩+1)\mathrm{\Psi })}$$
(99)
respectively
$$\underset{X}{sup}\{j_\beta (X)[j_21]h_{x_j}^\lambda _2\}\sqrt{(\mathrm{\Psi },(𝒩+1)\mathrm{\Psi })},$$
(100)
with $`j\beta ^c`$. Again this terms tend to zero as $`P\mathrm{}`$. The proof for the case where $`\beta =\{1,\mathrm{},N\}`$ is similar but simpler since the operator $`U_\beta `$ does not depend on $`X`$.
Finally, we have to compare the $`\sigma B`$ term with its localized counterparts. The estimates are similar to, but much easier than the estimates for $`(p+\sqrt{\alpha }A(x))^2`$ and are omitted for the convenience of the reader and authors who, by now, are exhausted.
A simple consequence of Lemma A.1 is the following.
###### A.2 COROLLARY.
Let $`\varphi `$ be a smooth function on $`^{3N}`$ such that $`j_\beta \varphi 0`$ for $`\beta =\{1,\mathrm{},N\}`$. Thus, $`\varphi `$ depends on $`R`$. Then, as operators,
$$\varphi H\varphi \left(\mathrm{\Sigma }^V(m,N)+o(1)\right)\varphi ^2.$$
(101)
Here, $`\mathrm{\Sigma }^V(m,N)=\mathrm{min}_{1N^{}<N}(E^V(N^{})+E^0(NN^{}))`$ and $`o(1)`$ vanishes as $`R\mathrm{}`$.
###### Proof.
By the IMS localization formula we have that
$$\varphi H\varphi =\underset{\beta }{}\varphi j_\beta Hj_\beta \varphi \varphi ^2\underset{\beta }{}|j_\beta |^2,$$
(102)
where the second term goes to 0 as $`R\mathrm{}`$. With our assumption on $`\varphi `$ only the sets $`\beta `$ with $`\beta ^c\mathrm{}`$ contribute. From Lemma A.1 we get that
$$\varphi H\varphi =\underset{\beta }{}U_\beta ^{}(X)\varphi j_\beta \left[H_\beta +H^{\beta ^c}\right]j_\beta \varphi U_\beta (X)+o(1)$$
(103)
as first $`R\mathrm{}`$ then $`P\mathrm{}`$. Certainly $`H_\beta E^V(m,|\beta |)`$ and $`H^{\beta ^c}E^0(m,|\beta ^c|)`$ from which the statement immediately follows. ∎
###### A.3 LEMMA.
Let $`\mathrm{\Psi }_n`$ be a normalized sequence in $``$ whose energy is uniformly bounded and such that for any $`\mathrm{\Phi }`$ with finite energy,
$$(\mathrm{\Psi }_n,\mathrm{\Phi })0,\mathrm{and}(\mathrm{\Psi }_n,H\mathrm{\Phi })0.$$
(104)
Then
$$(\varphi _R\mathrm{\Psi }_n,U^{}P_2U\varphi _R\mathrm{\Psi }_n)0.$$
(105)
Here $`U`$ is the Fock space localization $`U_\beta `$ that corresponds to $`\beta =\{1,\mathrm{},N\}`$.
###### Proof.
Since the energy of $`\mathrm{\Psi }_n`$ is uniformly bounded we also know that
$$(\mathrm{\Psi }_n,H^0(m)\mathrm{\Psi }_n)C.$$
(106)
is uniformly bounded
Let us describe the operator $`P_2U`$ in more detail. Recall that
$$Ua^{}(h_{i_1})\mathrm{}a^{}(h_{i_k})|0>=c^{}(h_{i_1})\mathrm{}c^{}(h_{i_k})U|0>$$
(107)
where $`|0>`$ denotes the vacuum vector in Fock space and $`U|0>=|0>|0>`$. Hence, using the definition of $`c^{}(h)`$, we find that
$$P_2Ua^{}(h_{i_1})\mathrm{}a^{}(h_{i_k})|0>=a^{}(h_{i_1}j_1)\mathrm{}a^{}(h_{i_k}j_1)|0>|0>.$$
(108)
The projection $`P_2`$ annihilates the photons in the second factor. In other words, the operator $`P_2U`$ when acting on a state
$$\mathrm{\Psi }=\{\mathrm{\Psi }^0,\mathrm{\Psi }^1(y_1),\mathrm{\Psi }^2(y_1,y_2),\mathrm{}\}$$
(109)
produces the localized state $`\mathrm{\Gamma }(j_1)\mathrm{\Psi }|0>`$ where
$$\mathrm{\Gamma }(j_1)\mathrm{\Psi }=\{\mathrm{\Psi }^0,j_1(y_1)\mathrm{\Psi }^1(y_1),j_1(y_1)j_1(y_2)\mathrm{\Psi }^2(y_1,y_2),\mathrm{}\}.$$
(110)
It follows that
$`(\varphi _R\mathrm{\Psi }_n,U^{}P_2U\varphi _R\mathrm{\Psi }_n)`$ $`=P_2U\varphi _R\mathrm{\Psi }_n^2`$
$`=\mathrm{\Gamma }(j_1)\varphi _R\mathrm{\Psi }_n^2`$
Next, we show that (106) implies that
$$\mathrm{\Gamma }(j_1)\varphi _R\mathrm{\Psi }_n0.$$
(111)
To achieve that we note first that on account of the positive mass we have that $`(\mathrm{\Psi }_n,𝒩\mathrm{\Psi }_n)`$ is uniformly bounded. Since $`\mathrm{\Psi }_n`$ is of the form
$$\{\mathrm{\Psi }_n^0,\mathrm{\Psi }_n^1(X,y_1),\mathrm{\Psi }_n^2(X,y_1,y_2),\mathrm{}\}$$
we know that $`_{kM}(\mathrm{\Psi }_n^k,\mathrm{\Psi }_n^k)\text{const}/M`$. It is therefore sufficient to prove (111) for each function
$$\mathrm{\Psi }_n^M(X,y_1,\mathrm{},y_M).$$
From the lemma below we learn that
$$\underset{j=1}{}(\mathrm{\Psi }_n,p_j^2\mathrm{\Psi }_n)$$
(112)
is uniformly bounded. Thus, we can write (111) as
$$\mathrm{\Gamma }(j_1)\varphi _R(1+\underset{j=1}{}p_j^2+H_f)^{1/2}(1+\underset{j=1}{}p_j^2+H_f)^{1/2}\mathrm{\Psi }_n.$$
(113)
which vanishes as $`n\mathrm{}`$ since $`(1+_{j=1}p_j^2+H_f)^{1/2}\mathrm{\Psi }_n`$ is uniformly bounded and since
$$\mathrm{\Gamma }(j_1)\varphi _R(1+\underset{j=1}{}p_j^2+H_f)^{1/2}$$
is compact on every finite particle subspace. Compactness follows from the fact that for continuous functions $`f`$ and $`g`$ vanishing at infinity the operator $`f(i)g(x)`$ is compact. ∎
###### A.4 LEMMA (Bound on $`A(x)^2`$).
For each $`x^3`$ and ultraviolet cutoff $`\mathrm{\Lambda }`$ write $`A(x)=D(x)+D^{}(x)`$ where $`D`$ contains the annihilation operators in $`A(x)`$ and $`D^{}`$ the creation operators. Similarly, write $`B(x)=E(x)+E^{}(x)`$. As operator bounds
$`H_f`$ $``$ $`{\displaystyle \frac{1}{8\pi \mathrm{\Lambda }}}D^{}(x)D(x)`$
$`H_f+{\displaystyle \frac{\mathrm{\Lambda }}{2}}`$ $``$ $`{\displaystyle \frac{1}{8\pi \mathrm{\Lambda }}}D(x)D^{}(x)`$
$`H_f+{\displaystyle \frac{\mathrm{\Lambda }}{8}}`$ $``$ $`{\displaystyle \frac{1}{32\pi \mathrm{\Lambda }}}A(x)^2`$
$`H_f`$ $``$ $`{\displaystyle \frac{3}{8\pi \mathrm{\Lambda }^3}}E^{}(x)E(x)`$
$`H_f+{\displaystyle \frac{3\mathrm{\Lambda }}{4}}`$ $``$ $`{\displaystyle \frac{3}{8\pi \mathrm{\Lambda }^3}}E(x)E^{}(x)`$
$`H_f+{\displaystyle \frac{3\mathrm{\Lambda }}{16}}`$ $``$ $`{\displaystyle \frac{3}{32\pi \mathrm{\Lambda }^3}}B(x)^2.`$ (114)
###### Proof.
We write $`A(x)=D(x)+D^{}(x)`$ with $`D(x)=_\lambda _{|k|<\mathrm{\Lambda }}|k|^{1/2}\epsilon _\lambda (k)\mathrm{exp}[ikx]a_\lambda ^{}(k)d^3k`$. There are thus four terms in $`A(x)^2`$ . Using the Schwarz inequality, the $`(DD)`$ term can be bounded above by $`(D^{}D)/2+(DD^{})/2`$. On the other hand, $`(DD^{})=(D^{}D)+\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }`$ is the commutator $`2/|k|=4\pi \mathrm{\Lambda }^2`$; the factor 2 comes from the two polarizations $`\lambda =1,2`$. Altogether, we obtain
$$A(x)^24D^{}(x)D(x)+4\pi \mathrm{\Lambda }^2.$$
Finally, we use the Schwarz inequality again to obtain
$$\underset{\lambda }{}\overline{h_\lambda (k)}a_\lambda ^{}(k)d^3k\underset{\lambda }{}h_\lambda (k)a_\lambda ^{}(k)d^3k\underset{\lambda }{}|h_\lambda (k)|^2/|k|d^3k\underset{\lambda }{}|k|a_\lambda ^{}(k)a_\lambda ^{}(k)d^3k.$$
In our case, $`h_\lambda (k)=\epsilon _\lambda (k)\mathrm{exp}[ikx]/\sqrt{|}k|`$, so $`_\lambda |h_\lambda (k)|^2/|k|d^3k=8\pi \mathrm{\Lambda }`$.
For $`B=\mathrm{curl}A`$, replace $`\mathrm{\Gamma }`$ by $`2\pi \mathrm{\Lambda }^4`$ and replace $`|h_\lambda (k)|`$ by $`\sqrt{|k|}`$. ∎
As a corollary of Lemma A.4 we have the following.
###### A.5 LEMMA (Bound on $`(p+A(x))^2`$).
For any $`\epsilon >0`$ there are constants $`\delta (\epsilon )>0`$ and $`C(\epsilon )<\mathrm{}`$ such that
$$\underset{j=1}{\overset{N}{}}\left\{(p_j+\sqrt{\alpha }A(x_j))^2+\frac{g}{2}\sqrt{\alpha }\sigma _jB(x_j)\right\}+\epsilon H_f\delta (\epsilon )\underset{j=1}{\overset{N}{}}p_j^2C(\epsilon ).$$
(115)
The constants $`\delta (\epsilon ),C(\epsilon )`$ depend on $`\alpha ,g,\mathrm{\Lambda },N`$.
###### Proof.
In addition to Lemma A.4, use the facts that for any $`0<\mu ,\nu <1`$, $`(p_j+\sqrt{\alpha }A(x_j))^2(1\mu )p^2+(11/\mu )\alpha A(x_j)^2`$ and $`2\sigma _jB(x_j)\nu B(x_j)^21/\nu `$. ∎
## Appendix B Appendix: VERIFICATION OF INFRARED BOUNDS
The proofs of the infrared bounds in Section 6 are somewhat formal. In particular, we carried out the calculations tacitly assuming that $`a_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m`$ (which is itself only defined for almost every $`k`$) is in the domain of the Operator $`\stackrel{~}{H}_m`$. One can actually prove this when $`\stackrel{~}{H}_m`$ is self-adjointly realized in terms of the Friedrichs’ extension and thereby make all the formal computations in Section 6 rigorous. Instead of doing so, we give here alternative proofs of the Theorems in Section 6 which avoid any reference to a domain of $`\stackrel{~}{H}_m`$. All the arguments can be carried out on the level of quadratic forms.
We recall that $`\stackrel{~}{H}_m`$ is the Hamiltonian $`H_m^V`$ after an “operator-valued gauge transformation”. Our remarks here about quadratic forms in relation to $`\stackrel{~}{H}_m`$ could just as well be applied to $`H_m^V`$ itself.
In order to keep the notation simple, we give the proof of the infrared bounds for the case of a single charged particle ($`N=1`$) with no magnetic moment, i.e., $`g=0`$. There is no difficulty in deriving these bounds for the general case.
Denote by $`𝒮`$ the set of all finite linear combinations of vectors that are products of $`C_c^{\mathrm{}}(^3)`$-functions and states in $``$ that have only a finite number of photons. It is well known that this set is dense in $``$, and that the quadratic form $`(\mathrm{\Psi },\mathrm{\Psi })_+:=(\mathrm{\Psi },(\stackrel{~}{H}_mE^V(m,1)+1)\mathrm{\Psi })`$ is defined for all $`\mathrm{\Psi }`$ in $`𝒮`$ and is bounded below by $`\mathrm{\Psi }^2`$. Hence this quadratic form is closable and the closure of $`𝒮`$ in this inner product is a Hilbert space $`Q(\stackrel{~}{H}_m)`$ with inner product $`(,)_+`$ and norm $`\mathrm{\Psi }_+=\sqrt{(\mathrm{\Psi },\mathrm{\Psi })_+}`$.
An eigenfunction $`\stackrel{~}{\mathrm{\Phi }}_m`$ of $`\stackrel{~}{H}_m`$ in the weak sense is a vector in $`Q(\stackrel{~}{H}_m)`$ such that
$$(\mathrm{\Psi },\stackrel{~}{\mathrm{\Phi }}_m)_+=e(\mathrm{\Psi },\stackrel{~}{\mathrm{\Phi }}_m)$$
(116)
for some real number $`e`$ and for all $`\mathrm{\Psi }Q(\stackrel{~}{H}_m)`$. It is in this sense that we proved in Section 4 that a ground state exists for the model with massive photons. (This implies that $`\stackrel{~}{\mathrm{\Phi }}_m`$ is in an eigenstate of the Friedrichs’ extension of $`\stackrel{~}{H}_m`$).
Define the smeared operators
$$a(f)=\underset{\lambda }{}a_\lambda (k)\overline{f(k,\lambda )}dk,$$
(117)
where $`f(k,\lambda )`$ is any function in $`L^2(^3;^2)`$. It is not difficult to show that $`a(f)\stackrel{~}{\mathrm{\Phi }}_m`$ is in the form domain of $`\stackrel{~}{H}_m`$. To this end define $`a_R(f)=R[𝒩+R]^1a(f)`$. Here $`R`$ is some large real number (which we eventually take towards infinity) and $`𝒩`$ is the number operator.
It is straightforward to see that $`a_R(f)`$ and $`a_R^{}(f)`$ are bounded operators on $`Q(\stackrel{~}{H}_m)`$ for every $`R>0`$, i.e.,
$$a_R(f)\mathrm{\Psi }_+C(R)\mathrm{\Psi }_+,$$
(118)
and similarly for $`a_R^{}(f)`$.
Generally, the constant $`C(R)`$ tends to $`\mathrm{}`$ as $`R`$ tends to $`\mathrm{}`$. For an eigenfunction of $`\stackrel{~}{H}_m`$, however, this is not the case. Simple but tedious commutator estimates reveal that for any eigenfunction $`\stackrel{~}{\mathrm{\Phi }}_m`$ there exists a constant $`C`$ independent of $`R`$ such that
$$(a_R(f)\stackrel{~}{\mathrm{\Phi }}_m,a_R(f)\stackrel{~}{\mathrm{\Phi }}_m)_+C(a_R(f)\stackrel{~}{\mathrm{\Phi }}_m,a_R(f)\stackrel{~}{\mathrm{\Phi }}_m).$$
(119)
The point is that $`(a_R(f)\stackrel{~}{\mathrm{\Phi }}_m,a_R(f)\stackrel{~}{\mathrm{\Phi }}_m)_+=(a_R^{}(f)a_R(f)\stackrel{~}{\mathrm{\Phi }}_m,\stackrel{~}{\mathrm{\Phi }}_m)_+`$ plus terms that are uniformly bounded in $`R`$. By the previous statement we know that $`a_R^{}(f)a_R(f)\stackrel{~}{\mathrm{\Phi }}_m`$ is in $`Q(\stackrel{~}{H}_m)`$ and hence
$$(a_R^{}(f)a_R(f)\stackrel{~}{\mathrm{\Phi }}_m,\stackrel{~}{\mathrm{\Phi }}_m)_+=e(a_R^{}(f)a_R(f)\stackrel{~}{\mathrm{\Phi }}_m,\stackrel{~}{\mathrm{\Phi }}_m)=e(a_R(f)\stackrel{~}{\mathrm{\Phi }}_m,a_R(f)\stackrel{~}{\mathrm{\Phi }}_m).$$
(120)
The last expression, however, is bounded uniformly in $`R`$, since the condition $`\stackrel{~}{\mathrm{\Phi }}_mQ(\stackrel{~}{H}_m)`$ implies that the expectation value of the field energy in $`\stackrel{~}{\mathrm{\Phi }}_m`$ is finite which in turn bounds the last expression in (120). Here we use the fact that the photons have a mass.
From this it follows easily that for a subsequence of $`R`$’s tending to infinity, $`a_R(f)\stackrel{~}{\mathrm{\Phi }}_m`$ has a weak limit in $`Q(\stackrel{~}{H}_m)`$. Since $`a_R(f)\stackrel{~}{\mathrm{\Phi }}_ma(f)\stackrel{~}{\mathrm{\Phi }}_m`$ strongly this shows that $`a(f)\stackrel{~}{\mathrm{\Phi }}_mQ(\stackrel{~}{H}_m)`$.
###### Proof of Theorem 6.1.
We shall use the abreviation
$$\underset{\lambda }{}\mathrm{}dk=\mathrm{\Sigma }\mathrm{}dk.$$
(121)
For our special choice of gauge
$$\stackrel{~}{A}^i(x)=a(G^i)+a^{}(G^i),i=1,2,3,$$
(122)
where we set
$$G_\lambda ^i(k,x)=\epsilon _\lambda ^i(k)|k|^{1/2}(e^{ikx}1)\chi _\mathrm{\Lambda }(k).$$
(123)
Next, pick any $`\mathrm{\Psi }`$ in $`𝒮`$ and calculate (recalling the definition of $`w`$ in Section 6 equation (48))
$$(\mathrm{\Psi },(\stackrel{~}{H}_mE^V(m,1))a(f)\stackrel{~}{\mathrm{\Phi }}_m)=2(\mathrm{\Psi },(f,G)(p+\stackrel{~}{A})\stackrel{~}{\mathrm{\Phi }}_m)(\mathrm{\Psi },a(\omega f)\stackrel{~}{\mathrm{\Phi }}_m)+i(\mathrm{\Psi },(f,\omega w)\stackrel{~}{\mathrm{\Phi }}_m)$$
(124)
with $`\omega (k)=\sqrt{k^2+m^2}`$. This extends, using an approximation argument, to all $`\mathrm{\Psi }Q(\stackrel{~}{H}_m)`$ and, in particular, to $`a(f)\stackrel{~}{\mathrm{\Phi }}_m`$. Here we note that, on account of Lemma A.4 and the assumption on the potential, $`\mathrm{\Psi }Q(\stackrel{~}{H}_m)`$ implies that $`(p+\stackrel{~}{A})\mathrm{\Psi }`$. Hence
$`0`$ $`(a(f)\stackrel{~}{\mathrm{\Phi }}_m,(\stackrel{~}{H}_mE^V(m,1))a(f)\stackrel{~}{\mathrm{\Phi }}_m)`$
$`=`$ $`2(a(f)\stackrel{~}{\mathrm{\Phi }}_m,(f,G)(p+\stackrel{~}{A})\stackrel{~}{\mathrm{\Phi }}_m)(a(f)\stackrel{~}{\mathrm{\Phi }}_m,a(\omega f)\stackrel{~}{\mathrm{\Phi }}_m)+i(a(f)\stackrel{~}{\mathrm{\Phi }}_m,(f,\omega w)\stackrel{~}{\mathrm{\Phi }}_m),`$
which yields the inequality
$$(a(f)\stackrel{~}{\mathrm{\Phi }}_m,a(\omega f)\stackrel{~}{\mathrm{\Phi }}_m)2(a(f)\stackrel{~}{\mathrm{\Phi }}_m,(f,G)(p+\stackrel{~}{A})\stackrel{~}{\mathrm{\Phi }}_m)+i(a(f)\stackrel{~}{\mathrm{\Phi }}_m,(f,\omega w)\stackrel{~}{\mathrm{\Phi }}_m)$$
(126)
for all $`f`$ in $`L^2(^3;^2)`$. Pick $`f`$ of the form $`\omega (k)^{1/2}q(k,\lambda )g_i(k,\lambda )`$ where $`g_i`$ is an orthonormal basis of $`L^2(^3;^2)`$ and $`q(k,\lambda )`$ a bounded function. Summing over this basis, we get on the left side of (126)
$$\underset{i}{}(a(\omega ^{1/2}qg_i)\stackrel{~}{\mathrm{\Phi }}_m,a(\omega ^{1/2}qg_i)\stackrel{~}{\mathrm{\Phi }}_m)=\mathrm{\Sigma }|q(k,\lambda )|^2a_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m^2dk,$$
(127)
and on the right side
$$2(a(\omega ^1|q|^2G)\stackrel{~}{\mathrm{\Phi }}_m,(p+\stackrel{~}{A})\stackrel{~}{\mathrm{\Phi }}_m)+i(a(|q|^2w)\stackrel{~}{\mathrm{\Phi }}_m,\stackrel{~}{\mathrm{\Phi }}_m).$$
(128)
Hence
$$\mathrm{\Sigma }|q(k,\lambda )|^2a_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m^2dk2(a(\omega ^1|q|^2G)\stackrel{~}{\mathrm{\Phi }}_m,(p+\stackrel{~}{A})\stackrel{~}{\mathrm{\Phi }}_m)+i(a(|q|^2w)\stackrel{~}{\mathrm{\Phi }}_m,\stackrel{~}{\mathrm{\Phi }}_m).$$
(129)
The right side can be written as
$$2\mathrm{\Sigma }\frac{|q(k,\lambda )|^2}{\omega (k)}(a_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m,G_\lambda (k)(p+\stackrel{~}{A})\stackrel{~}{\mathrm{\Phi }}_m)dk+i\mathrm{\Sigma }|q(k,\lambda )|^2(a_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m,w_\lambda \stackrel{~}{\mathrm{\Phi }}_m)dk,$$
(130)
and, applying Schwarz’s inequality, this is bounded above by
$`2`$ $`[\mathrm{\Sigma }{\displaystyle }|q(k,\lambda )|^2a_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m^2\mathrm{d}k]^{1/2}\times `$
$`\left[\left[\mathrm{\Sigma }{\displaystyle _{|k|\mathrm{\Lambda }}}\omega (k)^2|q(k,\lambda )|^2G_\lambda (p+\stackrel{~}{A})\stackrel{~}{\mathrm{\Phi }}_m^2dk\right]^{1/2}+\left[\mathrm{\Sigma }{\displaystyle _{|k|\mathrm{\Lambda }}}|q(k,\lambda )|^2w_\lambda \stackrel{~}{\mathrm{\Phi }}_m^2dk\right]^{1/2}\right].`$
Hence we obtain the bound
$`\mathrm{\Sigma }{\displaystyle |q(k,\lambda )|^2a(k)\stackrel{~}{\mathrm{\Phi }}_m^2dk}`$ (132)
$`8{\displaystyle \underset{\lambda }{}}\left[\mathrm{\Sigma }{\displaystyle _{|k|\mathrm{\Lambda }}}\omega (k)^2|q(k,\lambda )|^2G_\lambda (k)(p+\stackrel{~}{A})\stackrel{~}{\mathrm{\Phi }}_m^2dk+\mathrm{\Sigma }{\displaystyle _{|k|\mathrm{\Lambda }}}|q(k,\lambda )|^2w_\lambda \stackrel{~}{\mathrm{\Phi }}_m^2dk\right].`$
Since, $`\mathrm{div}_xG_\lambda =0`$ we have that $`G_\lambda (p+\stackrel{~}{A})\stackrel{~}{\mathrm{\Phi }}_m=(p+\stackrel{~}{A})G_\lambda \stackrel{~}{\mathrm{\Phi }}_m`$. Moreover, $`(p+\stackrel{~}{A})^2`$ is relatively form bounded with respect to $`\stackrel{~}{H}_m`$. But, as in the proof of exponential decay (Lemma 6.2), we have for each $`i=1,2,3`$
$$(G_\lambda ^i\stackrel{~}{\mathrm{\Phi }}_m,(\stackrel{~}{H}_mE^V(m,1))G_\lambda ^i\stackrel{~}{\mathrm{\Phi }}_m)=(\stackrel{~}{\mathrm{\Phi }}_m,|_xG_\lambda ^i|^2\stackrel{~}{\mathrm{\Phi }}_m)$$
(133)
and we arrive at the bound
$$\mathrm{\Sigma }|q(k,\lambda )|^2a_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m^2dkC\mathrm{\Sigma }_{|k|\mathrm{\Lambda }}\frac{|q(k,\lambda )|^2}{\omega (k)^2}\left[G_\lambda \stackrel{~}{\mathrm{\Phi }}_m^2+|_xG_\lambda |\stackrel{~}{\mathrm{\Phi }}_m^2+\omega (k)^2w_\lambda \stackrel{~}{\mathrm{\Phi }}_m^2\right]dk,$$
(134)
where $`C`$ is some constant independent of $`m`$. Since $`q(k,\lambda )`$ is arbitrary we obtain for almost every $`k`$ and each $`\lambda `$ that
$$a_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m^2C\omega (k)^2\left[G_\lambda \stackrel{~}{\mathrm{\Phi }}_m^2+|_xG_\lambda |\stackrel{~}{\mathrm{\Phi }}_m^2+\omega (k)^2w_\lambda \stackrel{~}{\mathrm{\Phi }}_m^2\right]\chi _\mathrm{\Lambda }(k).$$
(135)
The right side is bounded by
$$\frac{C}{|k|}|x|\stackrel{~}{\mathrm{\Phi }}_m^2\chi _\mathrm{\Lambda }(k),$$
(136)
which is finite on account of the exponential decay of $`\stackrel{~}{\mathrm{\Phi }}_m`$. ∎
###### Proof of Theorem 6.3.
First some notation: For any function $`f(k)`$ define
$$(\mathrm{\Delta }_hf)(k)=f(k+h)f(k),$$
(137)
and
$$\mathrm{\Delta }_ha(f)=a(\mathrm{\Delta }_hf)$$
(138)
Returning to (126) with $`f`$ replaced by $`\mathrm{\Delta }_hf`$ we have
$`(\mathrm{\Delta }_ha(f)\stackrel{~}{\mathrm{\Phi }}_m,a(\omega \mathrm{\Delta }_hf)\stackrel{~}{\mathrm{\Phi }}_m)`$
$``$ $`2(\mathrm{\Delta }_ha(f)\stackrel{~}{\mathrm{\Phi }}_m,(\mathrm{\Delta }_hf,G)(p+\stackrel{~}{A})\stackrel{~}{\mathrm{\Phi }}_m)+i(\mathrm{\Delta }_ha(f)\stackrel{~}{\mathrm{\Phi }}_m,(\mathrm{\Delta }_hf,\omega w)\stackrel{~}{\mathrm{\Phi }}_m)`$
which can be rewritten as
$`(\mathrm{\Delta }_ha(f)\stackrel{~}{\mathrm{\Phi }}_m,\mathrm{\Delta }_ha(\omega f)\stackrel{~}{\mathrm{\Phi }}_m)`$
$``$ $`(\mathrm{\Delta }_ha(f)\stackrel{~}{\mathrm{\Phi }}_m,a((\mathrm{\Delta }_h\omega )f(+h))\stackrel{~}{\mathrm{\Phi }}_m)`$
$``$ $`2(\mathrm{\Delta }_ha(f)\stackrel{~}{\mathrm{\Phi }}_m,(f,\mathrm{\Delta }_hG)(p+\stackrel{~}{A})\stackrel{~}{\mathrm{\Phi }}_m)+i(\mathrm{\Delta }_ha(f)\stackrel{~}{\mathrm{\Phi }}_m,(f,\mathrm{\Delta }_h(\omega w))\stackrel{~}{\mathrm{\Phi }}_m).`$
Notice that without the first term on the right of the inequality sign, the structure of this inequality is the same as (126), except that, of course, $`\mathrm{\Delta }_ha(f)`$ plays the role of $`a(f)`$, $`\mathrm{\Delta }_hG`$ plays the role of $`G`$ and $`\mathrm{\Delta }_h(\omega w)`$ plays the role of $`\omega w`$ . Thus, without this term we would obtain immediately the estimate analogous to (134),
$`\mathrm{\Sigma }{\displaystyle |q(k,\lambda )|^2(\mathrm{\Delta }_ha_\lambda )(k)\stackrel{~}{\mathrm{\Phi }}_m^2dk}`$ (141)
$``$ $`C\mathrm{\Sigma }{\displaystyle \frac{|q(k,\lambda )|^2}{\omega (k)^2}\left[\mathrm{\Delta }_hG_\lambda \stackrel{~}{\mathrm{\Phi }}_m^2+|_x\mathrm{\Delta }_hG_\lambda |\stackrel{~}{\mathrm{\Phi }}_m^2+\omega ^2\mathrm{\Delta }_h(\omega w_\lambda )\stackrel{~}{\mathrm{\Phi }}_m^2\right]dk}.`$
The remaining term in equation (B), after summing over the functions $`qg_i/\sqrt{\omega }`$, turns into
$$\mathrm{\Sigma }((\mathrm{\Delta }_ha_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m,a_\lambda (kh)\stackrel{~}{\mathrm{\Phi }}_m)\frac{|q(k,\lambda )|^2}{\omega (k)}(\mathrm{\Delta }_h\omega )(kh)\mathrm{d}k$$
(142)
which, by Schwarz’s inequality, is bounded above by
$$\left[\mathrm{\Sigma }|q(k,\lambda )|^2\mathrm{\Delta }_ha_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m^2dk\right]^{1/2}\left[\mathrm{\Sigma }a_\lambda (kh)\stackrel{~}{\mathrm{\Phi }}_m^2\frac{|q(k,\lambda )|^2}{\omega (k)^2}|\mathrm{\Delta }_h\omega (kh)|^2dk\right]^{1/2}.$$
(143)
This, together with (141), yields
$$\begin{array}{c}\mathrm{\Sigma }|q(k,\lambda )|^2(\mathrm{\Delta }_ha_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m^2\mathrm{d}k\hfill \\ \hfill C\mathrm{\Sigma }\frac{|q(k,\lambda )|^2}{\omega (k)^2}\left[\mathrm{\Delta }_hG_\lambda \stackrel{~}{\mathrm{\Phi }}_m^2+|_x\mathrm{\Delta }_hG_\lambda |\stackrel{~}{\mathrm{\Phi }}_m^2+\omega (k)^2\mathrm{\Delta }_h(\omega w_\lambda )\stackrel{~}{\mathrm{\Phi }}_m^2\right]dk\\ \hfill +C\mathrm{\Sigma }\frac{|q(k,\lambda )|^2}{\omega (k)^2}a_\lambda (kh)\stackrel{~}{\mathrm{\Phi }}_m^2|\mathrm{\Delta }_h\omega (kh)|^2dk.\end{array}$$
(144)
Again, since $`q`$ is arbitrary we obtain for every fixed $`\lambda `$
$$\begin{array}{c}(\mathrm{\Delta }_ha_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m^2\frac{C}{\omega (k)^2}[\mathrm{\Delta }_hG_\lambda \stackrel{~}{\mathrm{\Phi }}_m^2+|_x\mathrm{\Delta }_hG_\lambda |\stackrel{~}{\mathrm{\Phi }}_m^2\hfill \\ \hfill +\omega (k)^2\mathrm{\Delta }_h(\omega w_\lambda )\stackrel{~}{\mathrm{\Phi }}_m^2+a_\lambda (kh)\stackrel{~}{\mathrm{\Phi }}_m^2|\mathrm{\Delta }_h\omega (kh)|^2].\end{array}$$
(145)
Combining this with (136) we get
$$\begin{array}{c}(\mathrm{\Delta }_ha_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m^2\frac{C}{\omega (k)^2}[\mathrm{\Delta }_hG_\lambda \stackrel{~}{\mathrm{\Phi }}_m^2+|_x\mathrm{\Delta }_hG_\lambda |\stackrel{~}{\mathrm{\Phi }}_m^2+\omega (k)^2\mathrm{\Delta }_h(\omega w_\lambda )\stackrel{~}{\mathrm{\Phi }}_m^2]\hfill \\ \hfill +\frac{C}{\omega (k)^2|kh|}|x|\stackrel{~}{\mathrm{\Phi }}_m^2\chi _\mathrm{\Lambda }(kh)|\mathrm{\Delta }_h\omega (kh)|^2.\end{array}$$
(146)
The polarization vectors defined in (59), (60), are differentiable away from the 3-axis. The same straightforward estimates as in Section 6 lead to
$$(\mathrm{\Delta }_ha_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m^2C[\frac{1}{|k|(k_1^2+k_2^2)}+\frac{1}{|kh|\left((k_1h_1)^2+(k_2h_2)^2\right)}]|h|^2|x|\stackrel{~}{\mathrm{\Phi }}_m^2$$
(147)
which hold for all $`|k|<\mathrm{\Lambda }`$ and small $`|h|`$ with a constant $`C`$ that is independent of $`m`$.
Next, we observe that for $`k0`$ fixed, there exist a sequence of $`h`$ values, say $`h_l`$, tending to zero so that $`h^1(\mathrm{\Delta }_{he_j}a_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m`$ converges weakly to some element $`v_j(k)`$ which satisfies the estimate
$$v_j(k)^2C\frac{1}{|k|(k_1^2+k_2^2)}|x|\stackrel{~}{\mathrm{\Phi }}_m^2.$$
(148)
Here $`e_j`$ is the $`j`$-th canonical basis vector. Next we identify $`v_j(k)`$ as the weak derivative of $`a_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m`$. This weak derivative, by definition, can be computed via the expression
$$(\mathrm{\Psi },a(_j\varphi )\stackrel{~}{\mathrm{\Phi }}_m)$$
(149)
where $`\mathrm{\Psi }`$ is any state in $``$ and $`\varphi `$ is any test function in $`C_c^{\mathrm{}}(^3)`$. Clearly the above expression equals
$$\underset{h0}{lim}\mathrm{\Sigma }(\mathrm{\Psi },\mathrm{\Delta }_{he_j}a_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m)\varphi (k)dk.$$
(150)
But along the sequence $`h_l`$
$$\underset{l\mathrm{}}{lim}\mathrm{\Sigma }(\mathrm{\Psi },\mathrm{\Delta }_{h_le_j}a_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m)\varphi (k)dk=(\mathrm{\Psi },v_j(k))\varphi (k)dk,$$
(151)
which identifies $`v_j(k)`$ as the (negative) weak derivative of $`a_\lambda (k)\stackrel{~}{\mathrm{\Phi }}_m`$. ∎
marcel@math.uab.edu
lieb@princeton.edu
loss@math.gatech.edu
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# Symmetry Nonrestoration in a Gross-Neveu Model with Random Chemical Potential
## I Introduction
Intuitively, when heated, a system with initially broken symmetry will recover its symmetry because thermal fluctuations are able to overcome potential barriers. But a counterexample was noticed by Weinberg : For the four-dimensional O($`N`$)$`\times `$O($`N`$) scalar $`\varphi ^4`$ model, he showed that the system can remain in the broken phase even at sufficiently high temperature. This phenomenon is called inverse symmetry breaking or symmetry nonrestoration (SNR), depending on whether the system was in a symmetric or a broken phase at zero temperature.
Since Weinberg’s observation, SNR has been a subject of academic curiosity or a candidate way out of cosmological problems caused by topological defects like monopoles and domain walls (see Ref. for review). According to Bajc’s classification, there are three classes of SNR mechanisms in field theory: (i) a prototype case like the two-scalar model, (ii) flat directions in supersymmetric theories, and (iii) large charge density (or chemical potential). Here we restrict ourselves to class (iii).
If a large enough charge can not be stored in thermally excited modes at high temperature, it must reside in the vacuum, and this is a sign of SNR. In field theory, a scalar field (order-parameter field) gets a positive mass term by thermal effects, but a negative one by the effects due to the chemical potential. For a fixed charge (i.e., in the canonical formalism), the chemical potential is temperature dependent. In this case, if the effect of the chemical potential on the mass exceeds the thermal effects at sufficiently high temperature, the scalar field acquires a nonzero vacuum expectation value (i.e., SNR).
However, in an open system of the model which does not belong to the class (i) or (ii), the symmetry may always be restored at high temperature. In the grand canonical formalism, the chemical potential and the temperature are independent parameters and so the thermal effect on the mass always surpasses the effect of the chemical potential at sufficiently high temperature, for a fixed chemical potential. For example, consider the Gross-Neveu (GN) model with chiral symmetry. At finite chemical potential the initially broken chiral symmetry is always restored at high temperature.
In order to find a new kind of SNR in four-fermion models, we will extend the GN model at finite chemical potential to a disordered model with random chemical potential. Recently, disordered nonrelativistic Dirac fermions in two spatial dimensions have been studied in relation to the integer quantum Hall transition. Pure fermions exhibit such a transition as the value of the mass is tuned through zero, but its universality class is different from the one observed in actual experiments. Usually three types of (static) disorder are considered for a more realistic model: random gauge potential, random chemical potential and random mass.
Motivated by the SNR mechanism (iii), we introduce the (relativistic) GN model with random chemical potential in Sec. II. If the chemical potential has a Gaussian distribution at each site, our model is equivalent to the four-fermion model with two kinds of four-fermion interaction, $`(\overline{\mathrm{\Psi }}\mathrm{\Psi })^2`$ and $`(\overline{\mathrm{\Psi }}\gamma _0\mathrm{\Psi })^2`$ (see Eq.(3)), and has charge conjugation symmetry in addition to the $`Z_2`$ chiral symmetry. In Sec. III we examine the behavior of these symmetries as the temperature or disorder strength is varied using the $`1/N`$ expansion in three and two dimensions. While $`Z_2`$ chiral symmetry is always restored at high temperature, the charge conjugation symmetry exhibits SNR. In addition, we check the validity of the mean field approximation (the leading approximation in the $`1/N`$ expansion) in two dimensions. In Sec. IV the fundamental origin of SNR for charge conjugation symmetry is discussed conceptually. Our conclusions are presented in Sec. V.
## II Gross-Neveu model with random chemical potential
The Euclidean Lagrangian of the GN model at finite chemical potential $`\mu `$ is given by
$$=\overline{\mathrm{\Psi }}(\partial ̸+\mu \gamma _0)\mathrm{\Psi }\frac{g^2}{2N}(\overline{\mathrm{\Psi }}\mathrm{\Psi })^2,$$
(1)
where $`g^2(>0)`$ is the coupling constant of the four-fermion interaction $`(\overline{\mathrm{\Psi }}\mathrm{\Psi })^2`$ and $`N`$ is the number of flavors of the Dirac fermion $`\mathrm{\Psi }`$. The $`\gamma `$ matrices are $`4\times 4`$ and hermitian. Let us consider the system under the influence of a random chemical potential $`\rho (x)`$ with the Gaussian distribution $`\mathrm{exp}(d^dx\frac{N}{2R^2}\rho ^2)`$ at each site where $`R^2(>0)`$ is the strength of disorder and $`d`$ the dimension of the Euclidean space. The Gaussian noise is characterized by correlation functions
$$\rho (x)=0,\rho (x)\rho (x^{})=\frac{R^2}{N}\delta ^d(xx^{}).$$
(2)
After integrating out the random chemical potential, our model is equivalent to the four-fermion model
$$=\overline{\mathrm{\Psi }}\partial ̸\mathrm{\Psi }\frac{1}{2N}\left[g^2(\overline{\mathrm{\Psi }}\mathrm{\Psi })^2+R^2(\overline{\mathrm{\Psi }}\gamma _0\mathrm{\Psi })^2\right],$$
(3)
with the $`Z_2`$ chiral symmetry $`\{\mathrm{\Psi }\gamma _5\mathrm{\Psi },\overline{\mathrm{\Psi }}\overline{\mathrm{\Psi }}\gamma _5\}`$ and the charge conjugation symmetry $`\{\mathrm{\Psi }C\overline{\mathrm{\Psi }}^T,\overline{\mathrm{\Psi }}\mathrm{\Psi }^TC^{}\}`$. Here the matrix $`C`$ satisfies $`C^{}C=1,C^{}\gamma _\mu C=\gamma _\mu ^T`$. Under charge conjugation, $`\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ and $`\overline{\mathrm{\Psi }}\gamma _0\mathrm{\Psi }`$ transform to $`\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ and $`\overline{\mathrm{\Psi }}\gamma _0\mathrm{\Psi }`$ respectively. Hence, the Lagrangian Eq.(1) with definite chemical potential $`\mu `$ does not possess the charge conjugation symmetry (i.e., fermion-antifermion symmetry). Note that in Eq.(3) the chemical potential term does not appear explicitly.
We will study this model by the leading approximation of the $`1/N`$ expansion in three and two dimensions. To easily incorporate the $`1/N`$ expansion, let us rewrite the Lagrangian Eq.(3) by introducing scalar auxiliary fields $`\sigma (x)`$ and $`\rho (x)`$:
$$=\overline{\mathrm{\Psi }}(\partial ̸+\sigma +\rho \gamma _0)\mathrm{\Psi }+\frac{N}{2g^2}\sigma ^2+\frac{N}{2R^2}\rho ^2.$$
(4)
The random chemical potential $`\rho (x)`$ plays the role of a scalar auxiliary field. The $`Z_2`$ chiral symmetry and the charge conjugation symmetry are now expressed as $`\{\mathrm{\Psi }\gamma _5\mathrm{\Psi },\overline{\mathrm{\Psi }}\overline{\mathrm{\Psi }}\gamma _5,\sigma \sigma \}`$ and $`\{\mathrm{\Psi }C\overline{\mathrm{\Psi }}^T,\overline{\mathrm{\Psi }}\mathrm{\Psi }^TC^{},\rho \rho \}`$, respectively.
## III The behavior of $`Z_2`$ chiral symmetry and charge conjugation symmetry at zero and high temperature
For finite-temperature field theory we adopt the imaginary-time formalism. At inverse temperature $`\beta (=T^1)`$, the fermion fields are antiperiodic on $`R^{d1}\times [0,\beta ]`$, while the scalar auxiliary fields are periodic. Let us introduce the notation: $`_p^{(T)}T_{n=\mathrm{}}^{\mathrm{}}\frac{d^{d1}𝐩}{(2\pi )^{d1}},_p^{(0)}\frac{d^dp}{(2\pi )^d}`$. Integrating out the fermion fields in the partition function for Eq.(4) we obtain the effective action for the auxiliary fields $`\sigma `$ and $`\rho `$. In order to investigate the vacuum structure we need to find the finite-temperature effective potential $`V_T(\sigma ,\rho )`$ by taking $`\sigma `$ and $`\rho `$ as constant fields: To leading order in the $`1/N`$ expansion,
$$\frac{V_T(\sigma ,\rho )}{N}=\frac{\sigma ^2}{2g^2}+\frac{\rho ^2}{2R^2}2_p^{(T)}\mathrm{ln}[(p_0i\rho )^2+𝐩^2+\sigma ^2],$$
(5)
where $`p_0=(2n+1)\pi /\beta \omega _n(n=\mathrm{integer})`$ at nonzero temperature. Note that the effect of the chemical potential $`\rho `$ is to shift the energy by $`i\rho `$.
To evaluate the integration in Eq.(5), we need some mathematical formulae. By the standard method of contour integration,
$`T{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(\omega _ni\rho )^2+\sigma ^2}}`$ $`=`$ $`{\displaystyle \frac{1}{2|\sigma |}}\left[1{\displaystyle \frac{1}{1+\text{e}^{\beta (|\sigma |+|\rho |)}}}{\displaystyle \frac{1}{1+\text{e}^{\beta (|\sigma ||\rho |)}}}\right],`$ (6)
$`T{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\omega _ni\rho }{(\omega _ni\rho )^2+\sigma ^2}}`$ $`=`$ $`i\rho T{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\omega _n^2(\sigma ^2\rho ^2)}{(\omega _n^2+\sigma ^2\rho ^2)^2+(2\rho \omega _n)^2}}`$ (7)
$`=`$ $`{\displaystyle \frac{i}{2}}\left[{\displaystyle \frac{\mathrm{sinh}(\beta \rho )}{\mathrm{cosh}(\beta \sigma )+\mathrm{cosh}(\beta \rho )}}\right].`$ (8)
When the GN model is studied in the canonical formalism (i.e., with a fixed charge), similar calculations appear with imaginary chemical potential. In this case a regulating factor of the form $`\text{e}^{i\omega _n\tau }`$ is needed in evaluating the summation in Eq.(7) and ensures a finite result in the limit $`\tau 0`$ after the Matsubara sum has been performed. By using Eqs.(6) and (7), we obtain
$$T\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{ln}[(\omega _ni\rho )^2+\sigma ^2]=T\left[\mathrm{ln}2+\mathrm{ln}(\mathrm{cosh}(\beta \sigma )+\mathrm{cosh}(\beta \rho ))\right],$$
(9)
where the $`\zeta `$-function regularization was used to determine the field-independent constant. At zero temperature, these formulae reduce to
$`{\displaystyle \frac{dp_0}{2\pi }\frac{1}{(p_0i\rho )^2+\sigma ^2}}`$ $`=`$ $`{\displaystyle \frac{\theta (|\sigma ||\rho |)}{2|\sigma |}},`$ (10)
$`{\displaystyle \frac{dp_0}{2\pi }\frac{p_0i\rho }{(p_0i\rho )^2+\sigma ^2}}`$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{sgn}(\rho )\theta (|\rho ||\sigma |),`$ (11)
$`{\displaystyle \frac{dp_0}{2\pi }\mathrm{ln}[(p_0i\rho )^2+\sigma ^2]}`$ $`=`$ $`\mathrm{max}(|\sigma |,|\rho |).`$ (12)
Using Eqs.(6), (7) and $`E_\sigma \sqrt{𝐩^\mathrm{𝟐}+\sigma ^2}`$, we have
$`{\displaystyle \frac{}{\sigma }}\left({\displaystyle \frac{V_T}{N}}\right)`$ $`=`$ $`{\displaystyle \frac{\sigma }{g^2}}2\sigma {\displaystyle \frac{d^{d1}𝐩}{(2\pi )^{d1}}\frac{1}{E_\sigma }\left[1\frac{1}{1+\text{e}^{\beta (E_\sigma +\rho )}}\frac{1}{1+\text{e}^{\beta (E_\sigma \rho )}}\right]},`$ (13)
$`{\displaystyle \frac{}{\rho }}\left({\displaystyle \frac{V_T}{N}}\right)`$ $`=`$ $`{\displaystyle \frac{\rho }{R^2}}2\mathrm{sinh}(\beta \rho ){\displaystyle \frac{d^{d1}𝐩}{(2\pi )^{d1}}\frac{1}{\mathrm{cosh}(\beta E_\sigma )+\mathrm{cosh}(\beta \rho )}}.`$ (14)
To renormalize the effective potential $`V_T(\sigma ,\rho )`$ let us consider the GN model at zero temperature and in the absence of the random chemical potential because the effects of temperature and chemical potential do not change the ultraviolet behavior. Define $`1/G^21/g^21/g_c^2`$ with $`1/g_c^24\frac{d^dp}{(2\pi )^d}\frac{1}{p^2}`$. In $`2<d<4`$, the GN model is in the broken phase of the $`Z_2`$ chiral symmetry for negative $`G^2`$, corresponding to strong coupling ($`g^2>g_c^2`$) in the cutoff regularization, while it is in the symmetric phase for $`G^20`$, corresponding to weak coupling ($`0<g^2g_c^2`$). In particular, in two dimensions, the $`Z_2`$ chiral symmetry of the GN model must be broken no matter how we choose the coupling $`g^2`$. In the broken phase,
$$\frac{1}{g^2}=4\frac{d^dp}{(2\pi )^d}\frac{1}{p^2+M^2},$$
(15)
where $`M=|\sigma |(>0)`$ is the dynamically generated fermion mass at zero temperature.
From now on, we will adopt dimensional regularization, where $`G^2`$ is equal to the regularized $`g^2`$.
### A Three dimensions
In this case, renormalization is not needed to the leading order of the $`1/N`$ expansion (in dimensional regularization). By making use of Eq.(11), we can find the zero-temperature effective potential $`V_0(\sigma ,\rho )`$ directly:
$$\frac{V_0(\sigma ,\rho )}{N}=\frac{\sigma ^2}{2G^2}+\frac{\rho ^2}{2R^2}\frac{1}{6\pi }\left[\stackrel{3}{\mathrm{max}}(|\sigma |,|\rho |)3\sigma ^2\mathrm{max}(|\sigma |,|\rho |)\right],$$
(16)
where $`1/G^2=M/\pi `$ for broken $`Z_2`$ chiral symmetry. The gap equations have four kinds of solution $`(|\sigma |,|\rho |)`$: (i) $`(0,0)`$, (ii) $`(M,0)`$, (iii) $`(0,2\pi /R^2)`$ and (iv) $`(\sqrt{M(M2\pi /R^2)},M)`$. The solution (iv) exists only for $`M>2\pi /R^2`$ and corresponds to saddle points. Fig. 1 shows the zero-temperature effective potential as a function of $`\sigma /M`$ and $`\rho /M`$, for broken $`Z_2`$ chiral symmetry. $`\rho =0`$ is metastable, irrespective of the values of $`G^2`$ and $`R^2`$. For $`|\rho |>|\sigma |`$, however, $`V_0(\sigma ,\rho )`$ is unbounded from below due to the $`|\rho |^3/(6\pi )`$ term, which indicates breaking of the charge conjugation symmetry. This result stems from the fact that the term $`|\rho |^3/(6\pi )`$ arising from quantum effects surpasses the effect of the probability distribution ($`\rho ^2/(2R^2)`$) for large $`|\rho |`$.
At finite temperature, using Eqs.(6), (7) and dimensional regularization, we obtain
$`{\displaystyle _p^{(T)}}{\displaystyle \frac{1}{(\omega _ni\rho )^2+E_\sigma ^2}}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi \beta }}\left[\beta |\sigma |+\mathrm{ln}\left(1+2\text{e}^{\beta |\sigma |}\mathrm{cosh}(\beta \rho )+\text{e}^{2\beta |\sigma |}\right)\right],`$ (17)
$`{\displaystyle _p^{(T)}}{\displaystyle \frac{\omega _ni\rho }{(\omega _ni\rho )^2+E_\sigma ^2}}`$ $`=`$ $`{\displaystyle \frac{i\mathrm{sgn}(\rho )}{4\pi \beta ^2}}\left[\beta |\sigma |\mathrm{ln}\left({\displaystyle \frac{1+\text{e}^{\beta (|\sigma |+|\rho |)}}{1+\text{e}^{\beta (|\sigma ||\rho |)}}}\right)+\text{Li}_2(\text{e}^{\beta (|\sigma |+|\rho |)})\text{Li}_2(\text{e}^{\beta (|\sigma ||\rho |)})\right].`$ (18)
Here the polylogarithm $`\text{Li}_\nu (z)`$ is defined (for $`\nu >0`$) as $`\text{Li}_\nu (z)=_{k=1}^{\mathrm{}}\frac{z^k}{k^\nu }`$ (see Ref. for useful properties). From these formulae, the finite-temperature effective potential $`V_T(\sigma ,\rho )`$ is given by
$`{\displaystyle \frac{V_T(\sigma ,\rho )}{N}}`$ $`=`$ $`{\displaystyle \frac{\sigma ^2}{2G^2}}+{\displaystyle \frac{\rho ^2}{2R^2}}{\displaystyle \frac{\sigma ^3}{3\pi }}+{\displaystyle \frac{1}{\pi \beta ^3}}[\text{Li}_3(\text{e}^{\beta (\sigma +\rho )})+\text{Li}_3(\text{e}^{\beta (\sigma \rho )})`$ (20)
$`\beta \sigma \{\text{Li}_2(\text{e}^{\beta (\sigma +\rho )})+\text{Li}_2(\text{e}^{\beta (\sigma \rho )})\}],`$
up to a field-independent constant. At sufficiently high temperature,
$$\frac{V_T(\sigma ,\rho )}{N}\left(\frac{\mathrm{ln}2}{\pi }\right)T(\sigma ^2\rho ^2).$$
(21)
While the initially broken $`Z_2`$ chiral symmetry is restored at high temperature, charge conjugation symmetry is not. Hence our model exhibits nonrestoration of charge conjugation symmetry irrespective of the values of $`G^2`$ and $`R^2`$. We may interpret this phenomenon as an inverse symmetry breaking because $`\rho =0`$ is metastable at zero temperature. Intuitively, SNR is related to the tachyon-like behavior of the random chemical potential (see Eqs.(19) and (24)). In the quantum correction term of Eq.(5) the chemical potential acts as a negative mass term ($`\rho ^2`$) contrary to the usual positive mass term ($`\sigma ^2`$). From a different point of view, we will discuss the origin of SNR conceptually in Sec. IV.
### B Two dimensions
For dimensional regularization, we work in $`2+ϵ`$ dimensions. In terms of the fermion mass $`M`$, the zero-temperature effective potential is given by
$`{\displaystyle \frac{V_0(\sigma ,\rho )}{N}}`$ $`=`$ $`{\displaystyle \frac{\sigma ^2}{2\pi }}\left[1+2\mathrm{ln}\left({\displaystyle \frac{\mathrm{max}(|\sigma |,|\rho |)+\sqrt{\mathrm{max}^2(|\sigma |,|\rho |)\sigma ^2}}{M}}\right)\right]`$ (23)
$`+{\displaystyle \frac{\rho ^2}{2R^2}}{\displaystyle \frac{|\rho |}{\pi }}\sqrt{\stackrel{2}{\mathrm{max}}(|\sigma |,|\rho |)\sigma ^2},`$
where Eqs.(11) and (14) were used. Unlike in three dimensions, for large $`|\rho |`$ the effect of the probability distribution ($`\rho ^2/(2R^2)`$) is comparable to the last term ($`\rho ^2/\pi `$) in Eq.(20) arising from quantum effects. The charge conjugation symmetry can be controlled by the strength of disorder $`R^2`$. The system is in the symmetric state for $`0<R^2<\pi /2`$, while in the broken state for $`R^2>\pi /2`$. Fermions and antifermions are equally probable in the symmetric state ($`\rho =0`$), but only fermions (or antifermions) are allowed in the broken state ($`\rho =\pm \mathrm{}`$). Our system suffers from a quantum phase transition at $`R^2=\pi /2(R_c^2)`$. The gap equations have solutions $`(|\sigma |,|\rho |)`$: (i) $`(0,0)`$, (ii) $`(M,0)`$ and (iii) $`(\sqrt{(2R^2\pi )/(2R^2+\pi )}M,2R^2M/(2R^2+\pi ))`$ for $`R^2R_c^2`$, and (i) $`(0,|\rho |)`$ and (ii) $`(M,0)`$ for $`R^2=R_c^2`$. The solution (iii) exists only for $`R^2>R_c^2`$ and corresponds to saddle points. Fig. 2 shows the zero-temperature effective potential as a function of $`\sigma /M`$ and $`\rho /M`$ in (a) the symmetric and (b) the broken phase for the charge conjugation symmetry.
To examine the high-temperature ($`\beta 0`$) behavior, let us introduce dimensionless quantities: $`\stackrel{~}{V}_T=\beta ^2V_T,\stackrel{~}{\sigma }=\beta \sigma ,\stackrel{~}{\rho }=\beta \rho ,\stackrel{~}{M}=\beta M`$. We want to expand the finite-temperature effective potential in $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{\rho }`$ at high temperature (i.e., for small $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{\rho }`$). In terms of the dimensionless quantities, Eqs.(12) and (13) are reduced to
$`{\displaystyle \frac{}{\stackrel{~}{\sigma }}}\left({\displaystyle \frac{\stackrel{~}{V}_T}{N}}\right)`$ $`=`$ $`{\displaystyle \frac{2\stackrel{~}{\sigma }}{\pi }}\left[\mathrm{ln}\left({\displaystyle \frac{|\stackrel{~}{\sigma }|}{\stackrel{~}{M}}}\right)+{\displaystyle _0^{\mathrm{}}}𝑑x{\displaystyle \frac{1}{\sqrt{x^2+\stackrel{~}{\sigma }^2}}}\left({\displaystyle \frac{1}{1+\text{e}^{\sqrt{x^2+\stackrel{~}{\sigma }^2}+\stackrel{~}{\rho }}}}+{\displaystyle \frac{1}{1+\text{e}^{\sqrt{x^2+\stackrel{~}{\sigma }^2}\stackrel{~}{\rho }}}}\right)\right]`$ (24)
$`=`$ $`{\displaystyle \frac{2\stackrel{~}{\sigma }}{\pi }}\left[\left\{\mathrm{ln}\left({\displaystyle \frac{\pi }{\stackrel{~}{M}}}\right)\gamma +\mathrm{O}(\stackrel{~}{\sigma }^2)\right\}+\left\{{\displaystyle \frac{7\zeta (3)}{4\pi ^2}}+\mathrm{O}(\stackrel{~}{\sigma }^2)\right\}\stackrel{~}{\rho }^2+\mathrm{O}(\stackrel{~}{\rho }^4)\right],`$ (25)
$`{\displaystyle \frac{}{\stackrel{~}{\rho }}}\left({\displaystyle \frac{\stackrel{~}{V}_T}{N}}\right)`$ $`=`$ $`\stackrel{~}{\rho }\left[{\displaystyle \frac{1}{R^2}}{\displaystyle \frac{2\mathrm{sinh}(\stackrel{~}{\rho })}{\pi \stackrel{~}{\rho }}}{\displaystyle _0^{\mathrm{}}}𝑑x{\displaystyle \frac{1}{\mathrm{cosh}(\sqrt{x^2+\stackrel{~}{\sigma }^2})+\mathrm{cosh}(\stackrel{~}{\rho })}}\right]`$ (26)
$`=`$ $`\stackrel{~}{\rho }\left[{\displaystyle \frac{1}{R^2}}{\displaystyle \frac{2}{\pi }}+\left\{{\displaystyle \frac{7\zeta (3)}{2\pi ^3}}+\mathrm{O}(\stackrel{~}{\rho }^2)\right\}\stackrel{~}{\sigma }^2+\mathrm{O}(\stackrel{~}{\sigma }^4)\right].`$ (27)
To obtain the O($`\stackrel{~}{\rho }^0`$) term in the bracket of Eq.(21) we used the integration formula for small $`\stackrel{~}{\sigma }^2`$:
$$_0^{\mathrm{}}𝑑x\frac{1}{\sqrt{x^2+\stackrel{~}{\sigma }^2}\left(1+\text{e}^{\sqrt{x^2+\stackrel{~}{\sigma }^2}}\right)}=\frac{1}{2}\left[\mathrm{ln}\left(\frac{|\stackrel{~}{\sigma }|}{\pi }\right)+\gamma +\mathrm{O}(\stackrel{~}{\sigma }^2)\right].$$
(28)
Integrating Eqs.(21) and (22), at sufficiently high temperature, we obtain
$$\frac{V_T(\sigma ,\rho )}{N}=\frac{1}{\pi }\left[\mathrm{ln}\left(\frac{\pi T}{M}\right)\gamma \right]\sigma ^2+\frac{1}{2}\left(\frac{1}{R^2}\frac{2}{\pi }\right)\rho ^2+\frac{7\zeta (3)}{4\pi ^3T^2}\sigma ^2\rho ^2+\mathrm{},$$
(29)
up to a field-independent constant. Since the O($`\stackrel{~}{\sigma }^0`$) term in the bracket of Eq.(22) is exact, $`V_T(\sigma ,\rho )`$ has no term higher than the $`\rho ^2`$ term that consists of $`\rho `$ fields only. At high temperature, the $`Z_2`$ chiral symmetry is always restored. However, the behavior of the charge conjugation symmetry is the same as that at zero temperature. Hence, charge conjugation symmetry is not restored at high temperature, for $`R^2>R_c^2`$.
Until now in order to find the effective potential $`V_T(\sigma ,\rho )`$ we have used the mean field approximation (MFA) by taking $`\sigma `$ and $`\rho `$ as constant fields. This corresponds to the leading approximation in the $`1/N`$ expansion where the $`\sigma `$ and $`\rho `$ loops (i.e., the fluctuations of the $`\sigma `$ and $`\rho `$ fields) are not included. Let us check the validity of our calculations. In the large $`N`$ limit MFA is good, while for large but finite $`N`$ it may fail due to the contribution from kinks in two dimensions.
At first consider the case of the usual two-dimensional GN model (without the random chemical potential) where MFA predicts a wrong critical temperature $`T_0(0)`$. For $`0T<T_0`$ the MFA effective potential is the double well with two degenerate minima at $`M(T)`$ and $`M(T)`$, the solutions of the gap equation. Due to quantum tunneling between two degenerate minima, the system has kink solutions alternating between $`M(T)`$ and $`M(T)`$. They have higher energies than the constant solutions $`M(T)`$ and $`M(T)`$. The Helmholtz free energy $`F`$ is related to the internal energy $`U`$ and the entropy $`S`$ by $`F=UTS`$. Since $`F=U`$ at zero temperature, the constant solution is favored and MFA is expected to be valid. If we consider the contribution from kinks explicitly, we find the average (or blocking) potential which has a plateau between $`M(M(0))`$ and $`M`$ . Even in the presence of an arbitrarily small external field this potential is tilted to favor $`M`$ or $`M`$. Hence the contribution from kinks does not change the physics materially and MFA is qualitatively valid at zero temperature. For $`0<T<T_0`$ the number of kink configurations is sufficiently large to gain enough entropy and so their probability is overwhelming. Since the region of $`\sigma =M(T)`$ will, on the average, have the same weight as those of $`\sigma =M(T)`$, we have $`\sigma =0`$, which indicates the breakdown of MFA. For $`TT_0`$ the system is in the symmetric phase ($`M(T)=0`$) in MFA and thus has no kink solutions. Consequently, for the two-dimensional GN model MFA is good only at zero and high temperature ($`TT_0`$) and, by the formation of kinks, the true critical temperature turns out to be zero.
Now let us examine the validity of MFA for the two-dimensional GN model with the random chemical potential Eqs.(3) or (4). For the purpose of the present paper we consider only the cases of zero and sufficiently high temperature. At zero temperature the MFA effective potential $`V_0(\sigma ,\rho )`$ has degenerate minima at ($`|\sigma |,|\rho |`$): ($`M,0`$) for $`0<R^2R_c^2`$ and ($`0,\mathrm{}`$) for $`R^2>R_c^2`$ (see Fig. 2 (a) and (b)). Hence for $`0<R^2R_c^2`$ our system may have kink solutions for $`\sigma `$ alternating between $`M`$ and $`M`$, but no kinks for $`\rho `$. In this case the situation is similar to that of the usual GN model in the previous paragraph. Thus it is expected that MFA is good at zero temperature and for $`0<R^2R_c^2`$. For $`R^2>R_c^2`$ the MFA effective potential is unbounded from below and there is no tunneling between two degenerate ground states because of infinitely high barrier. So $`\sigma `$ and $`\rho `$ do not have kink solutions. At sufficiently high temperature the $`\sigma `$ and $`\rho `$ fields are decoupled from each other in $`V_T(\sigma ,\rho )`$ and can be treated separately. By the restoration of the $`Z_2`$ chiral symmetry in MFA, $`\sigma `$ has no kink solutions, irrespective of the value of the disorder strength. For $`0<R^2<R_c^2`$ the charge conjugation symmetry is preserved in MFA and so no kinks for $`\rho `$. For $`R^2>R_c^2`$ the situation is the same as that at zero temperature. For $`R^2=R_c^2`$ the MFA effective potential for $`\rho `$ vanishes. As $`R^2`$ is tuned through $`R_c^2`$, the charge conjugation symmetry undergoes a first-order phase transition from the symmetric phase ($`\rho =0`$) to the broken phase ($`\rho =\mathrm{}`$ or $`\mathrm{}`$), following positive or negative values of $`\rho `$ according to the value of $`\rho `$ in the broken phase. So it is reasonable to assume that $`\rho `$ has no kink solution at $`R^2=R_c^2`$. As a result, MFA is reliable at zero and high temperature for all values of the disorder strength.
## IV Origin of charge conjugation symmetry nonrestoration
In this section we discuss the mechanism of SNR for the charge conjugation symmetry conceptually. For convenience set $`\sigma =0`$ in Eqs.(4) and (5) because SNR is the effect of the random chemical potential. That is, we neglect the four-fermion interaction $`(\overline{\mathrm{\Psi }}\mathrm{\Psi })^2`$ and consider the system of free massless Dirac fermions in the presence of the random chemical potential. To leading order in the $`1/N`$ expansion the finite-temperature effective potential $`V_T(0,\rho )/N`$ consists of two parts: (i) a term from probability distribution ($`\rho ^2/(2R^2)P_R(\rho )`$) and (ii) the grand free energy (or grand potential) for free massless Dirac fermions at constant chemical potential $`\rho `$ ($`\mathrm{\Omega }_T(\rho )`$).
So it is essential to conceptually determine the value of the chemical potential that $`\mathrm{\Omega }_T(\rho )`$ favors. By the symmetry of $`\mathrm{\Omega }_T(\rho )`$ we can restrict ourselves to the positive chemical potential without loss of generality. At first, consider the case of zero temperature. According to Fermi-Dirac statistics, fermions fill all the energy levels to the Fermi energy (=chemical potential) and antifermions are suppressed. Hence, the larger chemical potential, the larger charge (number) density (= fermion number density$``$antifermion number density). This result is retained at nonzero temperature. Since the grand free energy density is minus the pressure, it is a decreasing function of the (positive) charge density. Consequently, the large charge density (i.e., large chemical potential) is preferred and, for all temperatures, $`\mathrm{\Omega }_T(\rho )`$ is minimized at large chemical potential ($`|\rho |\mathrm{}`$). This implies SNR for the charge conjugation symmetry in the open system of free massless Dirac fermions. Moreover, we can guess the functional form of $`\mathrm{\Omega }_T(\rho )`$ by dimensional analysis: At zero temperature, $`\mathrm{\Omega }_0(\rho )|\rho |^d`$ in $`d`$ dimensions. At sufficiently high temperature, $`\mathrm{\Omega }_T(\rho )T\rho ^2+\mathrm{O}(\rho ^4/T)`$ in three dimensions and $`\mathrm{\Omega }_T(\rho )\rho ^2+\mathrm{O}(\rho ^4/T^2)`$ in two dimensions, up to a field-independent constant. These qualitative results can be checked explicitly from Eqs.(15), (19), (20) and (24).
Now let us consider the contribution $`P_R(\rho )`$ from the probability distribution of the random chemical potential. In three dimensions, for large $`|\rho |`$, $`\mathrm{\Omega }_T(\rho )`$ exceeds $`P_R(\rho )`$ at zero and sufficiently high temperature, irrespective of the value of the disorder strength. Therefore, the initially broken charge conjugation symmetry is not restored at high temperature. In two dimensions $`P_R(\rho )`$ is comparable to $`\mathrm{\Omega }_T(\rho )`$. The probability distribution of the random chemical potential for weak disorder (small $`R^2`$) is dominated at $`\rho =0`$ and the charge conjugation symmetry is preserved at zero and high temperature. However, for strong disorder (large $`R^2`$) all values of the chemical potential have small probability density and so the fate of the charge conjugation symmetry is determined by $`\mathrm{\Omega }_T(\rho )`$. Thus, in this case, the initially broken charge conjugation symmetry is not restored at high temperature.
## V Conclusions
In the present paper we examined the symmetry behavior of the Gross-Neveu model with random chemical potential which is equivalent to the four-fermion model Eq.(3). We used the leading approximation in the $`1/N`$ expansion (i.e., the mean field approximation). Our model has the charge conjugation symmetry as well as the $`Z_2`$ chiral symmetry. The initially broken $`Z_2`$ chiral symmetry is always restored at high temperature. In three dimensions, the charge conjugation symmetry that is broken at zero temperature, is not restored at high temperature, irrespective of the value of the disorder strength $`R^2`$. In two dimensions, at zero temperature, the charge conjugation symmetry is not broken for weak disorder ($`0<R^2<R_c^2(=\pi /2)`$), but broken for strong disorder ($`R^2>R_c^2`$). Therefore, our system exhibits a quantum phase transition at $`R^2=R_c^2`$ as the value of $`R^2`$ is varied. For any given value of $`R^2`$ the high-temperature behavior of the charge conjugation symmetry is the same as its zero-temperature behavior. Hence the charge conjugation symmetry remains broken at high temperature (i.e., symmetry nonrestoration) for $`R^2>R_c^2`$. By examining the existence of the kink solutions we checked that the mean field approximation is reliable even in two dimensions at zero and high temperature.
In addition, we discussed our results on charge conjugation symmetry nonrestoration conceptually, after neglecting the four-fermion interaction $`(\overline{\mathrm{\Psi }}\mathrm{\Psi })^2`$ for convenience because symmetry nonrestoration is the effect of the random chemical potential. The behavior of the charge conjugation symmetry is determined by the competition of two terms in the finite-temperature effective potential: (i) the term $`\rho ^2/(2R^2)`$ from the Gaussian distribution for the chemical potential $`\rho `$ that favors the symmetric phase ($`\rho =0`$) and (ii) the grand free energy for free massless Dirac fermions which favors the broken phase ($`|\rho |\mathrm{}`$).
As further work, it would be worthwhile to perform next-to-leading order calculations and consider a non-Gaussian distribution for the random chemical potential.
## ACKNOWLEDGMENTS
We thank Dr. P. Vranas for helpful conversations and careful reading of the manuscript. This work was supported in part by the National Science Foundation under grant NSF-PHY96-05199. S.I.H. was also supported by the Korea Science and Engineering Foundation.
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# Obstructions to nonnegative curvature and rational homotopy theory 2000 Mathematics Subject classification. Primary 53C20, 55P62. Keywords: nonnegative curvature, soul, derivation, Halperin’s conjecture.
## 1 Introduction
According to the soul theorem of J. Cheeger and D. Gromoll \[CG72\], a complete open manifold of nonnegative sectional curvature is diffeomorphic to the total space of the normal bundle of a compact totally geodesic submanifold, called the soul.
A natural problem is to what extent the converse to the soul theorem holds. In other words, one asks which vector bundles admit complete nonnegatively curved metrics. Various aspects of this problem have been studied in \[Che73, Rig78, ÖW94, Yan95, GZ00, BK01b\]. In this paper we only deal with bundles over closed manifolds diffeomorphic to $`C\times T`$, where $`C`$ is simply-connected and $`T`$ is a standard torus. By \[CG72\] any soul has a finite cover of this form, where $`\mathrm{sec}(C)0`$.
Until recently, obstructions to the existence of metrics with $`\mathrm{sec}0`$ on vector bundles were only known for flat souls \[ÖW94\], which corresponds to the case when $`C`$ is a point. In \[BK01b\] we produced a variety of examples of vector bundles which do not admit complete metrics with $`\mathrm{sec}0`$. For instance, we showed that for any $`C`$ and $`T`$ with $`dim(T)4`$ and $`k2`$, there are infinitely many rank $`k`$ vector bundles over $`C\times T`$ whose total spaces admit no complete metrics with $`\mathrm{sec}0`$. In all these examples $`dim(T)>0`$, in fact no obstructions are known to the existence of complete metrics with $`\mathrm{sec}0`$ on vector bundles over simply-connected nonnegatively curved manifolds.
We first explain our approach to finding obstructions in case $`dim(T)>0`$. The main geometric ingredient is the splitting theorem in \[Wil00, BK01b\], which says that after passing to a finite cover, the normal bundle to the soul can be taken, by a base-preserving diffeomorphism, to the product $`\xi _C\times T`$ of a vector bundle $`\xi _C`$ over $`C`$ with the torus $`T`$ (see Theorem 3.1). Then one is faced with the purely topological problem of recognizing whether a given vector bundle over $`C\times T`$ has this property. In other words, one needs to study the orbit of $`\xi _C\times T`$ under the action of the diffeomorphism group of $`C\times T`$. Since vector bundles are rationally classified by the Euler and Pontrjagin classes, the problem reduces to analyzing the action of $`\mathrm{Diffeo}(C\times T)`$ on the rational cohomology algebra $`H^{}(C\times T,)`$ of $`C\times T`$. The “Taylor expansion” in $`T`$-coordinates of any self-diffeomorphism of $`C\times T`$ gives rise to a negative degree derivation of $`H^{}(C,)`$. One of the main points of this paper is that the orbit of $`\xi _C\times T`$ consists of bundles of the same form, unless there exists a negative degree derivation of $`H^{}(C,)`$ that does not vanish on the Euler or Pontrjagin classes of $`\xi _C`$. In particular, if $`H^{}(C,)`$ has no nonzero negative degree derivations, the above topological problem gets solved, which immediately implies that “most” bundles over $`C\times T`$ admit no complete metric of $`\mathrm{sec}0`$.
To state our main results we need the following technical definition. Given a vector bundle $`\xi `$ over $`C\times T`$, we say that $`\xi `$ virtually comes from $`C`$ if for some finite cover $`p:TT`$, the pullback of $`\xi `$ by $`\mathrm{id}_C\times p`$ is isomorphic to the product $`\xi _C\times T`$ where $`\xi _C`$ is a bundle over $`C`$.
If $`\xi `$ virtually comes from $`C`$, then no known method can rule out the existence of a complete metric with $`sec0`$ on the total space $`E(\xi )`$ of $`\xi `$, and potentially all such bundles might be nonnegatively curved.
In this paper we show that the converse is often true, namely, under various assumptions on $`C`$, we show that, if $`\xi `$ is a vector bundle over $`C\times T`$ such that $`E(\xi )`$ admits a complete metric with $`\mathrm{sec}0`$, then $`\xi `$ virtually comes from $`C`$. This happens for any $`C`$ if $`\xi `$ has rank two.
###### Theorem 1.1.
Let $`C`$ be a closed smooth simply-connected manifold, and $`T`$ be a torus. Let $`\xi `$ be a rank two vector bundle over $`C\times T`$. If $`E(\xi )`$ admits a complete metric with $`\mathrm{sec}0`$, then $`\xi `$ virtually comes from $`C`$.
Oriented $`^2`$-bundles over $`C\times T`$ are in one-to-one correspondence via the Euler class with $`H^2(C\times T,)H^2(C,)H^2(T,)`$. Thus, any oriented $`^2`$-bundle $`\xi `$ over $`C\times T`$ can be written uniquely as $`c_\xi +t_\xi `$ where $`c_\xi H^2(C,)`$, $`t_\xi H^2(T,)`$. Theorem 1.1 implies that if $`\mathrm{sec}(E(\xi ))0`$, then $`t_\xi =0`$.
More generally, it is easy to see that “most” vector bundles over $`C\times T`$ do not virtually come from $`C`$, at least when $`dim(T)`$ is large enough (for a precise result, see \[BK01b, 4.4, 4.6\] and Lemma B.1 below). In fact, $`\xi `$ virtually comes from $`C`$ iff all rational characteristic classes of $`\xi `$ lie in the $`H^{}(C,)H^0(T,)`$-term of the Künneth decomposition $`_iH^{}(C,)H^i(T,)`$ of $`H^{}(C\times T,)`$.
One of the main sources of examples of closed manifolds of nonnegative curvature is given by homogeneous spaces or, more generally, biquotients of compact Lie groups. In this case we prove
###### Theorem 1.2.
Let $`C=G//H`$ be a simply connected biquotient of compact Lie groups such that $`H`$ is semi-simple, and let $`T`$ be a torus. Let $`\xi `$ be a vector bundle over $`C\times T`$ of rank $`4`$. If $`E(\xi )`$ admits a complete metric with $`\mathrm{sec}0`$, then $`\xi `$ virtually comes from $`C`$.
Let $``$ be the class of simply-connected finite CW-complexes whose rational cohomology algebras have no nonzero derivations of negative degree.
###### Theorem 1.3.
Let $`C`$ be a closed smooth manifold, and $`T`$ be a torus. If $`\xi `$ is a vector bundle over $`C\times T`$ such that $`E(\xi )`$ admits a complete metric with $`\mathrm{sec}0`$, then $`\xi `$ virtually comes from $`C`$.
For example, $``$ contains any compact simply-connected Kähler manifold \[Mei83\], and any compact homogeneous space $`G/H`$ such that $`G`$ is a compact connected Lie group and $`H`$ is a closed subgroup with $`\mathrm{rank}(H)=\mathrm{rank}(G)`$ \[ST87\]. It was proved in \[Mar90\] that the total space of a fibration belongs to $``$ provided so do the base and the fiber.
Recall that a finite simply-connected cell complex $`C`$ is called elliptic if all but finitely many homotopy groups of $`C`$ are finite. If $`C`$ is elliptic, then $`C`$ has nonnegative Euler characteristic, and the sum of the Betti numbers of $`C`$ is $`2^m`$ where $`m`$ is the cohomological dimension of $`C`$ \[Fél89\].
Any compact simply connected homogeneous space or biquotient is elliptic, and more generally, all known closed simply connected nonnegatively curved manifolds are elliptic. In fact, it is conjectured \[GH83\] that any closed simply connected nonnegatively curved manifold is elliptic. If true, the conjecture would imply a classical conjecture of Chern-Hopf that nonnegatively curved manifolds have nonnegative Euler characteristic, and a conjecture of Gromov that the sum of the Betti numbers of a compact nonnegatively curved $`m`$-manifold is $`2^m`$.
Halperin conjectured that any elliptic space $`C`$ of positive Euler characteristic belongs to $``$. This conjecture, which is considered one of the central problems in rational homotopy theory, has been confirmed in several important cases \[FHT01, page 516\]. Note that if the above conjectures are true, then $``$ contains any simply-connected compact nonnegatively curved manifold of positive Euler characteristic.
We refer to the body of the paper for other results similar to Theorems 1.11.3. In particular, in Section 6 we establish analogs of Theorem 1.3 for $`C`$’s that belong to several classes of sphere bundles. Note that sphere bundles over closed nonnegatively curved manifolds are potentially a good source of compact manifolds with $`\mathrm{sec}0`$, because the unit sphere bundle of the normal bundle to the soul is nonnegatively curved \[GW00\]. Also in Section 8, we prove an analog of Theorem 1.3 where $`C`$ is any currently known simply connected positively curved manifold.
Not every nonnegatively curved vector bundle over $`C\times T`$ virtually comes from $`C`$, even though finding an explicit counterexample is surprisingly difficult. In fact, the conclusion of Theorem 1.2 fails already for rank six bundles over homogeneous spaces.
###### Theorem 1.4.
Let $`C=SU(6)/(SU(3)\times SU(3))`$ and $`dim(T)2`$, then there exists a rank six vector bundle $`\xi `$ over $`C\times T`$ which does not virtually come from $`C`$, but $`E(\xi )`$ admits a complete metric of $`\mathrm{sec}0`$ such that the zero section is a soul.
To prove the above theorem, we find a nonnegatively curved vector bundle $`\xi _C`$ over $`C`$ with the zero section being a soul, and a negative degree derivation $`D`$ of $`H^{}(C,)`$ that is induced by a derivation of the minimal model of $`C`$, and furthermore such that $`D`$ does not vanish on the Euler class of $`\xi _C`$, but vanishes on the Pontrjagin classes of the tangent bundle of $`C`$. Finding such $`\xi _C`$ and $`D`$ is not easy, and what makes it work here is some very special properties of the minimal model of $`C`$. Incidentally, $`SU(6)/(SU(3)\times SU(3))`$ is one of the simplest non-formal homogeneous spaces.
Now since $`D`$ is induced by a derivation of the minimal model, a multiple of $`D`$ can be “integrated” to a self-homotopy equivalence $`f`$ of $`C\times T`$. Furthermore, $`f`$ preserves the Pontrjagin classes of the tangent bundle of $`C\times T`$, because $`D`$ vanishes on the Pontrjagin classes of $`TC`$. Then by a surgery-theoretic argument, some iterated power of $`f`$ is homotopic to a diffeomorphism. Finally, since $`D`$ does not vanish on the Euler class of $`\xi _C`$, the $`f`$-pullback of $`\xi _C\times T`$ does not virtually come from $`C`$, yet it carries the pullback metric of $`\mathrm{sec}0`$ with zero section being a soul.
Structure of the paper.
Section 2 is a list of notations and conventions we use throughout the paper.
In Section 3 we introduce a purely topological property of a triple $`(C,T,k)`$, which we call splitting rigidity. As a link to nonnegative curvature, we show that if $`(C,T,k)`$ is splitting rigid, and $`\xi `$ is a rank $`k`$ vector bundle over $`C\times T`$ such that $`E(\xi )`$ admits a complete metric with $`\mathrm{sec}0`$, then $`\xi `$ virtually comes from $`C`$.
In Section 4 we relate splitting rigidity to the absence of negative degree derivations of the cohomology algebra of $`C`$. Sections 5 – 8 are devoted to applications, in particular, here we prove the results stated in Section 1, as well as splitting rigidity for certain sphere bundles and for all known positively curved manifolds.
Section 9 is an in depth study of splitting rigidity. In particular, we prove that if $`k`$ is sufficiently large, then splitting rigidity can be expressed in rational homotopy-theoretic terms. Section 10 contains the proof of Theorem 1.4.
In Section 11 we show by example that if $`k`$ is small, then changing $`C`$ within its homotopy type may turn a splitting rigid triple into a non-splitting rigid one. In Section 12 we obtains stronger obstructions to nonnegative curvature on a vector bundle under the assumption that the zero section is a soul.
In Section 13 we pose and discuss several open problems. The appendix contains a surgery-theoretic lemma and an existence result for vector bundles with prescribed Euler and Pontrjagin classes.
Much of the paper can be read without any knowledge of rational homotopy theory. In fact, rational homotopy background is only needed for Sections 78910 and 13.
Acknowledgments.
The first author is grateful to McMaster University and California Institute of Technology for support and excellent working conditions.
It is our pleasure to thank Gregory Lupton and Samuel Smith for insightful discussions on rational homotopy theory, Stefan Papadima for Lemma 5.1, Ian Hambleton for Lemma A.1, Toshihiro Yamaguchi for Example 9.7, Alexander Givental for incisive comments on deformation theory, and Burkhard Wilking and Wolfgang Ziller for countless discussions and insights related to this work. We are grateful to the referee for helpful advice on the exposition. As always, the authors are solely responsible for possible mistakes.
The present paper grew out of our earlier preprint \[BK\] written in the summer of 2000. In \[BK\] we proved much weaker results, say, Theorem 1.3 is stated there as an open question. Most of the results of the present paper were obtained in May and early June of 2001, and reported by the first author during the Oberwolfach geometry meeting on June 12, 2001. On June 29, 2001 we received a preprint by Jianzhong Pan where he independently proves Theorem 1.3 in response to our question in \[BK\]. We are very grateful to Berhard Hanke for noticing an error in the original proof of Lemma B.1 and to Berhard Hanke and Neil Strickland for their help in fixing the error.
## 2 Notations and conventions
Unless stated otherwise, all (co)homology groups have rational coefficients, all characteristic classes are over rationals, all manifolds and vector bundles are smooth.
Given a cell complex $`X`$, define $`Char(X,k)`$ to be the subspace of $`H^{}(X)`$ equal to $`_{i=1}^mH^{4i}(X)`$ if $`k=2m+1`$, and equal to $`(_{i=1}^{m1}H^{4i}(X))H^{2m}(X)`$ if $`k=2m`$. If $`\xi `$ is a (real) oriented rank $`k`$ vector bundle over $`X`$, then in case $`k`$ is odd, $`Char(X,k)`$ contains the Pontrjagin classes $`p_1(\xi ),\mathrm{},p_m(\xi )`$, and in case $`k`$ is even, $`Char(X,k)`$ contains the Pontrjagin classes $`p_1(\xi ),\mathrm{},p_{m1}(\xi )`$, and the Euler class $`e(\xi )`$. The total Pontrjagin class $`_{i0}p_i(\xi )`$ is denoted by $`p(\xi )`$.
For the product $`X\times Y`$ of pointed spaces $`X`$, $`Y`$, we denote the projections of $`X\times Y`$ onto $`X`$, $`Y`$ by $`\pi _X`$, $`\pi _Y`$. The basepoints define the inclusions of $`X`$, $`Y`$ into $`X\times Y`$ which we denote by $`i_X`$, $`i_Y`$. For a map $`f:X\times YX^{}\times Y^{}`$, we define $`f_{XX^{}}=\pi _X^{}fi_X`$, and $`f_{YY^{}}=\pi _Y^{}fi_Y`$.
For the rest of the paper, $`C`$ stands for a closed, connected, simply-connected, smooth manifold, and $`T`$ stands for a torus of some positive dimension. We use the Künneth isomorphism $`H^{}(C\times T)H^{}(C)H^{}(T)`$ to identify $`\pi _C^{}(H^{}(C))`$ with the subalgebra $`H^{}(C)1`$. We denote the total space of a vector bundle $`\xi `$ by $`E(\xi )`$.
## 3 Splitting criterion
The main geometric ingredient used in this paper is the following splitting theorem proved in \[BK01b\] (cf. \[Wil00\]).
###### Theorem 3.1.
Given a soul $`S`$ of an open complete nonnegatively curved manifold $`M`$, there is a finite cover $`p:\stackrel{~}{M}M`$, a soul $`\stackrel{~}{S}`$ of $`\stackrel{~}{M}`$ satisfying $`p(\stackrel{~}{S})=S`$, and a diffeomorphism $`f:\stackrel{~}{S}C\times T`$, where $`C`$ is a simply-connected manifold with $`\mathrm{sec}(C)0`$ and $`T`$ is a torus, such that the normal bundle to $`\stackrel{~}{S}`$ is the $`f`$-pullback of the bundle $`\xi _C\times T`$, where $`\xi _C`$ is a vector bundle over $`C`$ whose total space admits a metric of nonnegative curvature with the zero section being a soul.
Let $`\xi `$ be a vector bundle over $`C\times T`$. We say that $`\xi `$ satisfies ($``$) if
$`\begin{array}{c}E(\xi )\text{ has a finite cover diffeomorphic to the product of T and the total }\hfill \\ \text{space of a vector bundle over a closed simply-connected manifold.}\hfill \end{array}`$ $`()`$
We seek to understand how assumption ($``$) restricts $`\xi `$. In particular, we want to find conditions on $`C`$ ensuring that if $`\xi `$ satisfies ($``$), then $`\xi `$ virtually comes from $`C`$.
###### Definition 3.2.
A triple $`(C,T,k)`$, where $`k>0`$ is an integer, is called splitting rigid, if any rank $`k`$ vector bundle $`\xi `$ over $`C\times T`$ that satisfies ($``$) virtually comes from $`C`$.
By Theorem 3.1, if $`\mathrm{sec}(E(\xi ))0`$, then $`\xi `$ satisfies ($``$), so we get:
###### Proposition 3.3.
If $`(C,T,k)`$ is splitting rigid, and $`\xi `$ is a rank $`k`$ vector bundle over $`C\times T`$ such that $`E(\xi )`$ has a complete metric with $`\mathrm{sec}0`$, then $`\xi `$ virtually comes from $`C`$.
Thus if $`(C,T,k)`$ is splitting rigid, then the total spaces of “most” rank $`k`$ vector bundles over $`C\times T`$ do not admit complete metrics with $`\mathrm{sec}0`$.
As we prove in Section 9, splitting rigidity can be often expressed in rational homotopy-theoretic terms. For example, if $`kdim(C)`$, then a triple $`(C,T,k)`$ is splitting rigid if and only if, for any derivation of the minimal model of $`C`$ that commutes with the differential and has degree within $`[dim(T),0)`$, the induced derivation on $`H^{}(C)`$ vanishes on $`Char(C,k)`$.
The same statement holds with some other assumptions in place of $`kdim(C)`$, such as “$`p_i(TC)Char(C,k)`$ for all $`i>0`$”. In Section 11, we give an example when the “only if” part fails for $`k<dim(C)`$. On the other hand, the “if” part is true without any assumptions on $`k`$. As a first step towards these results, we prove the following proposition, whose weak converse is obtained in Section 9.
###### Proposition 3.4.
If any self-homotopy equivalence of $`C\times T`$ maps $`Char(C,k)1`$ to itself, then $`(C,T,k)`$ is splitting rigid.
###### Proof.
To check that $`(C,T,k)`$ is splitting rigid, we need to start with an arbitrary rank $`k`$ vector bundle $`\xi `$ over $`C\times T`$ that satisfies ($``$), and prove that $`\xi `$ virtually comes from $`C`$. Without loss of generality, we can pass to a finite cover to assume that $`E(\xi )`$ is the total space of a vector bundle $`\eta `$, which is the product of $`T`$ and a vector bundle over a closed smooth simply-connected manifold $`C^{}`$. In other words, $`\eta `$ is the $`\pi _C^{}`$-pullback of a vector bundle over $`C^{}`$.
Fix base points in $`C`$, $`C^{}`$, $`T`$ so that the inclusions $`i_C`$, $`i_C^{}`$, $`i_T`$ are defined, and let $`B=C\times T`$, $`S=C^{}\times T`$. We think of $`\xi `$ and $`\eta `$ as two vector bundle structures on a fixed manifold $`N`$, and use the zero sections to identify $`B`$ and $`S`$ with smooth submanifolds of $`N`$, and identify $`\xi `$, $`\eta `$ with the normal bundles to $`B`$ and $`S`$. Note that both $`C`$ and $`C^{}`$ are homotopy equivalent to the universal cover of $`N`$.
Let $`g:BS`$ be the homotopy equivalence induced by the zero section of $`\xi `$ followed by the projection of $`\eta `$. Note that $`\eta `$ is orientable, as the pullback of a bundle over a simply-connected manifold. Fix orientations of $`S`$ and $`\eta `$, which defines an orientation on $`E(\eta )=E(\xi )`$. We orient $`B`$ so that $`\mathrm{deg}(g)=1`$, which defines an orientation on $`\xi `$.
To simplify notations, assume that $`k=2m`$; the case of odd $`k`$ is similar. Since $`\xi `$ has rank $`k`$, $`e(\xi ),p_i(\xi )Char(B,k)`$ for $`0<i<m`$. As we remarked in \[BK01b, 4.4\], the bundle $`\xi `$ virtually comes from $`C`$ iff $`e(\xi )`$ and $`p(\xi )`$ lie in $`H^{}(C)1H^{}(C\times T)`$. Alternatively, since $`Char(B,k)(H^{}(C)1)=Char(C,k)1`$, we see that $`\xi `$ virtually comes from $`C`$ iff $`e(\xi ),p_i(\xi )Char(C,k)1`$ for $`0<i<m`$.
By Whitehead’s theorem, the maps $`g_{CC^{}}:CC^{}`$, $`g_{TT}:TT`$ are homotopy equivalences. Fix their homotopy inverses $`g_{CC^{}}^1`$, $`g_{TT}^1`$. Consider a self-homotopy equivalence $`h=(g_{CC^{}}^1\times g_{TT}^1)g`$ of $`C\times T`$. Note that each of the maps $`h_{CC}`$, $`h_{TT}`$ is homotopic to the identity because, say, $`h_{CC}`$ is equal to
$$\pi _Chi_C\pi _C(g_{CC^{}}^1\times g_{TT}^1)gi_Cg_{CC^{}}^1\pi _C^{}gi_Cg_{CC^{}}^1g_{CC^{}}\mathrm{id}_C$$
Now it is routine to check that $`g^{}`$ maps $`Char(C^{},k)1`$ into $`Char(C,k)1`$ if and only if $`h^{}`$ maps $`Char(C,k)1`$ to itself.
By assumption $`h^{}`$ maps $`Char(C,k)1`$ to itself, so $`g^{}(Char(C^{},k)1)=Char(C,k)1`$. It was observed in \[BK01b, section 3\] that $`g^{}`$ maps $`e(\eta )`$ to $`e(\xi )`$. Thus $`e(\xi )Char(C,k)1`$, as needed, and it remains to show that $`Char(C^{},k)1`$ also contains $`p_i(\xi )`$ for $`0<i<m`$.
Since $`g`$, viewed as a map $`BN`$, is homotopic to the inclusion $`BN`$, we have that $`TN|_Bg^\mathrm{\#}(TN|_S)`$. By the Whitney sum formula
$$p(TB)p(\xi )=p(TN|_B)=p(g^\mathrm{\#}(TN|_S))=p(g^\mathrm{\#}(TS\eta ))=g^{}(p(TS))g^{}(p(\eta ))$$
Since $`T`$ is parallelizable, $`p(TB)=p(TC)1H^{}(C)1`$ and $`p(TS)=p(TC^{})1H^{}(C^{})1`$. Since $`p(TB)`$ is a unit in $`H^{}(C)1`$, we can write $`p(TB)^1=_ja_j1`$ for some $`a_jH^{}(C)`$, and
$$p(\xi )=p(TB)^1g^{}(p(TS))g^{}(p(\eta ))=\left(\underset{j0}{}a_j1\right)\left(\underset{l0}{}g^{}(p_l(TC)1)\right)\left(\underset{n0}{}g^{}(p_n(\eta ))\right)$$
Also $`p_m(\eta )H^{}(C^{})1`$, because $`\eta `$ is a product of torus and a bundle over $`C^{}`$. By definition of “Char”, this means that $`p_l(TC^{})1`$ and $`p_n(\eta )`$ lie in $`Char(C^{},k)1`$ for any $`0<l,n<m`$, and therefore, $`g^{}(p_l(TC^{})1))`$, $`g^{}(p_n(\eta ))Char(C,k)1`$ for any $`0<l,n<m`$. The above formula now implies that $`p_i(\xi )Char(C,k)1`$ for $`0<i<m`$, as promised. This completes the proof that $`\xi `$ virtually comes from $`C`$. ∎
###### Remark 3.5.
It is easy to see that the assumption of Proposition 3.4 that $`Char(C,k)1`$ is invariant under any self-homotopy equivalence of $`C\times T`$ is equivalent to the formally weaker assumption that $`Char(C,k)1`$ is invariant under any self-homotopy equivalence of $`C\times T`$ satisfying $`h_{CC}\mathrm{id}_C`$, $`h_{TT}\mathrm{id}_T`$.
###### Remark 3.6.
Proposition 3.4 implies that if $`Char(C,k)1`$ is invariant under any graded algebra automorphism of $`H^{}(C\times T)`$, then $`(C,T,k)`$ is splitting rigid, and this is how we establish splitting rigidity in this paper, except for one example in Section 11 where deeper manifold topology gets involved.
###### Example 3.7.
Of course, 3.4 applies if $`Char(C,k)=0`$. Thus, $`(S^{2m+1},T,k)`$ is splitting rigid for any $`k`$, $`T`$.
In a special case $`e(\xi )H^{}(C)1`$, the very same proof of Proposition 3.4 yields the following stronger statement, in which $`ϵ^m`$ denotes the trivial rank $`m`$ bundle over $`C\times T`$.
###### Proposition 3.8.
Let $`\xi `$ be a rank $`k`$ vector bundle over $`C\times T`$ with $`e(\xi )H^{}(C)1`$ such that $`\xi ϵ^m`$ satisfies ($``$) for some $`m0`$. If any self-homotopy equivalence of $`C\times T`$ maps $`Char(C,k)1`$ to itself, then $`\xi `$ virtually comes from $`C`$.
The assumption on $`e(\xi )`$ cannot be dropped: say if $`\xi `$ is a rank two bundle over the $`2`$-torus with $`e(\xi )0`$, then $`\xi ϵ^1`$ becomes trivial in a finite cover, but $`\xi `$ does not. Of course, if $`e(\xi )=0`$ (which for example is always true if $`k`$ is odd), then $`e(\xi )H^{}(C)1`$.
## 4 Taylor expansion in cohomology
Let $`A=_pA_p`$ be a graded $``$-algebra. If $`aA_p`$, we refer to $`p`$ as the degree of $`a`$ and denote it by $`|a|`$. In this paper we only consider graded commutative algebras with an identity element $`1A_0`$, and such that $`A_p=0`$ for $`p<0`$.
Let $`B`$ be a subalgebra of $`A`$, and let $`n`$. A degree $`n`$ derivation of $`B`$ with values in $`A`$ is a linear map $`D:BA`$ such that if $`aA_p`$, then $`|D(a)|=p+n`$, and $`D(ab)=D(a)b+aD(b)(1)^{np}`$ for any $`aA_p`$, $`bA`$. We refer to $`n`$ as the degree of $`D`$, and denote it by $`|D|`$. If $`B=A`$, we just say that $`D`$ is a degree $`n`$ derivation of $`A`$.
Let $`\mathrm{Der}_n(B,A)`$ be the $``$-vector space of degree $`n`$ derivation of $`B`$ with values in $`A`$, and write $`\mathrm{Der}_n(A)`$ for $`\mathrm{Der}_n(A,A)`$. Let $`\mathrm{Der}_{}(A)=_{n<0}\mathrm{Der}_n(A)`$. We refer to derivations of $`A`$ of negative degree (i.e. to elements of $`\mathrm{Der}_{}(A)`$) as negative derivations of $`A`$.
The cohomology algebra $`H^{}(T)`$ of $`T`$ is an exterior algebra on degree one generators $`x_j`$ with $`j=1,\mathrm{},dim(T)`$. The $``$-vector space $`H^{}(T)`$ has an obvious basis $`\{t_i\}`$, $`i=0,\mathrm{},n`$ of square-free monomials in variables $`x_j`$ where $`n=2^{dim(T)}1`$. We order $`\{t_i\}`$ lexicographically so that $`t_0=1`$, $`t_j=x_j`$ for $`j=1,\mathrm{},dim(T)`$, and $`t_n=x_1\mathrm{}x_{dim(T)}`$. Thus, $`t_i^2=0`$ for $`i>0`$. We write $`H^{}(T)=_it_i`$. Then $`H^{}(C\times T)=H^{}(C)H^{}(T)=_iH^{}(C)t_i`$ is a free $`H^{}(C)`$-module.
Let $`h`$ be a self-homotopy equivalence of $`C\times T`$. Since $`H^{}(C\times T)`$ is free $`H^{}(C)`$-module with basis $`\{t_i\}`$, given $`aH^{}(C)`$, there is a unique sequence of elements $`\frac{h^{}}{t_i}(a)H^{}(C)`$ such that $`h^{}(a1)=_i(1t_i)(\frac{h^{}}{t_i}(a)1)`$. We think of $`\frac{h^{}}{t_i}`$ as $``$-linear self-maps of $`H^{}(C)`$. Informally, it is useful to interpret the above formula as a Taylor expansion of $`h^{}`$ at $`a1`$.
Since $`h^{}`$ is an algebra isomorphism, the maps $`\frac{h^{}}{t_i}`$ satisfy certain recursive identities, obtained from $`h^{}(ab)=h^{}(a)h^{}(b)`$ by collecting the terms next to $`1t_i`$’s. For example, $`\frac{h^{}}{t_0}`$ is an algebra isomorphism of $`H^{}(C)`$, and
$$\frac{h^{}}{t_1}(ab)=\frac{h^{}}{t_1}(a)\frac{h^{}}{t_0}(b)+(1)^{|a|}\frac{h^{}}{t_0}(a)\frac{h^{}}{t_1}(b).$$
(4.1)
If $`h_{CC}\mathrm{id}_C`$, which can always be arranged in our case by Remark 3.5, then $`\frac{h^{}}{t_0}(a)=a`$, and therefore, $`\frac{h^{}}{t_1}`$ is a derivation of $`H^{}(C)`$ of degree $`1`$.
Let $`d=dimT`$. Suppose all degree $`1`$ partial derivatives $`\frac{h^{}}{t_1},\mathrm{},\frac{h^{}}{t_d}`$ vanish. Then we claim that $`\frac{h^{}}{t_{d+1}}`$ is a derivation. Indeed, since $`h^{}`$ is a homomorphism we get
$$ab1+\frac{h^{}}{t_{d+1}}(ab)t_{d+1}+\text{higher order terms}=h^{}(ab1)=h^{}(a1)h^{}(b1)=$$
$$=(a1+\frac{h^{}}{t_{d+1}}(a)t_{d+1}+\text{h. o. terms})(b1+\frac{h^{}}{t_{d+1}}(b)t_{d+1}+\text{h. o. terms})$$
(4.2)
$$=ab1+[\frac{h^{}}{t_{d+1}}(a)b+a\frac{h^{}}{t_{d+1}}(b)]t_{d+1}+\text{h. o. terms}$$
which proves our assertion. Similarly, if $`\frac{h^{}}{t_i}=0`$ for $`0<i<k`$, then $`\frac{h^{}}{t_k}`$ is a derivation of degree $`|t_k|`$.
More generally, if $`h_{CC}`$ is not homotopic to $`\mathrm{id}_C`$ then $`\frac{h^{}}{t_1}(\frac{h^{}}{t_0})^1`$ is a derivation of $`H^{}(C)`$ of degree $`1`$, and if $`\frac{h^{}}{t_i}=0`$ for $`0<i<k`$, then $`\frac{h^{}}{t_k}(\frac{h^{}}{t_0})^1`$ is a derivation of $`H^{}(C)`$ of degree $`|t_k|`$.
Thus, if $`H^{}(C)`$ has no nonzero negative derivations, then $`h^{}(a1)=a1`$ for any $`aH^{}(C)`$, and any self-homotopy equivalence $`h`$ of $`C\times T`$. Thus, $`(C,T,k)`$ is splitting rigid for any $`T`$, $`k`$. Combining with Propositions 3.33.4, we deduce Theorem 1.3. We actually need the following stronger statement.
###### Proposition 4.1.
If every negative derivation of $`H^{}(C)`$ vanishes on $`Char(C,k)`$, then $`(C,T,k)`$ is splitting rigid for any $`T`$.
###### Proof.
Let $`h`$ be a self-homotopy equivalence of $`C\times T`$ with $`h_{CC}\mathrm{id}_C`$, so that $`\frac{h^{}}{t_0}(a)=a`$, and $`\frac{h^{}}{t_1}`$ is a derivation of $`H^{}(C)`$ of degree $`1`$. By Proposition 3.4, Remark 3.5, it suffices to show that $`h^{}(b1)=b1`$, for all $`bChar(C,k)`$.
Let $`\varphi _1`$ be a self-map of $`H^{}(C\times T)`$ defined by $`\varphi _1(at)=at(1t_1)(\frac{h^{}}{t_1}(a)t)`$ and for $`tH^{}(T)`$, $`aH^{}(C)`$. The fact that $`\frac{h^{}}{t_1}`$ is a derivation and $`t_1^2=0`$ implies that $`\varphi _1`$ is a homomorphism (cf. (4.2) above). It is also easy to check that the map $`atat+(1t_1)(\frac{h^{}}{t_1}(a)t)`$ is the inverse to $`\varphi _1`$ and therefore $`\varphi _1`$ is an automorphism of $`H^{}(C\times T)`$.
Then
$$\varphi _1h^{}(a1)=a1+\underset{i2}{}(1t_i)(\varphi _1^i(a)1),$$
where $`\varphi _1^i`$ are linear self-maps of $`H^{}(C)`$.
Now $`\varphi _1^2`$ is a derivation of $`H^{}(C)`$, so the formulas $`\varphi _2(a1)=a1(1t_2)(\varphi _1^2(a)1)`$ and $`\varphi _2(1t)=1t`$ for $`tH^{}(T)`$, $`aH^{}(C)`$ define an automorphism $`\varphi _2`$ of $`H^{}(C\times T)`$. Then
$$\varphi _2\varphi _1h^{}(a1)=(a1)+\underset{i3}{}(1t_i)(\varphi _2^i(a)1),$$
where $`\varphi _2^i`$ are linear self-maps of $`H^{}(C)`$, and $`\varphi _2^3`$ is a derivation. Continuing in this fashion, we get automorphisms $`\varphi _k`$ with $`\varphi _k(a1)=a1(1t_i)(\varphi _{k1}^k(a)1)`$ where $`\varphi _{k1}^k`$ is a derivation of $`H^{}(C)`$, and such that $`\varphi _n\mathrm{}\varphi _1h^{}(a1)=a1`$.
Thus $`h^{}(a1)=\varphi _1^1\mathrm{}\varphi _n^1(a1)`$. Also $`\varphi _k^1(a1)=a1+(1t_i)(\varphi _{k1}^k(a)1)`$. Now if $`bChar(C,k)`$, then by assumption $`\varphi _{k1}^k(b)=0`$ so that $`\varphi _k^1(b1)=b1`$. Thus, $`h^{}(b1)=b1`$ as desired. ∎
## 5 Splitting rigidity for rank two bundles
The following well-known lemma is the key ingredient in the proof of Theorem 1.1.
###### Lemma 5.1.
Let $`AB`$ be a finite dimensional subalgebra of a commutative graded $``$-algebra $`B`$ satisfying $`B_0`$. Let $`D`$ be a derivation of $`A`$ of degree $`2n<0`$. Then $`D`$ vanishes on $`A_{2n}`$.
###### Proof.
Let $`aA^{2n}`$. Since $`B_0`$ and $`|D|=2n`$, $`D(a)B_0`$ is a rational multiple of $`1`$. Choose a positive integer $`m`$ such that $`a^m=0`$ but $`a^{m1}0`$ (which exists since $`A`$ is finite-dimensional). Since $`|a|`$ is even, we get $`0=D(a^m)=ma^{m1}D(a)`$ so that $`D(a)=0`$. ∎
###### Proof of Theorem 1.1.
By Proposition 3.3, it is enough to show that $`(C,T,2)`$ is splitting rigid. Then by Proposition 4.1, it suffices to show that any negative derivation of $`H^{}(C)`$ vanishes on $`H^2(C)`$. Since $`C`$ is simply connected, negative derivations of degree $`2`$ automatically vanish on $`H^2(C)`$ for degree reasons. Since $`H^{}(C)`$ is finite dimensional, Lemma 5.1 implies that all derivations of degree $`2`$ vanish on $`H^2(C)`$ as well. ∎
## 6 Splitting rigidity for sphere bundles
In this section we establish splitting rigidity for various classes of sphere bundles. We need the following standard lemma.
###### Lemma 6.1.
Let $`A,B,C`$ be graded commutative $``$-algebras such that $`A`$ is finite dimensional, $`B`$ is a subalgebra of $`C`$, and the algebras $`AB`$ and $`C`$ are isomorphic as $`B`$-modules. Let $`Der(C)|_B`$ be the image of the restriction map $`Der(C)Der(B,C)`$. Then $`Der(C)|_B`$ and $`ADer(B)`$ are isomorphic as $``$-vector spaces.
###### Sketch of the proof.
Let $`\{a_i\}`$ be a basis of the vector space $`A`$. Given a derivation $`D\mathrm{Der}(C)|_B`$, define linear self-maps $`D_i^B`$ of $`B`$ by $`D(1b)=_ia_iD_i^B(b)`$. It is routine to check that $`D_i^B\mathrm{Der}(B)`$, and the correspondence $`D_ia_iD_i^B`$ gives the promised isomorphism. ∎
###### Remark 6.2.
If in the above proof $`D\mathrm{Der}_{}(C)`$, then $`D_i^B\mathrm{Der}_{}(B)`$ for all $`i`$.
###### Theorem 6.3.
Let $`T`$ be a torus, and $`C`$ be a closed, simply-connected, smooth manifold which is the total space of a sphere bundle with zero Euler class and base $`B`$ satisfying $`H^{odd}(B,)=0`$. Then if $`\xi `$ is a vector bundle over $`C\times T`$ such that $`E(\xi )`$ admits a complete metric with $`\mathrm{sec}0`$, then $`\xi `$ virtually comes from $`C`$.
###### Proof.
If $`l`$ is even, both the fiber and the base of the $`S^l`$-fibration $`p:CB`$ belong to $``$, and hence $`C`$ by \[Mar90\]. Thus, we are done by Theorem 1.3.
Suppose that $`l`$ is odd. Since the Euler class of $`CB`$ is trivial, the Serre spectral sequence of the fibration collapses at the $`E_2`$-term, so $`H^{}(C)`$ and $`H^{}(B)H^{}(S^l)`$ are isomorphic, as $`H^{}(B)`$-modules. Since $`H^{odd}(B)=0`$, this implies that $`p:H^{even}(B)H^{even}(C)`$ is an isomorphism. Note that under the above identification of $`H^{}(C)`$ and $`H^{}(B)H^{}(S^l)`$, $`p^{}(H^{}(B))=H^{ev}(C)`$ corresponds to $`H^{}(B)1`$.
By Proposition 4.1, to prove splitting rigidity of $`(C,T,k)`$ for any $`T,k`$, it is enough to show that all negative derivations of $`H^{}(C)`$ vanish on $`H^{ev}(C)`$. Let $`D\mathrm{Der}_{}(H^{}(C))`$ be a negative derivation. By Lemma 6.1 we can write $`D|_{H^{ev}(C)}`$ as $`aD_1+D_2`$ where $`aH^l(S^l)`$ and $`D_1,D_2\mathrm{Der}_{}(H^{}(B))`$. Now $`B`$ implies that $`D_1=D_2=0`$, and hence $`D|_{H^{ev}(C)}=0`$, as needed. ∎
In case $`dim(S^l)dim(B)`$, the Euler class is automatically zero, therefore we have
###### Corollary 6.4.
Let $`S^lCB`$ be a sphere bundle over $`B`$ with $`H^{odd}(B)=0`$ where $`ldim(B)`$. Then $`(C,T,k)`$ is splitting rigid for any $`T,k`$.
###### Remark 6.5.
If $`B`$ is elliptic, then $`H^{odd}(B,)=0`$ is equivalent to $`\chi (B)>0`$ \[FHT01, Proposition 32.10\], where $`\chi `$ is the Euler characteristic. In particular, all homogeneous spaces $`G/H`$ with $`\mathrm{rank}(H)=\mathrm{rank}(G)`$ have $`H^{odd}(G/H,)=0`$. As we mentioned in the introduction, conjecturally, any nonnegatively curved manifold $`B`$ of positive Euler characteristic has $`H^{odd}(B,)=0`$ and lies in $``$.
Let $`(n)`$ be the class of simply-connected finite CW-complexes whose rational cohomology algebras are generated in dimension $`n`$. By Lemma 5.1, $`(2n)`$ for any $`n`$. Notice that this implies that any $`B(2n)`$ satisfies the assumptions of Theorem 6.3. This also implies that the class $`(2n)`$ is closed under fibrations because if $`FEB`$ is a fibration with the fiber in $``$, then by \[Mar90\] there is an $`H^{}(B)`$-module isomorphism $`H^{}(E)H^{}(F)H^{}(B)`$. A simple Mayer-Vietoris arguments implies that if closed simply-connected manifolds $`C`$, $`C^{}`$ belong to $`(2n)`$, then so does the connected sum $`C\mathrm{\#}C^{}`$.
###### Example 6.6.
Besides being closed under fibrations and connected sums, the class $`(2n)`$ contains the following nonnegatively curved manifolds:
$``$ $`(2)`$ contains $`S^2`$, $`CP^n`$, $`CP^n\mathrm{\#}CP^n`$, $`CP^n\mathrm{\#}\overline{CP^n}`$ \[Che73\], all nontrivial $`S^2`$-bundles over $`S^4`$ \[GZ00\], biquotients $`G//T`$ of a compact connected Lie group $`G`$ by a maximal torus $`T`$ \[Sin93\], $`SO(2n+1)/(SO(2n1)\times SO(2))`$ which is a homology $`CP^{2n1}`$ \[MZ87\], projectivized tangent bundle to $`CP^n`$ \[Wil02\], and the exceptional space $`G_2/U(2)`$ which is a homology $`CP^5`$ \[MZ87\];
$``$ $`(4)`$ contains $`S^4`$, $`HP^n`$, $`HP^n\mathrm{\#}HP^n`$, $`HP^n\mathrm{\#}\overline{HP^n}`$ \[Che73\], $`G_2/SO(4)`$ which is a homology $`HP^2`$, projectivized tangent bundle to $`HP^n`$ \[Wil02\], and $`Sp(n)/K`$ where $`K`$ is the product of $`n`$ copies of $`Sp(1)`$;
$``$ $`(8)`$ contains $`S^8`$, $`F_4/Spin(9)`$, and $`F_4/Spin(8)`$.
###### Theorem 6.7.
If $`S^lCB`$ is a sphere bundle such that $`B(2n)`$ for some $`n`$, then $`(C,T,k)`$ is splitting rigid for any $`T,k`$.
###### Proof.
If $`l`$ is even, then $`S^l,B`$ so that $`C`$ by \[Mar90\], and we are done by Theorem 1.3. If the Euler class of $`CB`$ vanishes, then the result follows by Theorem 6.3. Thus, we can assume that $`l`$ is odd and $`p:CB`$ has nonzero Euler class.
It follows from the Gysin sequence that $`p^{}:H^{ev}(B)H^{ev}(C)`$ is onto, so $`H^{ev}(C)`$ is generated in dimension $`2n`$. Again, we see from the Gysin sequence that $`H^i(C)=0`$ for $`0<i<2n`$. Hence, Lemma 5.1 implies that $`\mathrm{Der}_{}(H^{ev}(C),H^{}(C))=0`$. So by proposition 4.1, $`(C,T,k)`$ is splitting rigid for any $`T,k`$. ∎
## 7 Splitting rigidity for biquotients
Let $`G`$ be a compact Lie group and $`HG\times G`$ be a compact subgroup. Then $`H`$ acts on $`G`$ on the left by the formula $`(h_1,h_2)g=h_1gh_2^1`$. The orbit space of this action is called a biquotient of $`G`$ by $`H`$ and denoted by $`G//H`$. If the action of $`H`$ on $`G`$ is free, then $`G//H`$ is a manifold. This is the only case we consider in this paper. In the special case when $`H`$ has the form $`K_1\times K_2`$ where $`K_1G\times 1G\times G`$ and $`K_21\times GG\times G`$ we will sometimes write $`K_1\backslash G/K_2`$ instead of $`G//(K_1\times K_2)`$. For any biinvariant metric on $`G`$ the above action of $`H`$ is isometric, and therefore, $`G//H`$ can be equipped with a submersion metric, which by O’Neill’s submersion formula is nonnegatively curved. Thus, biquotients form a large class of examples of nonnegatively curved manifolds.
As was observed by Eschenburg \[Esc92a\], any biquotient $`G//H`$ is diffeomorphic to a biquotient of $`G\times G`$ by $`G\times H`$ written as $`\mathrm{\Delta }G\backslash G\times G/H`$, where $`\mathrm{\Delta }G`$ stands for the diagonal embedding of $`G`$ into $`G\times G`$. Let $`p:G//HB_H`$ be the classifying map of the principle $`H`$ bundle $`HGG//H`$. Then it is easy to see (cf. \[Esc92a\]) that $`GG//HBH`$ is a Serre fibration (which need not be principal!). Moreover, this fibration fits into the following fibered square (see \[Esc92a\] and \[Sin93\])
(7.3)
where both vertical arrows are fibrations with fiber $`G`$ and both horizontal arrows are fibrations with fiber $`(G\times G)/H`$. In particular, the fibration $`G//HB_H`$ is the pullback of the fibration $`GB_GB_{G\times G}`$. Following Eschenburg, we call the fibration $`B_GB_{G\times G}`$ the reference fibration.
Next we are going to construct a Sullivan model of the biquotient $`G//H`$.
We refer to \[TO97, Chapter1\] for a gentle introduction to rational homotopy theory, and use the textbook \[FHT01\] as a comprehensive reference.
Recall that a free DGA $`(\mathrm{\Lambda }V,d)`$ is called pure if $`V`$ is finite-dimensional and $`d|_{V^{ev}}=0`$ and $`d(V^{odd})V^{ev}`$. It is well-known that homogeneous spaces admit natural pure Sullivan models given by their Cartan algebras. The next proposition shows that the same remains true for biquotients.
###### Proposition 7.1.
Let $`G//H`$ be a biquotient. Then it admits a pure Sullivan model.
###### Proof.
We begin by constructing the canonical model of the reference fibration $`\varphi :B_GB_{G\times G}=B_G\times B_G`$. Since this fibration is induced by the diagonal map $`\mathrm{\Delta }:GG\times G`$, it follows that $`\varphi `$ is the diagonal embedding $`\mathrm{\Delta }_{B_G}:B_GB_G\times B_G`$. Consider the map $`\varphi ^{}:H^{}(B_G\times B_G)H^{}(B_G)`$. It is well-known that $`G`$ is rationally homotopy equivalent to $`S^{2m_11}\times \mathrm{}\times S^{2m_n1}`$ and the minimal model of $`B_G`$ is isomorphic to $`H^{}(BG,)[x_1,\mathrm{},x_n]`$ with zero differentials and with $`|x_i|=2m_i`$. Similarly the minimal model of $`B_G\times B_G`$ is isomorphic to its cohomology ring $`B=[x_1,\mathrm{},x_n,y_1,\mathrm{},y_n]`$ with $`|x_i|=|y_i|=2m_i`$. Thus $`\varphi ^{}`$ can be viewed as a DGA-homomorphism of minimal models of $`B_G`$ and $`B_{G\times G}`$.
Let us construct a Sullivan model of $`\varphi ^{}`$. Since $`\varphi =\mathrm{\Delta }_{B_G}`$ we compute that $`\varphi ^{}(x_i)=\varphi ^{}(y_i)=x_i`$ for all $`i=1,\mathrm{},n`$. Consider the relative Sullivan algebra $`(B\mathrm{\Lambda }(q_1,\mathrm{}q_n),d)`$ where $`dx_i=dy_i=0`$ and $`dq_i=x_iy_i`$. Then it is immediate to check that this relative algebra is a Sullivan model (in fact, a minimal one) of $`\varphi ^{}`$ with the quasi-isomorphism $`B\mathrm{\Lambda }(q_1,\mathrm{}q_n)H^{}(B_G)`$ given by $`x_ix_i`$, $`y_ix_i`$, $`q_i0`$.
By the naturality of models of maps \[FHT01, page 204, Proposition 15.8\], from the fibered square (7.3), we obtain that a Sullivan model of the map $`G//HB_H`$ can be given by the pushout of $`(B\mathrm{\Lambda }(q_1,\mathrm{}q_n),d)`$ via the homomorphism $`f^{}:BH^{}(B_H)`$; i.e. it can be written as
$$(H^{}(B_H),0)_{(B,d)}(B\mathrm{\Lambda }(q_1,\mathrm{}q_n),d)=(H^{}(B_H)\mathrm{\Lambda }(q_1,\mathrm{}q_n),\overline{d})$$
where $`\overline{d}`$ is given by $`\overline{d}|_{H^{}(B_H)}=0`$ and $`\overline{d}(q_i)=f^{}(x_iy_i)`$. In particular, $`M(G//H)=(H^{}(B_H)\mathrm{\Lambda }(q_1,\mathrm{}q_n),\overline{d})`$ is a model for $`G//H`$. Notice that $`H^{}(B_H)`$ is a free polynomial algebra on a finite number of even-dimensional generators, and thus the model $`M(G//H)`$ is pure. ∎
###### Remark 7.2.
It is easy to see that the minimal model of a pure Sullivan model is again pure. Therefore, Proposition 7.1 implies that minimal models of biquotients are pure.
###### Remark 7.3.
The pure model $`M(G//H)`$ constructed in the proof of Proposition 7.1 provides an effective way of computing rational cohomology of biquotients. We refer to $`M(G//H)`$ as the Cartan model of $`G//H`$. This method of computing $`H^{}(G//H)`$ is essentially equivalent to the method developed by Eschenburg \[Esc92b\] who computed the Serre spectral sequence of the fibration $`GG//HB_H`$. In fact, it is easy to recover this spectral sequence by introducing the standard bigrading on the Cartan model $`M(G//H)`$. Also note that in case when $`G//H`$ is an ordinary homogeneous space (i.e. when $`HG\times G`$ has the form $`H\times 1G\times G`$) this model is easily seen to reduce to the standard Cartan model of $`G/H`$.
###### Proof of Theorem 1.2.
By Theorem 1.1, we only have to consider the case when $`k=3`$ or $`4`$. By Proposition 3.3 it is enough to show that $`(C,T,k)`$ is splitting rigid for any $`T`$ and $`k=3`$ or $`4`$. Since $`Char(C,3)=Char(C,4)=H^4(C)`$ and according to Proposition 4.1, to insure splitting rigidity it is enough to show that all negative derivations of $`H^{}(G//H)`$ vanish on $`H^4(G//H)`$.
By passing to a finite cover we can assume that both $`G`$ and $`H`$ are connected. Since $`H`$ is semisimple and $`G//H`$ is simply connected, the long exact sequence of the fibration $`HGG//H`$ implies that $`G`$ is also semisimple.
Let $`(\mathrm{\Lambda }(V),d)`$ be the minimal model of $`G//H`$. Since $`G//H`$ is $`2`$-connected, we have that $`V^1=V^2=0`$. According to Proposition 7.1, the model $`(\mathrm{\Lambda }(V),d)`$ is a pure. By minimality of $`(\mathrm{\Lambda }(V),d)`$, we have that $`d|_{V_3}=0`$ and $`V^3H^3(G//H)`$. By the structure theorem for pure DGAs, \[Oni94, p 141, Prop 3\], this implies that $`\mathrm{\Lambda }(V)\mathrm{\Lambda }V^3\mathrm{\Lambda }(\widehat{V})`$ and $`H^{}(\mathrm{\Lambda }(V))\mathrm{\Lambda }V^3H^{}(\mathrm{\Lambda }(\widehat{V}))`$ for some differential subalgebra $`\mathrm{\Lambda }(\widehat{V})\mathrm{\Lambda }(V)`$ such that $`\widehat{V}^3=0`$. Let $`A=\mathrm{\Lambda }V^3`$ and $`B=H^{}(\mathrm{\Lambda }(\widehat{V}))`$. By above $`B^1=B^2=B^3=0`$. Therefore, by Lemma 5.1, all negative derivations of $`B`$ vanish on $`B^4`$. Notice that $`H^4(G//H)`$ corresponds to $`1B^4`$ under the isomorphism $`H^{}(G//H)AB`$, and hence applying Lemma 6.1, we conclude that negative derivations of $`H^{}(G//H)`$ vanish on $`H^4(G//H)`$. ∎
###### Remark 7.4.
Using Proposition 3.8, we get a stronger version of Theorem 1.2. Namely, if $`\xi `$ is a vector bundle over $`C\times T`$ of rank $`3`$ or $`4`$ with $`e(\xi )H^{}(C)1`$ such that for some $`m0`$, the manifold $`E(\xi )\times ^m`$ admits a complete metric with $`\mathrm{sec}0`$, then $`\xi `$ virtually comes from $`C`$. If $`\xi `$ has rank $`3`$, then $`e(\xi )=0`$, so the assumption $`e(\xi )H^{}(C)1`$ is automatically true.
###### Remark 7.5.
It would be interesting to see whether Theorem 1.2 remains true if $`H`$ is not assumed to be semi-simple. In that case a slight modification of the proof of Theorem 1.2 still shows that derivations of degree $`4,3`$ and $`1`$ vanish on $`H^4(G//H)`$, and therefore $`(G//H,S^1,k)`$ is splitting rigid for $`k4`$. However, as of this writing, we are unable to see whether degree $`2`$ derivations have to vanish, and thus the general case remains unclear.
###### Proposition 7.6.
Let $`C=G//H`$ be a simply connected biquotient of a compact group $`G`$ by a torus $`H`$ satisfying $`\mathrm{rank}(H)=\mathrm{rank}(G)1`$. Then $`(C,T,k)`$ is splitting rigid for any $`T,k`$.
###### Proof.
First, we show that $`H^{even}(G//H)`$ is generated in dimension $`2`$. Consider the Cartan model $`M(G//H)=(H^{}(B_H)\mathrm{\Lambda }(q_1,\mathrm{}q_n),\overline{d})`$. It admits a natural grading by the wordlength in $`q_i`$’s given by $`M(G//H)_k=H^{}(B_H)\mathrm{\Lambda }^k(q_1,\mathrm{}q_n)`$. Since the differential decreases the wordlength in $`q_i`$’s by $`1`$, this grading induces a natural grading in the cohomology $`H^{}(G//H)=H_k^{}(G//H)`$; this is the so-called lower grading on $`H^{}(G//H)`$.
According to \[Hal77, Theorem 2\], (also cf. \[Sin93, Proposition 6.4\]), $`H_k^{}(G//H)=0`$ for $`k>\mathrm{rank}(G)\mathrm{rank}(H)`$. In our case $`\mathrm{rank}(G)\mathrm{rank}(H)=1`$, and hence $`H_k^{}(G//H)=0`$ for $`k>1`$. Next observe that $`H_1^{ev}(G//H)=0`$ since $`|q_i|`$ is odd for any $`i`$ and $`H^{odd}(B_H)=0`$. Similarly $`H_0^{odd}(G//H)=0`$. Therefore, $`H^{ev}(G//H)=H_0^{}(G//H)`$ which is equal to the quotient of $`H^{}(B_H)`$ by $`\overline{d}(\mathrm{\Lambda }(q_1,\mathrm{}q_n))`$. Since $`H`$ is a torus, $`H^{}(B_H)`$ is generated by $`2`$-dimensional classes and hence by above the same is true for $`H^{ev}(G//H)`$.
By Lemma 5.1 this implies that $`Der_{}(H^{ev}(G//H),H^{}(G//H))=0`$ which by the splitting criterion 4.1 means that $`(G//H,T,k)`$ is splitting rigid for any $`T,k`$. ∎
## 8 Splitting rigidity for known positively curved manifolds
###### Proposition 8.1.
Let $`C`$ be a known closed simply-connected positively curved manifold, and $`T`$ be a torus. If $`\xi `$ is a vector bundle over $`C\times T`$ such that $`E(\xi )`$ admits a complete metric with $`\mathrm{sec}0`$, then $`\xi `$ virtually comes from $`C`$.
###### Proof.
First of all, note that all known even dimensional positively curved manifolds belong to $``$. Indeed, it follows from the classification theorem of positively curved homogeneous spaces \[Wal72\] that any even dimensional homogeneous space belongs to $`(2n)`$ for some $`n>0`$. The only known example of a positively curved even-dimensional manifold which is nondiffeomorphic to a positively curved homogeneous space is the space $`M^6=SU(3)//T^2`$ \[Esc92b\]. This space is an $`S^2`$-bundle over $`CP^2`$ and, therefore, it lies in $``$.
Now we establish splitting rigidity for all known odd-dimensional positively curved manifolds. Those are the standard spheres, the Berger $`7`$-dimensional homology sphere $`B^7=Sp(2)/Sp(1)_{\text{max}}`$, the Eschenburg $`7`$-manifolds $`E_{k,l,m,n}^7`$ \[Esc82\] all obtained as biquotients of $`SU(3)`$ by $`S^1`$, and the Bazaikin $`13`$-manifolds $`B_{k,l,m,n}^{13}`$ \[Baz96\] obtained as biquotients of $`SU(5)`$ by $`Sp(2)\times S^1`$.
By 3.7, any odd-dimensional rational homology sphere is splitting rigid, so it remains to deal with the Eschenburg and Bazaikin manifolds. By Lemma 8.2 below, all Eschenburg manifolds are rationally homotopy equivalent to $`S^2\times S^5`$, and all Bazaikin manifolds are rationally homotopy equivalent to $`CP^2\times S^9`$. Since both $`CP^2`$ and $`S^2`$ belong to $`(2)`$, the proof of Theorem 6.3 implies that any negative derivation of $`H^{}(S^2\times S^5)`$ or $`H^{}(CP^2\times S^9)`$ vanishes on even cohomology. Now Proposition 4.1 implies splitting rigidity of all Bazaikin and Eschenburg manifolds. ∎
###### Lemma 8.2.
All Eschenburg manifolds are rationally homotopy equivalent to $`S^2\times S^5`$; all Bazaikin manifolds are rationally homotopy equivalent to $`CP^2\times S^9`$
###### Proof.
We only give a proof for the Bazaikin manifolds; the Eschenburg manifolds are treated similarly. Let $`B^{13}`$ be a Bazaikin manifold. From the homotopy sequence of the fibration
$$Sp(2)\times S^1SU(5)B^{13}$$
one easily sees that $`B^{13}`$ has the same rational homotopy groups as $`CP^2\times S^9`$. Let $`M=(\mathrm{\Lambda }V,d)`$ be the minimal model of $`B^{13}`$. It is well-known (e.g. see \[FHT01, Theorem 15.11\]) that $`V\mathrm{Hom}_{}(\pi _{}(B),)`$. Therefore, $`M`$ is a free graded algebra on generators $`x_2,y_5,y_9`$ with the degrees of the generators given by the subscripts. Obviously, $`d(x_2)=0`$ and $`d(y_5)=kx_2^3`$ for some rational $`k`$. Note that $`k0`$, else we would have $`H^5(M)`$ which is known not to be the case by \[Baz96\]. Replacing $`y_5`$ with $`y_5/k`$, we can assume that $`k=1`$. It is clear that $`dy_9`$ must be equal to $`lx_2^5`$ for some rational $`l`$. We claim that any such minimal model is isomorphic to the one with $`l=0`$. Indeed, let $`M_l=\mathrm{\Lambda }(x_2,y_5,y_9)|dx_2=0,dy_5=x_2^3,dy_9=lx_2^5`$. Consider the map $`M_0M_l`$ given by $`x_2x_2,y_5y_5,y_9y_9ly_5x_2^2`$. This map is easily seen to be a DGA-isomorphism with the inverse given by $`x_2x_2,y_5y_5,y_9y_9+ly_5x_2^2`$. Since $`M_0`$ is a minimal model of $`CP^2\times S^9`$, the proof is complete. ∎
###### Remark 8.3.
As we explained above, any known closed (simply-connected) even-dimensional positively curved manifold $`C`$ belongs to $`(2n)`$ for some $`n`$. The same is true for all known \[Wil02\] even-dimensional manifolds with sectional curvature positive on an open dense subset, such as projectivized tangent bundles of $`HP^n`$ and $`CP^n`$. Therefore, all these examples are also splitting rigid for any $`T,k`$.
## 9 Splitting rigidity and derivations in minimal models
In this section, we study splitting rigidity using methods of rational homotopy theory. We prove that a triple $`(C,T,k)`$ is splitting rigid if, for any derivation of the minimal model of $`C`$ that commutes with differential and has degree within $`[dim(T),0)`$, the induced derivation on $`H^{}(C)`$ vanishes on $`Char(C,k)`$. The converse to this statement is proved in 9.4 under various assumptions on $`(C,T,k)`$, such as $`2kdim(C\times T)+3`$, or $`p_i(TC)Char(C,k)`$.
Thus, under either of the assumptions, the splitting rigidity is a phenomenon of rational homotopy theory, in other words, whether or not $`(C,T,k)`$ is splitting rigid depends only on $`k`$, $`dim(T)`$, and the minimal model (or equivalently, the rational homotopy type) of $`C`$. This is no longer true for smaller $`k`$, in fact in Section 11, we give an example of two triples $`(C,T,6)`$, $`(M,T,6)`$ with homotopy equivalent$`C`$ and $`M`$, such that $`(C,T,6)`$ is splitting rigid, while $`(M,T,6)`$ is not.
Let $`(M_C,d_C)`$, $`(M_T,d_T)`$ be (Sullivan) minimal models for $`C`$, $`T`$. Since $`H^{}(T)`$ is a free exterior algebra, we can assume that $`M_T=H^{}(T)`$ and $`d_T=0`$. Then $`(M_CM_T,d)`$ is a minimal model for $`C\times T`$, where for $`xM_C`$, $`tM_T`$
$$d(xt)=d_C(x)t+(1)^{|x|}xd_T(t)=d_C(x)t.$$
First, we modify the arguments of Section 4 to produce the Taylor expansion in $`M_CM_T`$ of any self-homotopy equivalence of $`C\times T`$. Let $`h`$ be a self-homotopy equivalence of $`C\times T`$. Then the induced map of minimal models $`h^\mathrm{\#}`$ is an isomorphism \[FHT01, 12.10(i)\].
For $`xM_C`$, consider $`\frac{h^\mathrm{\#}}{t_i}(x)M_C`$ such that $`h^\mathrm{\#}(x1)=_i(1t_i)(\frac{h^\mathrm{\#}}{t_i}(x)1)`$. Since $`h^\mathrm{\#}`$ commutes with $`d`$ and $`d(1t_i)=0`$, we get that each $`\frac{h^\mathrm{\#}}{t_i}`$ commutes with $`d_C`$, up to sign, and therefore induces a linear self-map of $`H^{}(C)`$.
As in Section 4, since $`h^\mathrm{\#}`$ is an algebra isomorphism, the maps $`\frac{h^\mathrm{\#}}{t_i}`$ satisfy the same recursive identities, obtained from $`h^\mathrm{\#}(xy)=h^\mathrm{\#}(x)h^\mathrm{\#}(y)`$ by collecting the terms next to $`1t_i`$’s. Thus, $`\frac{h}{t_0}`$ is an isomorphism of $`(M_C,d_C)`$, and $`\frac{h^\mathrm{\#}}{t_0}(\frac{h^\mathrm{\#}}{t_0})^1`$ is a derivation of $`M_C`$ of degree $`1`$. More generally, if $`\frac{h^\mathrm{\#}}{t_i}(x)=0`$ for all $`xM_C`$ and all $`0<i<k`$, then $`\frac{h^\mathrm{\#}}{t_k}(\frac{h^\mathrm{\#}}{t_0})^1`$ is a derivation of $`M_C`$ of degree $`|t_k|`$. As we noted before, $`\frac{h^\mathrm{\#}}{t_k}(\frac{h^\mathrm{\#}}{t_0})^1`$ commutes with $`d_C`$, up to sign. Hence, it induces a degree $`|t_k|`$ derivation of $`H^{}(C)`$. A slight variation of the proof of Proposition 4.1, implies the following.
###### Proposition 9.1.
If each negative derivation of $`H^{}(C)`$ induced by a derivation of $`M_C`$ of degree $`dim(T)`$ vanishes on $`Char(C,k)`$, then $`(C,T,k)`$ is splitting rigid.
###### Remark 9.2.
The assumptions of Proposition 9.1 only involve the minimal model of $`C`$. Thus, if $`C^{}`$ is rationally homotopy equivalent to $`C`$ and $`(C,T,k)`$ satisfies the assumptions of Proposition 9.1, then so does $`(C^{},T,k)`$.
The following Lemma shows that an integer multiple of any negative derivation of $`M_C`$ can be “integrated” to a self-homotopy equivalence of $`C\times T`$. This gives many examples of triples which are not splitting rigid, and is a crucial ingredient in the proof of Theorem 1.4.
###### Lemma 9.3.
Let $`D`$ be a negative derivation of $`H^{}(C)`$ induced by a derivation of $`M_C`$ that commutes with $`d_C`$. Let $`T`$ be a torus with $`|D|dim(T)`$. Then there are integers $`i,m>0`$, and a self-homotopy equivalence $`h`$ of $`C\times T`$ such that $`h^{}(a1)=a1+(1t_i)(mD(a)1)`$ for $`aH^{}(C)`$, and $`h_{CC}\mathrm{id}_C`$, $`\pi _Th=\pi _T`$.
###### Proof.
Let $`\stackrel{~}{D}`$ be a derivation of $`M_C`$ that induces $`D`$. The fact $`\stackrel{~}{D}`$ is a derivation implies that the map $`\varphi :M_CM_CM_T`$ defined by $`\varphi (x)=(x1)+(1t_i)(\stackrel{~}{D}(x)1)`$ is a DGA-homomorphism.
Being a DGA-homomorphism, $`\varphi `$ defines a (unique up to homotopy) map of rationalizations $`f:C_0\times T_0C_0`$ where we can choose $`f`$ so that $`fi_{C_0}=\mathrm{id}_{C_0}`$. Look at the diagram
where $`r=r_T\times r_C:C\times TC_0\times T_0`$ is the rationalization, and try to find a finite covering $`p`$, and a map $`\stackrel{~}{f}`$ that makes the diagram commute, and satisfies $`\stackrel{~}{f}i_C=\mathrm{id}_C`$.
It follows from an obstruction theory argument as in \[BK01b, Section 4\] that such $`p`$, $`\stackrel{~}{f}`$ can be constructed, where the key point is that all obstructions are torsion since the homotopy fiber of $`r`$ has torsion homotopy groups, and all torsion obstruction vanish after precomposing with a suitable finite cover $`p`$. Furthermore (cf. the proof of Lemma B.1), one can choose $`p,\stackrel{~}{f}`$ satisfying $`\stackrel{~}{f}i_C=\mathrm{id}_C`$, $`p=\mathrm{id}_C\times (\times ^n)`$, where $`\times ^n:TT`$ is the $`n`$-th power map, so that, for some positive integer $`m`$, the induced map on $`H^{}(C)`$ satisfies
$$\stackrel{~}{f}^{}[a1]=[a1]+m(1t_i)([\stackrel{~}{D}(a)]1).$$
Finally, by Whitehead’s theorem, the map $`h(c,t)=(\stackrel{~}{f}(c,t),t)`$ is a homotopy equivalence with the desired properties. ∎
###### Proposition 9.4.
Assume that $`(C,T,k)`$ is splitting rigid and either $`2kdim(C\times T)+3`$, or $`p_i(TC)Char(C,k)`$ for all $`i>0`$. If $`D`$ is a negative derivation of $`H^{}(C)`$ of degree $`dim(T)`$ induced by a derivation of $`M_C`$, then $`D`$ vanishes on $`Char(C,k)`$.
###### Proof.
Arguing by contradiction, let $`D`$ be a negative derivation of $`H^{}(C)`$ with $`|D|dim(T)`$ such that $`D`$ does not vanish on $`Char(C,k)`$.
First, we show that $`D`$ is nonzero on either Euler or Pontrjagin class of a rank $`k`$ bundle $`\xi _C`$ over $`C`$. Assume, say that $`k=2m`$ so that $`Char(C,k)=(_{i=1}^{m1}H^{4i}(C))H^{2m}(C)`$. If $`aH^{2m}(C)`$ satisfies $`D(a)0`$, then by Lemma B.1, $`a`$ is proportional to the Euler class of some rank $`k`$ vector bundle $`\xi _C`$ over $`C`$. If $`aH^{4i}(C)`$, then by Lemma B.1, $`a`$ is proportional to the $`i`$th Pontrjagin class of some bundle of rank $`k`$ vector bundle $`\xi _C`$ over $`C`$.
Say, suppose that $`D`$ is nonzero on the Euler class $`e(\xi _C)`$. By Lemma 9.3, there is a self-homotopy equivalence $`f`$ of $`C\times T`$ such that $`f^{}(e(\xi _C)1)`$ does not lie in $`H^{}(C)1`$. Now there are two cases to consider.
If $`2kdim(C\times T)+3`$, then by Haefliger’s embedding theorem \[Hae61\] the homotopy equivalence $`f:C\times TE(\xi _C)\times T`$ is homotopic to a smooth embedding $`q`$. The normal bundle $`\nu _q`$ has rank $`k`$. Since $`q`$ and $`f`$ are homotopic, $`\nu _q`$ and $`q^\mathrm{\#}(\xi _C\times T)`$ have equal Euler and Pontrjagin classes. This follows from the intersection pairing interpretation of the Euler class, and the Whitney sum formula for the total Pontrjagin class (see \[BK01b, Section 3\] for details). Thus, $`e(\nu _q)=f^{}(e(\xi _C)1)`$ does not lie in $`H^{}(C)1`$. So $`\nu _q`$ does not virtually come from $`C`$, while $`E(\nu _q)`$ is diffeomorphic to $`T\times E(\xi _C)`$, and this means that $`(C,T,k)`$ is not splitting rigid.
If $`p_i(TC)Char(C,k)`$ for all $`i>0`$, then $`D(p(TC))=0`$. In other words, $`f^{}p(TC)=p(TC)`$, so by Lemma A.1, there is a diffeomorphism $`q`$ of $`C\times T`$ such that $`q^{}(e(\xi _C)1)`$ does not lie in $`H^{}(C)1`$. Then the pullback bundle $`f^\mathrm{\#}(\xi _C\times T)`$ does not virtually come from $`C`$, while its total space is diffeomorphic to $`T\times E(\xi _C)`$; thus $`(C,T,k)`$ is not splitting rigid. ∎
###### Remark 9.5.
Note that if $`kdim(C)`$, then $`p_i(TC)Char(C,k)`$ for all $`i>0`$.
###### Remark 9.6.
The proof of Proposition 9.4 implies the weak converse of the Proposition 3.4: if $`(C,T,k)`$ is splitting rigid, and either $`2kdim(C\times T)+3`$, or $`p_i(TC)Char(C,k)`$, then any homotopy equivalence of $`C\times T`$ maps $`Char(C,k)1`$ to itself.
The following example is due to T. Yamaguchi \[Yam, Example 3\] and it was constructed in response to a question in the previous version of this paper.
###### Example 9.7.
Consider a minimal Sullivan algebra $`M`$ with the generators $`\{x,y,z,a,b,c\}`$, whose degrees are respectively $`2,3,3,4,5,7`$, and the differential given by $`d(x)=d(y)=0`$, $`d(z)=x^2`$, $`d(a)=xy`$, $`d(b)=xa+yz`$, $`d(c)=a^2+2yb`$. A direct computation shows that the rank of $`H^i(M)`$ is equal to $`1`$ for $`i=0,2,3,11,12,14`$, is equal to $`2`$ for $`i=7`$, and is equal to 0 for other $`i`$’s. Also the generators of $`H^{}(C)`$ as an algebra are
$$x,y,e=ya,f=xbza,g=x^2cxab+yzb,h=3xyc+a^3.$$
The products are all trivial, except $`xh`$, $`yg`$, and $`ef`$ in $`H^{14}(M)`$. Thus $`H^{}(M)`$ satisfies the Poincare duality, and therefore by rational surgery \[Sul77, Theorem 13.2\], there exists a closed simply connected elliptic $`14`$-dimensional manifold $`C`$ with minimal model $`M`$.
Then $`\mathrm{Der}(H^{}(C))`$ is not zero since there is, for example, a non-zero derivation $`(g,y)`$ of degree $`8`$. Here $`(p,q)`$ stands for the derivation which send p to q and other generators to zero. One can check that $`(g,y)`$ is indeed a derivation. On the other hand, another direct computation (see \[Yam, Example 3\]) shows that all derivations induced from $`M`$ vanish on $`H^{}(C)`$. Therefore, the manifold $`X=C\times C\times C\times C`$ has the same property, but there exists a derivation $`D`$ of $`H^{}(X)`$ of degree $`8`$ which is non-zero on $`H^{44}(X)`$. Thus, $`(X,T,k)`$ is splitting rigid for any $`T,k`$, but this fact can not be seen by looking only at $`H^{}(X)`$.
The following lemma, combined with Proposition 9.1, shows that the property of being “splitting rigid for all $`k`$” depends only on $`dim(T)`$ and the rational homotopy type of $`C`$.
###### Lemma 9.8.
Let $`(C,T,k)`$ be splitting rigid for all $`k`$. Then if $`D`$ is a negative derivation of $`H^{}(C)`$ of degree $`dim(T)`$ induced by a derivation of $`M_C`$, then $`D`$ vanishes on $`H^{even}(C)`$.
###### Sketch of the proof.
As in the proof of Proposition 9.4, use $`D`$ to construct a self-homotopy equivalence $`f`$ of $`C\times T`$.
If $`D`$ is nonzero on $`H^{4i}(C)`$ for some $`i`$, then $`D`$ is nonzero on the $`i`$th Pontrjagin class of some bundle $`\xi _C`$ over $`C`$. Then for some large $`k`$, the map $`f:C\times TE(\xi _C)\times T`$ is homotopic to a smooth embedding. By the same argument as in the proof of Proposition 9.4, the normal bundle to this embedding does not virtually come from $`C`$ and, therefore, $`(C,T,k)`$ is not splitting rigid.
If $`D`$ vanishes on $`_iH^{4i}(C)`$, then $`D(p(TC))=0`$ so $`f`$ preserves $`p(TC)`$. Hence, by Lemma A.1, replacing $`f`$ with some power of $`f`$, we can assume that $`f`$ is homotopic to a diffeomorphism. If $`D`$ is nonzero on $`H^{2i}(C)`$ for some $`i`$, then $`D`$ is nonzero on the Euler class of some bundle $`\xi _C`$ over $`C`$. Looking at the bundle $`f^\mathrm{\#}(\xi _C\times T)`$ shows that $`(C,T,2i)`$ is not splitting rigid. ∎
## 10 Proof of Theorem 1.4
Let $`G=SU(6)`$ and $`H=SU(3)\times SU(3)`$. According to \[TO97, Chapter5, Example 4.14\] (cf. \[GHV76, Section 11.14\] and \[FHT01, Proposition 5.16\]), the minimal model of $`G/H`$ is given by $`(M,d)=(\mathrm{\Lambda }(y_4,y_6,x_7,x_9,x_{11}),d)`$ with the degrees given by the subscripts, and $`d(y_4)=0=d(y_6)`$, $`d(x_7)=y_4^2`$, $`d(x_9)=2y_4y_6`$, $`d(x_{11})=y_6^2`$.
Now it is straightforward to compute the cohomology algebra of $`G/H`$. In particular, $`G/H`$ has nonzero Betti numbers only in dimensions $`0,4,6,13,15,19`$ and the cohomology groups in dimensions $`4,6`$ are generated by the classes $`[y_4],[y_6]`$.
Let $`\xi `$ be the rank $`3`$ complex vector bundle over $`G/H`$ classified by $`p:G/HBSU(3)`$ which is the composition of the classifying map $`G/HBH=BSU(3)\times BSU(3)`$ for the bundle $`GG/H`$ with the projection $`BSU(3)\times BSU(3)BSU(3)`$ on the first factor. Since $`G/(SU(3)\times 1)`$ is $`6`$-connected, from the Serre spectral sequence of the bundle $`G/SU(3)G/HBSU(3)`$ we see that the map $`p^{}:H^i(BSU(3))H^i(G/H)`$ is an isomorphism for $`i6`$. Hence, $`c_3(\xi )`$ is the generator of $`H^6(G/H)`$, and by rescaling $`y_6`$, we can assume that $`c_3(\xi )=[y_6]`$. Note that $`c_3(\xi )`$ is equal to the Euler class $`e(\xi _{})`$ of $`\xi _{}`$, the realification of $`\xi `$.
Alternatively, $`\xi `$ can be described as the associated bundle to the principal bundle $`GG/H`$ via the representation $`\rho `$ of $`SU(3)\times SU(3)`$ given by the projection onto the first factor followed by the standard action of $`SU(3)`$ on $`^6`$. Thus, $`E(\xi )`$ admits a complete metric with $`\mathrm{sec}0`$ such that the zero section is a soul.
To finish the proof it remains to find a torus $`T`$, and a self-diffeomorphism $`f`$ of $`C\times T`$ such that $`f^{}(c_3(\xi )1)H^{}(G/H)1`$, because then the bundle $`f^\mathrm{\#}\xi _{}`$ does not virtually come from $`C`$, and $`E(f^\mathrm{\#}\xi _{})`$ carries a complete metric with $`\mathrm{sec}0`$ and zero section being a soul.
Because $`M`$ is free, the linear map $`\stackrel{~}{D}:MM`$ defined by $`\stackrel{~}{D}(y_4)=\stackrel{~}{D}(x_7)=0`$, $`\stackrel{~}{D}(y_6)=y_4`$, $`\stackrel{~}{D}(x_9)=x_7`$, $`\stackrel{~}{D}(x_{11})=2x_9`$ is a derivation of $`M`$ of degree $`2`$ (see \[FHT01, page 141\]). By computing on the generators, it is straightforward to see that $`\stackrel{~}{D}`$ commutes with $`d`$, and hence induces a derivation $`D`$ of $`H^{}(G/H)`$ such that $`D([y_6])=[y_4]`$.
By Lemma 9.3, there is a positive integer $`m`$, and a self-homotopy equivalence $`h`$ of $`C\times T`$ where $`dim(T)=2`$ such that $`\pi _Th=\pi _T`$ and $`h^{}(a1)=a1+(1t_3)(mD(a)1)`$ for any $`aH^{}(G/H)`$.
Note that $`h^{}`$ preserves the total Pontrjagin class of the tangent bundle to $`G/H\times T`$. Indeed, since $`T`$ is parallelizable and $`H^{4i}(G/H)`$ are only nonzero if $`i=0,1`$, it suffices to show that $`h^{}`$ preserves $`p_1(G/H)1`$, or equivalently, that $`mD`$ vanishes on $`p_1(G/H)`$. In fact, more is true, namely, $`D`$ vanishes on $`H^4(G/H)`$ since $`D([y_4])=[\stackrel{~}{D}(y_4)]=0`$. By Lemma A.1 below, some power
$$f=h^k=\underset{k}{\underset{}{h\mathrm{}h}}$$
of $`h`$ is homotopic to a diffeomorphism. Since $`\pi _Th=\pi _T`$, we have $`\pi _Tf=\pi _T`$ so that $`f^{}(1t)=1t`$ for any $`tH^{}(T)`$. Combining with $`(t_3)^2=0`$, we get for $`aH^{}(C)`$
$$f^{}(a1)=a1+(1t_3)(kmD(a)1).$$
Since $`D(c_3(\xi ))=D([y_6])=[y_4]0`$, we get $`f^{}(c_3(\xi )1)H^{}(G/H)1`$, and the proof of Theorem 1.4 is complete.
## 11 Proving splitting rigidity without rational homotopy
This section contains an example of two triples $`(C,T,6)`$, $`(M,T,6)`$ such that $`C`$ and $`M`$ are homotopy equivalent and $`(C,T,6)`$ is splitting rigid, while $`(M,T,6)`$ is not.
Let $`M=S^3\times S^3\times S^{10}\times S^{11}`$. Note that $`(M,T,6)`$ is not splitting rigid if $`dim(T)3`$. Indeed, let $`\xi _M`$ be the pullback of $`TS^6`$ via the map $`MS^6`$, which is the composition of the projection $`MS^3\times S^3`$ followed by a degree one map $`S^3\times S^3S^6`$. Then $`\xi _M`$ has nonzero Euler class. If $`dim(T)3`$, then one can easily construct a self-diffeomorphism $`f`$ of $`M\times T`$ such that $`f^{}e(\xi _M\times T)`$ does not lie in $`H^6(M)1`$, in particular, the bundle $`f^\mathrm{\#}(\xi _M\times T)`$ does not virtually come from $`M`$.
###### Proposition 11.1.
There exists a smooth manifold $`C`$ homotopy equivalent to $`M`$ such that $`(C,T,6)`$ is splitting rigid for any $`T`$.
###### Proof.
By a surgery argument as in \[BK01a, A.1\], there is a closed smooth manifold $`C`$ which is homotopy equivalent to $`M`$ and has $`p(TC)=1+p_4(TC)`$ with $`p_4(TM)0`$.
To check that $`(C,T,6)`$ is splitting rigid, we need to start with an arbitrary rank $`6`$ vector bundle $`\xi `$ over $`B=C\times T`$ that satisfies ($``$), and prove that $`\xi `$ virtually comes from $`C`$, or equivalently, that the Euler and Pontrjagin classes of $`\xi `$ lie in $`Char(C,6)1`$.
We shall borrow notations and arguments from the proof of Proposition 3.4. As in 3.4, we can assume that $`E(\xi )`$ is the total space of a vector bundle $`\eta `$ which is the product of $`T`$ and a vector bundle $`\eta _C^{}`$ over a closed smooth simply-connected manifold $`C^{}`$. Let $`S=C^{}\times T`$ and $`g:BS`$ by the homotopy equivalence as in 3.4.
Note that $`g^{}p(TS)=p(TB)`$. Indeed, as in 3.4, $`g^{}p(TN|_S)=p(TN|_B)`$. So $`g^{}p(\eta )g^{}p(TS)=p(\xi )p(TB)`$. The only Pontrjagin classes of a rank $`6`$ vector bundle that have a chance of being nonzero are $`p_1`$, $`p_2`$, $`p_3`$. Since $`C^{}`$ has zero cohomology in dimensions $`4`$, $`8`$, $`12`$, we get $`p(\eta )=1`$. By the same argument, $`p(TS)=1+p_4(TS)+p_6(TS)`$. We have
$$g^{}(1+p_4(TS)+p_6(TS))=(1+p_1(\xi )+p_2(\xi )+p_3(\xi ))(1+p_4(TB)),$$
hence $`p_6(TS)=0=p_i(\xi )`$ for all $`i`$, and $`g^{}p(TS)=p(TB)`$. Then one easily sees that $`g_{CC^{}}`$ maps $`p(TC^{})`$ to $`p(TC)`$. Hence, $`g_{CC^{}}\times g_{TT}`$ maps $`p(TS)`$ to $`p(TB)`$, and therefore, $`h=(g_{CC^{}}^1\times g_{TT}^1)g`$ preserves $`p(TB)`$.
As in 3.4, it suffices to show that $`h`$ preserves $`Char(C,6)1=H^6(C)1`$. We think of $`H^{}(C)`$ as an exterior algebra on generators $`x,y,q,s`$ corresponding to spheres $`S^3`$, $`S^3`$, $`S^{10}`$, $`S^{11}`$ so we need to show that $`h^{}(xy)=xy`$. By rescaling we can assume that $`p(TB)=xyq1`$.
By dimension reasons $`\frac{h^{}}{t_i}(xy)=0`$ unless $`|t_i|`$ is $`3`$ or $`6`$. Similarly, $`\frac{h^{}}{t_i}(q)=0`$ unless $`|t_i|`$ is $`4`$, $`7`$ or $`10`$. Now collecting terms next to $`1t_i`$’s in the identity $`xyq1=h^{}(xyq1)`$, we conclude that $`q\frac{h^{}}{t_i}(xy)=0`$, and hence $`\frac{h^{}}{t_i}(xy)=0`$ for all $`i>0`$. Thus, $`h^{}(xy)=xy`$ as promised. ∎
###### Remark 11.2.
The above example shows that Proposition 9.4 fails without assuming either $`2kdim(C\times T)+3`$, or $`p_i(TC)Char(C,k)`$. Indeed, $`(H^{}(C),0)`$ is a minimal model of $`C`$ and there is a degree $`3`$ derivation of $`H^{}(C)`$ given by $`D(x)=1`$, $`D(y)=D(q)=D(s)=0`$ which does not vanish on $`Char(C,6)=H^6(C)`$. Namely, $`D(xy)=y`$. Yet $`(C,T,6)`$ is splitting rigid.
It is instructive to see where the proof of Proposition 9.4 fails. Using $`D`$, we produce a self-homotopy equivalence $`f`$ of $`C\times T`$, and an $`^6`$-bundle $`\xi _C`$ with $`xy=e(\xi _C)1`$. However, the homotopy equivalence $`f:C\times TE(\xi _C)\times T`$ is not homotopic to a smooth embedding.
###### Remark 11.3.
As always with splitting rigid triples, the total spaces of “most” $`^6`$-bundles over $`C\times T`$ do not admit complete metrics with $`\mathrm{sec}0`$. We do not know whether $`C`$ in Proposition 11.1 admits a metric with $`\mathrm{sec}0`$. Yet, no currently known method rules out the existence of $`\mathrm{sec}0`$ on $`C`$, because $`C`$ is homotopy equivalent to a closed nonnegatively curved manifold, and $`C`$ admits a metric of positive scalar curvature, for $`C`$ is spin and $`dim(C)=273(\mathrm{mod}8)`$, so \[Sto92\] applies.
## 12 Nonnegatively curved vector bundles with souls equal to the zero sections
The purpose of this section is to obtain restrictions on normal bundles to souls in nonnegatively curved manifolds. In other words, we look for conditions on a vector bundle $`\xi `$ ensuring that $`E(\xi )`$ admits no complete nonnegatively curved metric such that the zero section is a soul. The assumption that a given submanifold is a soul imposes a nontrivial restriction on the metric, so it is no surprise that we get stronger results on obstructions.
Our exposition is parallel to the one in Section 3. We say that a vector bundle $`\xi `$ over $`C\times T`$ satisfies condition ($``$) if
$`\begin{array}{c}\text{there is a finite cover }\pi :C\times TC\times T,\text{a closed manifold}C^{},\text{and}\hfill \\ \text{a diffeomorphism }f:C^{}\times TC\times T\text{ such that the bundle }f^\mathrm{\#}\pi ^\mathrm{\#}(\xi )\hfill \\ \text{virtually comes from}C^{}.\hfill \end{array}`$ $`()`$
###### Example 12.1.
According to Theorem 3.1, if $`S`$ is a soul in a complete nonnegatively curved manifold, then the normal bundle to $`S`$ satisfies ($``$).
Caution. Clearly, if $`\xi `$ satisfies ($``$), it also satisfies condition ($``$) from Section 3. The converse is generally false, as the following example shows.
###### Example 12.2.
Let $`M=S^3\times S^5\times S^7`$. By a surgery argument as in \[BK01a, A.1\], there is a closed $`15`$-dimensional manifold $`C`$ which is homotopy equivalent to $`M`$ and such that $`p_2(TC)`$, $`p_3(TC)`$ are nonzero. Clearly, there is a derivation of $`H^{}(C)`$ of degree $`3`$ which is nonzero on $`H^8(C)`$. Since $`(H^{}(C),0)`$ is the minimal model for $`C`$, Proposition 9.4 implies that $`(C,T,15)`$ is not splitting rigid if $`dim(T)3`$. Thus, there is a rank $`15`$ bundle $`\xi `$ over $`C\times T`$ which satisfies ($``$) but which does not virtually come from $`C`$. It remains to show that $`\xi `$ does not satisfy ($``$). If it does, then after passing to a finite cover, we can assume that, for some diffeomorphism $`f:C^{}\times TC\times T`$, $`f^\mathrm{\#}\xi `$ virtually comes from $`C^{}`$. Since $`f`$ is a diffeomorphism and $`T`$ is parallelizable, $`f^{}(p_i(TC)1)=p_i(TC^{})1`$ for all $`i`$. Since $`p_i(TC)`$ generates $`H^{4i}(C)`$ and $`f^{}`$ is an isomorphism, this means that $`f^{}(H^{4i}(C)1)=H^{4i}(C^{})1`$ for all $`i`$. Since $`f^{}p_i(\xi )=p_i(f^\mathrm{\#}(\xi ))H^{4i}(C^{})1`$ for all $`i`$, this implies that $`p_i(\xi )H^{4i}(C)1`$ for all $`i`$ and hence $`\xi `$ virtually comes from $`C^{}`$. This is a contradiction, so $`\xi `$ cannot satisfy ($``$).
Recall that, given a compact Lie group $`G`$, a rank $`n`$ vector bundle $`\xi `$ has structure group $`G`$ if $`\xi `$ is associated with a principal $`G`$-bundle via some representation $`GO(n)`$. We now have the following splitting criterion similar to Proposition 3.4:
###### Proposition 12.3.
Let $`\xi `$ be a vector bundle over $`C\times T`$ which has structure group $`O(k)`$ and satisfies ($``$). If any self homotopy equivalence of $`C\times T`$ maps $`Char(C,k)1`$ to itself, then $`\xi `$ virtually comes from $`C`$.
###### Proof.
Passing to a finite cover, we can assume that $`\xi `$ is orientable, and that $`\eta =f^\mathrm{\#}\xi `$ virtually comes from $`C^{}`$, where $`f:C^{}\times TC\times T`$ is a diffeomorphism. As before, to show that $`\xi `$ virtually comes from $`C`$ it is enough to check that its rational characteristic classes lie in $`H^{}(C)1`$. Proceeding exactly as in the proof of Proposition 3.4, we conclude that $`(f^{})^1`$ maps $`Char(C^{},k)1`$ to $`Char(C,k)1`$.
Next note that all $`e(\eta )`$, $`p_i(\eta )`$ lie in the subalgebra generated by $`Char(C^{},k)1`$. Indeed, $`\eta `$ has the structure group $`SO(k)`$, so $`\eta `$ is a pullback of a bundle over $`BSO(k)`$. Since the cohomology of $`BSO(k)`$ is generated by $`Char(BSO(k),k)`$, by naturality of the characteristic classes, we see that $`e(\eta ),p_i(\eta )Char(C^{},k)`$ for any $`i`$.
Since $`(f^{})^1(Char(C^{},k)1)Char(C,k)1`$, we conclude that all the characteristic classes of $`\xi `$ lie in $`Char(C,k)1H^{}(C)1`$, hence $`\xi `$ virtually comes from $`C`$. ∎
Now all the splitting rigidity results that relied on Proposition 3.4, can be adapted to this new setting. In particular, we obtain
###### Theorem 12.4.
Let $`\xi `$ be a vector bundle over $`C\times T`$ with structure group $`O(2)`$. If $`E(\xi )`$ admits a complete metric with $`\mathrm{sec}0`$ such that the zero section is a soul, then $`\xi `$ virtually comes from $`C`$.
###### Theorem 12.5.
Let $`C=G//H`$ be a simply connected biquotient of compact Lie groups such that $`H`$ is semi-simple. Let $`\xi `$ be a vector bundle over $`C\times T`$ whose structure group can be reduced to a subgroup of $`O(4)`$. If $`E(\xi )`$ admits a complete metric with $`\mathrm{sec}0`$ such that the zero section is a soul, then $`\xi `$ virtually comes from $`C`$.
## 13 Open problems
### 13.1 Induced derivations
As we proved in Section 9, for sufficiently large $`k`$, splitting rigidity of $`(C,T,k)`$ is equivalent to vanishing on $`Char(C,k)`$ of all negative derivations of degree $`\mathrm{dim}(T)`$ induced from the minimal model $`M_C`$. Yet all our geometric applications are proved by checking the stronger condition that all negative derivations vanish on $`Char(C,k)`$. This is mostly due to the fact that we do not know how to effectively check which derivations of $`H^{}(C)`$ are induced from the minimal model.
Let us restate the problem in purely rational homotopy theoretic terms. It is well known that the space $`\mathrm{Der}(M_C)`$ is a differential graded Lie algebra (DGLA) with the differential given by $`D_C=[,d_C]`$. It is trivial to check that closed derivations preserve $`\mathrm{ker}(d_C)`$ and exact ones send $`\mathrm{ker}(d_C)`$ to $`\mathrm{Im}(d_C)`$ (see \[Gri94\] for details). Therefore, we have a natural graded Lie algebra homomorphism $`m:H^{}(\mathrm{Der}(M_C))\mathrm{Der}(H^{}(C))`$. We seek to understand the image of the map $`m`$. Example 9.7 produces an elliptic smooth manifold $`C`$ such that $`\mathrm{Im}(m)=0`$ but $`\mathrm{Der}(H^{}(C))0`$.
To relate to our geometric applications we would like to find such examples when $`C`$ is nonnegatively curved. It is easy to see that $`m`$ is onto if $`C`$ is formal, but that is all we can generally say at the moment. (Recall that a space $`X`$ is called formal if $`X`$ and $`(H^{}(X),d=0)`$ have isomorphic minimal models).
Let us also mention that according to Sullivan \[Sul77\], the DGLA $`(\mathrm{Der}(M_C),D_C)`$ is a (Quillen) Lie algebra model for $`Baut_1(C)`$ (the classifying space for the identity component of the monoid of self-homotopy equivalences of $`C`$), and therefore, understanding the map $`m`$ can be helpful for computing the rational homotopy groups of $`Baut_1(C)`$.
### 13.2 Halperin’s conjecture for biquotients
As we mentioned in the introduction, the conjecture of Halperin that any elliptic space of positive Euler characteristic belongs to $``$ has been verified for all homogeneous spaces of compact Lie groups \[ST87\]. However, it remains open for the natural bigger class of elliptic spaces formed by biquotients. According to \[Sin93\], a biquotient $`G//H`$ has a positive Euler characteristic iff $`\mathrm{rank}(G)=\mathrm{rank}(H)`$. Therefore, we pose the following
###### Problem 13.1.
Prove that any biquotient $`G//H`$ belongs to $``$ if $`\mathrm{rank}(G)=\mathrm{rank}(H)`$.
This is unknown even for the simplest examples such as $`Sp(1)\backslash Sp(n)/SU(n)`$.
### 13.3 Nonnegatively curved vector bundles over rational H-spaces
Most explicit examples of nonnegatively curved bundles are given by homogeneous vector bundles, i.e. by vector bundles associated to principal $`H`$-bundles $`HGG/H`$ via some representations $`HO(k)`$. For any given $`k`$, there are only finitely many rank $`k`$ homogeneous vector bundles (because the number of nonequivalent irreducible representations $`HO(k)`$ is finite). However, homogeneous vector bundles can fill a substantial part of $`[G/H,BO]`$, the set of stable equivalence classes of bundles over $`G/H`$.
If $`H=1`$, or more generally if $`G/H`$ is rationally homotopy equivalent to the product of odd-dimensional spheres, then $`H^{}(BH)H^{}(G/H)`$ has trivial image \[GHV76, page 466\], and hence, any homogeneous vector bundle over $`G/H`$ has zero Euler and Pontrjagin classes. Motivated by the above discussion, we pose the following:
###### Problem 13.2.
Does there exist a nonnegatively curved vector bundle $`\xi `$ over a closed manifold $`C`$ such that $`C`$ is rationally homotopy equivalent to the product of odd-dimensional spheres, $`\mathrm{sec}(C)0`$, and $`e(\xi )0`$ or $`p_i(\xi )0`$ for some $`i>0`$?
### 13.4 Nonnegatively curved nonsplitting rigid examples
The example of a nonnegatively curved vector bundle $`\xi `$ that satisfies condition ($``$) but does not virtually come from $`C`$ provided by Theorem 1.4 is essentially the only example of this kind known to us. This is not very satisfactory, say, because this example is unstable; that is, $`\xi ϵ^1`$ does virtually come from $`C`$. More importantly, we want to understand how “generic” such examples are.
To construct a stable example, it suffices to find a vector bundle $`\xi `$ over $`C`$ with $`\mathrm{sec}(E(\xi ))0`$ and soul equal to the zero section, and a negative derivation $`D`$ of $`H^{}(C)`$ induced by a derivation of the minimal model such that $`D(p(\xi ))0`$ and $`D(p(TC))=0`$. While we think that many such examples exist, finding an explicit one, say among homogeneous vector bundles, seems to be an unpleasant task because
1. Pontrjagin classes of homogeneous vector bundles are often difficult to compute, and
2. there is no easy algorithm for computing the space of negative derivations of $`H^{}(G/H)`$ or of the minimal model of $`G/H`$.
One of the few cases when $`\mathrm{Der}_{}(H^{}(G/H))`$ is easily computable is when $`G/H`$ is formal. According to \[Oni94, Theorem 12.2\], any formal compact homogeneous space $`G/H`$ is rationally homotopy equivalent to the product of odd-dimensional spheres and an elliptic space $`X`$ of positive Euler characteristic. Again by \[Oni94, Theorem 12.2\], the image of the homomorphism $`H^{}(BH)H^{}(G/H)`$ is equal to the $`H^{}(X)`$-factor. If $`X`$ (i.e. if Halperin’s conjecture holds for $`X`$), then by Lemma 6.1, one concludes that if $`\xi `$ is a homogeneous vector bundle and $`D`$ is a negative derivation of $`H^{}(G/H)`$, then $`D`$ vanishes on $`e(\xi )`$, $`p_i(\xi )`$ for $`i>0`$. Thus, if Halperin’s conjecture is true, then homogeneous vector bundles over formal homogeneous spaces cannot be used to prove an analog of Theorem 1.4.
Finally, note that a positive solution to Problem 13.2 (for $`C`$ with $`p(TC)=1`$ which includes the case when $`C`$ is a compact Lie group or the product of odd-dimensional spheres) yields an analog of Theorem 1.4, because for any nontrivial element $`a`$ of $`H^{}(C)`$, there exists a negative derivation $`D`$ of $`H^{}(C)`$ with $`D(a)0`$, and $`D(p(TC))`$ vanishes by assumption.
## Appendix A Surgery-theoretic lemma
We are grateful to Ian Hambleton for sketching the proof of the following lemma.
###### Lemma A.1.
Let $`C`$ be a closed smooth simply-connected manifold, and $`T`$ be a torus such that $`dim(C)+dim(T)5`$. If $`h`$ is a self-homotopy equivalence of $`C\times T`$ that preserves the rational total Pontrjagin class of $`C\times T`$, then $`h^m`$ is homotopic to a diffeomorphism for some $`m>0`$.
###### Proof.
Let $`k=dim(T)`$, $`n=dim(C)`$, $`I=[0,1]`$, $`B=C\times T`$, $`\pi =\pi _1(B)`$ so that $`\pi ^k`$. Look at the following commutative diagram whose rows are the smooth and the topological surgery exact sequences:
First, note that $`[B\times I\text{rel},G/Top]L_{n+k+1}(\pi )`$ is onto. Indeed, by the Poincare duality with $`L`$-theory coefficients $`[B\times I\text{rel},G/Top]=H^0(B\times I;𝐋)H_{n+k+1}(B\times I;𝐋)`$, so it suffices to show that the homology assembly map (see e.g. \[Dav94, p216\])
$$A_{n+k+1}:H_{n+k+1}(B\times I;𝐋)L_{n+k+1}(\pi )$$
is onto. By naturality of the assembly and since $`T`$ is the classifying space for $`\pi _1(B\times I)`$, $`A_{n+k+1}`$ factors as the composition of the map $`H_{n+k+1}(B\times I;𝐋)H_{n+k+1}(T;𝐋)`$ induced by the projection $`B\times IT`$, and the universal assembly $`H_{n+k+1}(T;𝐋)L_{n+k+1}(\pi )`$. The former map is onto, since it has a section induced by a section of $`B\times IT`$, while the latter map is an isomorphism since $`\pi ^k`$ \[Wal99, Chapter 15B\]. Thus, $`A_{n+k+1}`$ is onto.
It is known that the map $`[B\times I\text{rel},G/O][B\times I\text{rel},G/Top]`$ has a finite cokernel (see e.g. \[Dav94, page 213\]), and hence so does the map $`[B\times I\text{rel},G/O]L_{n+k+1}(\pi )`$.
By exactness of the smooth surgery exact sequence, the $`L_{n+k+1}(\pi )`$-action on $`S_O(B)`$ has finite orbits.
Since $`h`$ preserves the total Pontrjagin class, $`h^m`$ is tangential for some $`m>0`$, so replacing $`h`$ by $`h^m`$, we can assume that $`h`$ is tangential. Hence for any integer $`l>0`$, $`h^l`$ is tangential so that the normal invariant of $`[B,h^l]S_O(B)`$ lies in the image of the map $`[B,SG][B,G/Top]`$ induced by the fibration $`STopSGG/Top`$. Since $`SG`$ is rationally contractible, $`[B,SG]`$ is a finite set, so there exists an infinite sequence of positive integers $`l_k`$ such that the elements $`[B,h^{l_k}]S_O(B)`$ have the same normal invariant. By exactness, $`[B,h^{l_k}]`$ lie in the same $`L_{n+k+1}(\pi )`$-orbit which is a finite set by above. In particular, for some $`p>q>0`$ we have $`[B,h^p]=[B,h^q]`$. Thus, for some self-diffeomorphism $`f`$ of $`B`$, we get that $`fh^q`$ and $`h^p`$ are homotopic, or $`f`$ is homotopic to $`h^{pq}`$, as wanted. ∎
## Appendix B Vector bundles with prescribed characteristic classes
The following lemma is probably well-known, yet there seems to be no reference available, so we include a complete proof.
###### Lemma B.1.
Let $`X`$ be a finite CW-complex, $`n`$ be a positive integer.
(i) If $`k=2n+1`$, then for any $`n`$-tuple $`(p_1,\mathrm{},p_n)`$ of cohomology classes with $`p_iH^{4i}(X)`$ for $`i=1,\mathrm{},n`$, there is an integer $`m>0`$ and an orientable rank $`k`$ vector bundle $`\xi `$ over $`X`$ such that $`p_i(\xi )=mp_i`$ for $`i=1,\mathrm{},n`$.
(ii) If $`k=2n`$, then for any $`n`$-tuple $`(p_1,\mathrm{},p_{n1},e)`$ of cohomology classes with $`p_iH^{4i}(X)`$ for $`i=1,\mathrm{},n1`$, and $`eH^{2n}(X)`$, there is an integer $`m>0`$ and an orientable rank $`k`$ vector bundle $`\xi `$ over $`X`$ such that $`e(\xi )=me`$, $`p_i(\xi )=mp_i`$ for $`i=1,\mathrm{},n1`$.
###### Proof.
We only give a proof for $`k=2n+1`$; the even case is similar. Let $`\gamma ^k`$ be the universal $`k`$-bundle over $`BSO(k)`$. We think of $`p_i(\gamma ^k)H^{4i}(BSO(k))[BSO(k),K(,4i)]`$ as a map $`BSO(k)K(,4i)`$. It is well-known that the map
$$c=(p_1(\gamma ^k),\mathrm{},p_n(\gamma ^k)):BSO(k)K=K(,4)\times \mathrm{}\times K(,4n)$$
is a rational homotopy equivalence and thus the homotopy groups of its homotopy fiber $`F`$ are torsion. Moreover since $`BSO(k)`$ and $`K(Z,m)`$ admit a CW structure with finitely many cells in every dimension, all homotopy groups of $`F`$ are finitely generated.
Similarly, consider the map $`f=(p_1,\mathrm{},p_n):XK`$ and try to lift it to $`BSO(k)`$. In other words, try to find $`\stackrel{~}{f}:XBSO(k)`$ which would make the following diagram commute up to homotopy:
We first study the auxiliary problem of trying to lift the identity map $`g=\mathrm{id}:KK`$ to a map $`\stackrel{~}{g}:KBSO(k)`$
The obstructions to lifting $`g`$ lie in finite groups $`H^{+1}(K;\pi _{}(F))`$, and are generally nonzero.
Each factor $`K(,i)`$ of $`K`$ is an $`H`$-space, and hence $`K`$ is too. Let $`\times _i^m:K(,i)K(,i)`$ (respectively $`\times ^m:KK`$) be the the $`m`$th power map.
It is easy to see that $`(\times _i^m)^{}:H^i(K(,i);)H^i(K(,i);)`$ is multiplication by $`m`$.
By \[HQ17, Lemma 4\] the induced map $`(\times _i^m)^{}:H^{>0}(K(,i);/m)H^{>0}(K(,i);/m)`$ is identically zero for any $`m>1,i>0`$.
Let $`o_j(g)H^{j+1}(K;\pi _j(F))`$ be the first nontrivial obstruction to lifting $`g`$. Since $`\pi _j(F)`$ is finite, by above we can find an $`m`$ such that $`(\times ^m)^{}(o_j(g))=0`$.
Naturality of obstructions under pullbacks gives $`o_j(\times ^m)=o_j(g\times ^m)=(\times _i^m)^{}(o_j(g))=0`$.
Let $`l`$ be the largest dimension of a cell in $`X`$. By repeating the above process finitely many times we find some $`m=m(l)`$ such that all the obstructions to lifting $`\times ^m`$ on the $`l`$-skeleton $`K_l`$ of $`K`$ vanish. Since $`f:XK`$ can be homotoped to have the image in $`K_l`$ the map $`\times ^mf`$ can be lifted as well.
Thus, there is a map $`\stackrel{~}{f}:XBSO(k)`$ such that $`c\stackrel{~}{f}`$ is homotopic to $`\times ^mf`$.
Now the bundle $`\xi =\stackrel{~}{f}^\mathrm{\#}(\gamma ^k)`$ has the desired properties. ∎
Igor Belegradek
School of Mathematics
Georgia Institute of Technology
Atlanta, GA, USA 30332-0160
email: ib@math.gatech.edu
Vitali Kapovitch
Department of Mathematics, University of Toronto
Toronto, ON, Canada M5S 2E4
email: vtk@math.toronto.edu
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# On the effect of coronal outflow on spectra formation in galactic black hole systems
## 1 Introduction
Hard X-ray spectra of the galactic black hole systems are well described by a power law primary emission along with the pronounced reflected component, which causes the observed flattening of the spectrum. The primary emission is likely to be produced by Compton upscattering of soft photons on thermal electrons in hot, optically thin medium close to a relatively cold accretion disc, being the source of seed photons for Comptonization (see e.g. review in Poutanen 1998). A fraction of the upscattered photons is directed towards the disk and can be reflected from its surface, giving the rise to reflected continuum and fluorescent iron line emission (Lightman & White 1988; George & Fabian 1991).
Observational data for Cyg X-1 and other black hole systems in their hard/low state show often rather hard spectra (photon spectral index $`\mathrm{\Gamma }1.51.9`$; Poutanen et al. 1997; Gierliński et al. 1997; Dove et al. 1997; see also Poutanen 1998), while the amplitude of reflection $`R`$ covers the broad range of values between 0 and 2. Moreover, $`R`$ and $`\mathrm{\Gamma }`$ are correlated (Zdziarski, Lubiński & Smith 1999; Revnivtsev, Gilfanov & Churazov 1999), in the sense that the harder the spectrum, the smaller the amplitude of reflection. The correlation exists both within the low/hard state and when sources change their spectral state (Życki, Done & Smith 1998).
These observations cannot be explained by the model in which static, continuous corona covers the cold accretion disc, as it predicts the power law slope $`\mathrm{\Gamma }2`$ (i.e. rather soft spectra) and the reflection amplitude $`R1.0`$. Among the possible models, which could reproduce the reflection amplitude in the range $`R=01`$ there are two competitive: (i) cold disc disrupted in the inner part (e.g. Poutanen, Krolik & Ryde 1997; Esin, McClintock & Narayan 1997) and (ii) a highly ionized, non-disrupted disc (Nayakshin, Kazanas & Kallman 2000; Ross, Fabian & Young 1999). Detailed shape of the reflected continuum depends on the geometry, ionization state and abundances of elements in the scattering medium. However, spectral fitting does not always allow to constrain these parameters independently. It is possible to explain the observed spectra of GBH in terms of weakly ionized, or neutral reflection from the disc which inner radius is of the order of 50 $`R_\mathrm{g}`$ (e.g. Done & Życki 1999 for the hard state of Cyg X-1), as well as with highly ionized reflection from the disc extending to the marginally stable orbit (as suggested by Ross et al. 1999 for the hard state of Cyg X-1; see also Done & Nayakshin 2000).
The third possible model, in which both the reflection amplitudes $`R>1`$ and $`R<1`$ are possible, is a mildly relativistic outflow/inflow in the corona (Beloborodov 1999). Relativistic aberration reduces the hard X-ray flux scattered towards the disc, which leads to reduction of the reflected component and the soft flux from reprocessing entering the corona. In order to obtain quantitative agreement with observed spectral indices and reflection amplitudes, the model requires additional reduction of the soft flux intercepted by the hot plasma. This leads to the ’active regions’ geometry (e.g. magnetic flares above the disc; Haardt, Maraschi & Ghisellini 1994) rather than a continuous corona. However, the coronal outflow may reproduce the reflection amplitude $`R<1`$ and in order to obtain $`R>1`$ an inflow of the plasma must be postulated.
In this article we reanalyze the non-static corona model, with plasma moving at relativistic speed in the direction perpendicular to the disc surface. We study the dependence of the amplitude of the reflection component on the bulk motion velocity, taking into account the thermal motion of electrons within the plasma determined by the electron temperature $`T_\mathrm{e}`$ and we discuss the resulting shape of the spectra. We also check if high reflection amplitudes ($`R>1`$) might be explained in the frame of the outflow model but taking into account possible high ionization of the disc surface and the anisotropy of Compton scattering. We emphasize the importance of the first Compton scattering, which plays crucial role at lower energies. For highly ionized disc surface the effect of absorption by heavy elements is reduced and the reflected spectrum makes a substantial contribution to the total spectrum in the $`0.55`$ keV band. Therefore we present firstly the semi-analytical calculations in the approximation of single scattering, and after that we perform numerical simulations of multiple scattering, which is responsible for the power law shape of the hard X-ray tail.
The contents of the paper is as following. In Section 2 we analyze semi-analytically the amplitude of the reflection in a single scattering approximation (after Ghisellini et al. 1991) but taking into account anisotropy of scattering within the rest frame of the electron, describing the thermal motion without the assumption of highly relativistic beaming and incorporating the systematic bulk motion (outflow) of the corona. Since the effect of multiple scattering is essential in a real situation, in Section 3 we repeat the computations for a corona of a given optical depth and an electron temperature using a Monte Carlo method for both outwards and inwards directions of bulk velocity. The dependence of the reflection amplitude on the outflow velocity is presented in Section 3.1. The spectra resulting from Monte Carlo computations are presented in Section 3.2. The discussion of the results is given in Section 4.
## 2 Single scattering approximation for the amplitude of reflection
In this section we generalize the determination of the amplitude of the reflection component derived by Ghisellini et al. (1991) for a slab geometry of a hot corona and anisotropic soft photon input from underlying disc. Those results were obtained assuming isotropic electron scattering in the rest frame of an electron and relativistic chaotic motion of electrons. We introduce subsequently the anisotropy of the Thomson cross-section (Section 2.1), we relax the assumption that the thermal motion of the plasma is relativistic (Section 2.2), and finally we introduce the systematic bulk motion of the corona (Section 2.3).
We show, that including the effect of anisotropic soft photon distribution in Compton scattering process the reflection amplitude $`R2`$ can be obtained. It’s value is somewhat larger when the angular dependence of Thomson cross-section is taken into account. On the other hand for low electron velocities ($`\gamma <2`$) the angular distribution of scattered radiation becomes important and the value of $`R`$ is reduced. Adding the systematic coronal outflow reduces the anisotropy effect only in the range of high bulk velocities. However all these results stand mostly for the first scattering, while for the multiple scattering in the cloud of optical depth $`\tau 1`$ the reflection is weaker (see Section 3).
### 2.1 The effect of angular dependence of the Thomson cross-section
In this section we calculate the power of Compton radiation scattered by electrons with isotropic relativistic velocity field $`v=\beta c`$, $`\beta 1`$. This approach is appropriate for a non-thermal plasma. The radiation field is anisotropic and incoming photons subtend a restricted solid angle. We assume the geometry is the same as in Ghisellini et al. (1991). However, in our calculations we take into account the angular dependence of Thomson cross-section (Rybicki & Lightman, 1979). Therefore the emitted power is given by
$`P(\alpha ,\gamma )={\displaystyle \frac{3}{8}}{\displaystyle \frac{P_{\mathrm{iso}}}{(1+\beta ^2/3)}}{\displaystyle _{\varphi _{\mathrm{min}}}^{2\pi \varphi _{\mathrm{min}}}}𝑑\varphi `$ (1)
$`{\displaystyle _{\theta _{\mathrm{min}}}^{\theta _{\mathrm{max}}}}(1\beta \mathrm{cos}\theta )^2(1+\mathrm{cos}^2\theta )\mathrm{sin}\theta d\theta `$
Here $`P_{\mathrm{iso}}`$ is the power of isotropic emission and $`P(\alpha )`$ is the power of radiation scattered by electron of the angle $`\alpha `$ between its velocity and the axis of symmetry of incoming photons. The limits $`\theta _{\mathrm{min}}`$, $`\theta _{\mathrm{max}}`$ and $`\varphi _{\mathrm{min}}`$are functions of $`\alpha `$ and are given by equations (see Ghisellini et al., 1991):
$`\theta _{\mathrm{min}}=\mathrm{max}(0,\alpha \pi /2)`$
$`\theta _{\mathrm{max}}=\mathrm{min}(\alpha +\pi /2,\pi )`$
$$\varphi _{\mathrm{min}}=\{\begin{array}{cc}0\hfill & \text{for }0<\theta <\pi /2\alpha \hfill \\ & \text{or }3\pi /2\alpha <0<\pi \hfill \\ \mathrm{arccos}[(\mathrm{tan}\theta \mathrm{tan}\alpha )^1]\hfill & \text{otherwise}\hfill \end{array}$$
Here we follow the assumption of Ghisellini et. al (1991), that all photons are emitted in the direction of electron’s velocity. This is justified in the case of $`\gamma 1`$, when the direction of the motion of an electron (i.e the angle $`\alpha `$) can be identified with the viewing angle. In Figure 1 we plot the ratio of the total power emitted by electrons moving downwards to that of electrons moving upwards:
$$R(\gamma )=\frac{_{\pi /2}^\pi P(\alpha )\mathrm{sin}(\alpha )𝑑\alpha }{_0^{\pi /2}P(\alpha )\mathrm{sin}(\alpha )𝑑\alpha }$$
(3)
It is worth seeing that this ratio saturates for large $`\gamma `$ at the value of $`R2.55`$, which is somewhat larger than that obtained under assumption of $`d\sigma /d\mathrm{\Omega }=\sigma _\mathrm{T}/4\pi `$ ($`R2.2`$). This result might be of possible importance for non-thermal Comptonization models. It seems however to be in conflict with the recent results of spectral fitting of the data for the soft state of Cyg X-1 (Gierliński et al. 1999).
In Figure 2 we plot the ratio of the power emitted downwards to that emitted by electron moving in one particular direction, given by
$$R(\alpha ,kT)=\frac{_{\pi /2}^\pi P(\alpha )\mathrm{sin}(\alpha )𝑑\alpha }{P(\alpha )}.$$
(4)
This can be identified with the reflection amplitude seen by the observer with the viewing angle $`\alpha `$. We plot the result versus the electron temperature instead of $`\gamma `$ in order to expose better the lower velocity part of the curve but we remind that those results (dashed line) are not supposed to be accurate for not very highly relativistic velocities. We see that for low inclination angle the expected enhancement of the reflection is very large, up to a factor 3.5 for typical plasma temperatures.
### 2.2 The effect of mildly relativistic thermal motion
If the plasma is rather thermal than non-thermal, the electron velocity is only mildly relativistic and the assumption of all the emission being beamed in the direction of motion of an electron used in Section 2.1 is not valid. Therefore, in this section we consider the angular distribution of the scattered radiation, given by
$$P(\alpha ,\gamma ,\mathrm{\Theta }_{\mathrm{out}})=\frac{P(\alpha ,\gamma )}{2\gamma ^4(1+\beta \mathrm{cos}\mathrm{\Theta }_{\mathrm{out}})^3}$$
(5)
where $`\mathrm{\Theta }_{\mathrm{out}}=\pi `$ means scattering in the direction of electron’s movement. The accretion disc receives a fraction of scattered radiation that depends on the direction and velocity of electron
$`P_{\mathrm{disc}}(\alpha ,\gamma )=`$ (6)
$`{\displaystyle _{\mathrm{\Theta }_{\mathrm{out}}^{\mathrm{min}}}^{\mathrm{\Theta }_{\mathrm{out}}^{\mathrm{max}}}}{\displaystyle _{\varphi _{\mathrm{out}}^{\mathrm{min}}}^{2\pi \varphi _{\mathrm{out}}^{\mathrm{min}}}}P(\alpha ,\gamma ,\mathrm{\Theta }_{\mathrm{out}})\mathrm{sin}\mathrm{\Theta }_{\mathrm{out}}d\mathrm{\Theta }_{\mathrm{out}}d\varphi _{\mathrm{out}}`$
where $`\mathrm{\Theta }_{\mathrm{out}}^{\mathrm{max}}`$, $`\mathrm{\Theta }_{\mathrm{out}}^{\mathrm{min}}`$ and $`\varphi _{\mathrm{out}}^{\mathrm{min}}`$ are given by equations (2.1).
The viewing angle $`i`$ in this case cannot be identified with the direction of electron motion given by $`\alpha `$. Instead, we have a relation
$$\mathrm{cos}\mathrm{\Theta }_{\mathrm{out}}=(\mathrm{sin}i\mathrm{sin}\phi \mathrm{sin}\alpha +\mathrm{cos}i\mathrm{cos}\alpha ).$$
(7)
Here $`\alpha `$ and $`\phi `$ determine the electron’s velocity vector direction:
$$\stackrel{}{v}=v(\mathrm{cos}\phi \mathrm{sin}\alpha ,\mathrm{sin}\phi \mathrm{sin}\alpha ,\mathrm{cos}\alpha )$$
(8)
In Figure 3 we show the assumed geometry scheme. Note, that the angle $`\mathrm{\Theta }_{\mathrm{out}}`$ is measured from the vector $`\stackrel{}{v}`$ to the direction of emitted photon, and therefore we keep the integration limits given by equations (2.1) unchanged (see also Fig. 1 and Fig. 2 in Ghisellini et al. 1991).
We calculate the reflection amplitude averaged over the whole range of $`\alpha `$ and $`\phi `$ as a function of electron velocity and viewing angle, according to the formula
$$R(\gamma ,i)=\frac{1}{2\pi }\frac{_0^\pi _0^{2\pi }P_{\mathrm{disc}}(\alpha ,\gamma )\mathrm{sin}\alpha d\alpha d\phi }{_0^\pi _0^{2\pi }P(\alpha ,\gamma ,\mathrm{\Theta }_{\mathrm{out}}(\alpha ,\phi ,i))\mathrm{sin}\alpha d\alpha d\phi }.$$
(9)
In Figure 4 we plot the reflection amplitude as a function of $`\gamma `$ for three different values of viewing angle. The plot for an inclination $`i=60^o`$ roughly corresponds to the value averaged over all inclinations, as presented in Figure 1 for non-thermal plasma. We see that the anisotropy is only slightly reduced if the total collimation assumption is relaxed. This effect is mostly seen in Figure 2 where we present with the solid curve low velocity (moderate temperature) tail of the distribution.
### 2.3 Mildly relativistic bulk motion
In this section we assume that bulk velocity vector is perpendicular to the disc surface and directed outwards. We calculate the net electron velocity as a sum of thermal chaotic motion and systematic (bulk) outflow. The angle between the net velocity vector and vertical axis, $`\alpha ^{}`$, is connected with angle $`\alpha `$ via relativistic velocity transformation. Therefore the net Lorentz factor, $`\gamma ^{}`$, depends on the angle $`\alpha ^{}`$ as well as on bulk and thermal velocities. In this case the amount of reflection is given by:
$`R(\beta _{\mathrm{bulk}},\beta _{\mathrm{therm}},i)={\displaystyle \frac{1}{2\pi }}\times `$ (10)
$`{\displaystyle \frac{_0^\pi _0^{2\pi }P_{\mathrm{disc}}(\alpha ^{},\beta _{\mathrm{bulk}},\beta _{\mathrm{therm}})\mathrm{sin}\alpha ^{}d\alpha ^{}d\phi ^{}}{_0^\pi _0^{2\pi }P(\alpha ^{},\beta _{\mathrm{bulk}},\beta _{\mathrm{therm}},\mathrm{\Theta }_{\mathrm{out}}(\alpha ^{},i))\mathrm{sin}\alpha ^{}d\alpha ^{}d\phi ^{}}}.`$
In Figure 5 we plot the dependence of the amount of reflection on the bulk velocity to the light velocity ratio, $`\beta _{\mathrm{bulk}}`$, for different values of electron temperature $`\beta _{\mathrm{therm}}`$ and viewing angle $`i`$. The electron temperature corresponds to a single value of electron velocity, as for velocity distribution the analytic approximation would not work.
In the case of $`v_{\mathrm{therm}}=0`$ we obtain the same solution as in Beloborodov (1999). In Figure 6 we plot this solution, calculated for different viewing angles.
The comparison of Figures 6 and 5 shows that single scattering approach predicts very strong dependence of the amplitude of reflection on the plasma temperature due to anisotropy of the Compton scattering. When the thermal motions are important ($`v_{\mathrm{therm}}/c\beta _{\mathrm{bulk}}`$) the first reflection is significantly enhanced, by a factor of a two, and the higher the temperature the stronger the effect. The values of the amplitude of reflection cover the whole range between 0 and $`2`$ for outflow solutions ($`\beta _{\mathrm{bulk}}>0`$).
The effect of reflection enhancement drops rapidly when the bulk velocity is dominant ($`\beta _{\mathrm{bulk}}>v_{\mathrm{therm}}/c`$). In this case the maximum angle $`\alpha ^{}`$ is much smaller than $`\pi `$ and the integration limits in the equation ( 10) must be changed. The reflection amplitude increases then with the value of assumed viewing angle $`i`$, while in the case of more significant thermal motion the trend is opposite.
## 3 Reflected spectra from Monte Carlo simulations
The roughly power law shape of the primary emission component in X-ray spectra in GBH is most likely due to the effect of multiple scatterings within the hot plasma.
It is well known from analytical solutions and Monte Carlo simulations that the contribution of the first scattering to the total spectrum is rather specific (see Stern et al. 1995, Svensson 1996, Haardt, Maraschi & Ghisellini 1997). It means, that when the hard X-ray spectrum forms via multiple scatterings, only the first one is influenced by the anisotropy of seed photon distribution. This is the reason why any anisotropy effects are present in the low energy part of the spectrum. In the thermal medium, even for small optical depths, the power law spectrum is shaped by multiple scatterings. Therefore the semi-analytical computations, dealing with the first scattering process, in the case of both thermal and systematic bulk motion of the electrons within the corona can only serve as a guide and a help to understand the numerical results. The fully reliable answer can only be provided by full Monte Carlo simulations of the Comptonization process within the corona.
In this Section we compute the amplitude of the Compton reflected component using a Monte Carlo comptonization code and we show the corresponding spectra. We concentrate on highly ionized reflector, as the approximations used in Beloborodov (1999) may not be valid in that case. The code employs standard algorithms for simulating the inverse Compton scattering, and it was written following descriptions by Pozdnyakov, Sobol & Sunyaev (1983) and Górecki & Wilczewski (1984). Modifications to the code necessary to implement the bulk motion are described in Appendix A. Following Beloborodov (1999) we assume that the comptonizing region as a whole is stationary and its geometry is that of a slab.
### 3.1 Reflection amplitude
The bulk velocity vector is assumed perpendicular to the plane of the disc but the direction of the velocity can be both outwards and towards the disc. The disc is a source of soft photons for the comptonization. Photons backscattered from the cloud form radiation illuminating the disc. The usual Compton reflection process is then simulated by another Monte Carlo routine (Życki & Czerny 1994), assuming the abundances given in Morrison & McCammon (1983). The opacities were computed using the code described in Done et al. (1992) for the ionization parameter $`\xi F_\mathrm{X}/(n_\mathrm{e}r^2)=10^4`$. Further scattering of the reflected photons in the hot comptonizing cloud is not considered, since the intercepted fraction would be geometry-dependent (e.g. factor $`\mu _\mathrm{s}`$ in Beloborodov 1999).
The reflection amplitude $`R`$ is defined here (cf. Beloborodov 1999 and Section 2.3), as the ratio of the energy integrated fluxes:
$$R(\beta _{\mathrm{bulk}})=\frac{F_{\mathrm{back}}(\beta _{\mathrm{bulk}})}{F_{\mathrm{dir}}(\beta _{\mathrm{bulk}},i)}.$$
(11)
Here $`F_{\mathrm{dir}}`$ is the flux directly escaping from the comptonizing cloud towards an observer, at the inclination angle $`i`$, and $`F_{\mathrm{back}}`$ is the flux directed towards the disc. We note that our definition of $`R`$ would only be equivalent to that used in spectral fitting, if the shapes of spectra used in models fit to the data were exactly the same as in our simulations (for both the primary and the reflected components). Since this is not necessarily the case in practice, translating an amplitude $`R`$ inferred from spectral modelling to a bulk velocity may not be accurate.
The parameters of the comptonizing cloud: the electron temperature $`kT_\mathrm{e}`$ and optical depth $`\tau _{\mathrm{es}}`$ are chosen so that the resulting comptonized spectra correspond to typical spectra of GBH. Figure 7 shows the resulting amplitude of the reprocessed component as a function of the plasma bulk velocity $`\beta _{\mathrm{bulk}}`$.
Comparison of the Figures 7 and 5 shows that the enhancement of the reflection due to the anisotropy of the first scattering is sharply reduced if the subsequent scatterings are taken into account. Only for rather high inclination angles is the obtained reflection amplitude larger than 1.0 for $`\beta _{\mathrm{bulk}}>0`$, as the larger the viewing angle, the fewer photons are able to escape from the slab towards the observer. This means, in particular, that bulk velocities directed towards the disc are still necessary in some cases in order to explain full range of the $`\mathrm{\Gamma }`$$`R`$ relation found by Zdziarski et al. (1999), as the enhancement of flux towards the disc due to plasma bulk motion is required to explain $`R>1`$ seen in some sources.
### 3.2 Radiation spectra of outflowing corona
In Figure 8 we present an example of the overall spectrum calculated from the model of the outflowing corona. We choose the following values for the model parameters: $`kT_\mathrm{e}=100`$ keV, $`\tau =0.8`$, $`kT_0=0.1`$ keV and the bulk velocity $`\beta _{\mathrm{bulk}}=0.3`$. Adopted soft photon temperature corresponds to a typical value for galactic black holes. Such a parameterization is convenient if we do not consider the energy balance within the corona. We assume that the corona is a continuous medium, i.e. we neglect the clumpiness of the corona described by the parameter $`\mu _\mathrm{s}`$ in the model of Beloborodov (1999) since we also neglect the secondary reprocessing of the reflected component through the corona.
The continuous line shows the radiation emitted towards an observer (averaged over the entire hemisphere, roughly corresponding to an inclination of 60) while the short- dashed line shows the component backscattered towards the disk. Since the first scattering dominates soft X-ray band for galactic sources, the backscattered radiation in this band is enhanced, as predicted by analytical results presented in Section 2. However, hard X-ray part is dominated by multiple scattered photons, anisotropy effect is smeared off and the backscattered component is not enhanced in this band. The net effect is therefore the systematic difference between the spectral slope of the back-scattered radiation and forward-scattered radiation. This effect was discussed for a corona without a bulk motion in a number of papers (e.g. Stern et al. 1995). Therefore, the continuum formed in the corona and emitted towards an observer is slightly curved, particularly in the soft X-ray band, instead of being a simple power law with a high energy cut-off, as frequently assumed in spectral analysis of the data.
The backscattered continuum is subsequently reflected by the disk surface which in our calculations is assumed to be ionized ($`\xi =10^4`$). We show this spectral component in the Figure 8.
The reflected component in the outflowing corona model is again partially reprocessed by the corona. The effect depends on the corona clumpiness since only a fraction of radiation $`\mu _\mathrm{s}`$ (following the notation of Beloborodov, 1999) would pass again through the hot plasma. In the present paper we neglect this secondary reprocessing since it is essential only if $`\mu _s`$ is close to 1 and the optical depth is close to 1. However, in detail modeling this effect should be rather taken into account.
## 4 Discussion
It is generally assumed that in galactic X-ray sources and active galactic nuclei soft X-ray radiation originates from the geometrically thin and optically thick accretion disc while the hard flux is produced by Comptonization in optically thin plasma outside the disc. Part of the hard X-ray flux that is directed towards the observer is detected in the form of power law continuum and the part of the flux scattered back to the disc surface produces so called ’reflection hump’ in the spectrum over 10 keV (Lightman & White 1988; Pounds et al. 1990; Done et al. 1992) as well as the iron $`K_\alpha `$ line near 6.4 keV (Życki & Czerny 1994).
The observed values of the amplitude of the reflection component cover broader range than initially expected ($`R1`$ for AGN in Pounds et al. 1990). For some X-ray sources the observed reflection seems to be weak, which was modelled either by the disruption of the inner disc and X-ray irradiation of its outer parts, or by the high ionization state of illuminated medium, which results in steepening of the reflected continuum and makes it indistinguishable from the primary power law (Ross et al. 1999).
Both interpretations can explain the observed correlation between $`R`$ and $`\mathrm{\Gamma }`$. In the model with a disrupted cold disc partially overlapping the innermost hot flow the observed correlation is due to radiative coupling of the two components, as the amount of overlap varies (Zdziarski et al. 1999). In the model with strongly ionized disc surface the correlation is governed by the variable optical depth of the ionized scattering layer (Nayakshin et al. 2000), which determines the effective albedo of the disc and the soft flux from thermalized fraction of the illuminating X-rays. Both models can explain the reflection amplitude values $`R1`$.
However, the strength of reflection was for some sources found to be $`R>1`$, which may either indicate that reflecting medium subtends a solid angle larger than $`2\pi `$, or that the radiation directed towards the disc is enhanced. The latter is possible when the scattering process is anisotropic or may be due to the velocity of systematic bulk motion directed towards the disc.
In this article we show that the anisotropy effect is important mostly for the first scattering and weakens in numerical Monte Carlo simulations performed for multiple scattering in the plasma parameterized by optical depth and electron temperature. The maximum amplitude reached in the case of highly ionized reflector is only $`R1.1`$. Therefore, the systematic mildly-relativistic motion towards the disc, or the ’coronal inflow’, is required to explain the higher values of $`R`$. On the other hand, in order to produce $`R>1`$ in the disrupted disc model it would be necessary to allow for reflection from outer regions of a thickened disc (e.g. Shakura & Sunyaev 1973) and/or from the dusty/molecular torus around central black holes of AGN. Since the absorption of irradiating X-rays below $`5`$ keV is reduced for a highly ionized disc surface, the reflected photons may contribute to the soft X-ray excesses observed in the spectra of GBH in the hard/low state (Cyg X-1; Di Salvo et al. 1999, and in preparation). We show that this contribution is higher when the radiation backscattered to the disc is enhanced as a result of the anisotropy of the first scattering.
The hard to low state transition characteristic for many accreting black hole systems is in the ’disc plus sphere’ model connected with the change of the inner radius of the cold disc. In the hard spectral state the cold disc is pushed outward while the hot inner plasma is responsible for hard X-ray radiation. In the soft state the thermal disc emission dominates as the cold disc extends almost to the marginally stable orbit. However, the physical mechanism of such behaviour is currently unclear. In the outflowing corona model the outflow velocity is the basic control parameter. The hard spectral state would then correspond to rapid coronal expansion, while the spectrum with dominating soft component would be produced during the vertical collapse of the hot gas. Again the physical mechanism of such dependence is unclear, and in particular it is unclear how the required changes of either $`R_{\mathrm{in}}`$ or $`\beta _{\mathrm{bulk}}`$ could be driven by changing accretion rate.
In conclusion, the basic predictions of all the proposed models are similar and only detailed computations of spectral and temporal behaviour and comparison with the high quality data, as expected from Chandra and XMM, may allow for a distinction between them.
## Acknowledgments
This work was supported in part by grant 2P03D01816, 2P03D01718 and 2P03D01519 of the Polish State Committee for Scientific Research.
## Appendix A The Monte Carlo code
Our initial comptonization code was written following closely descriptions given by Pozdnyakov, Sobol & Sunyaev (1983) and Górecki & Wilczewski (1984), and applied to model X-ray spectra of accretion disc with accreting advective corona by Janiuk, Życki & Czerny (2000). In order to include the effect of bulk motion of the plasma, we needed to make two major modifications.
Firstly, the average scattering cross section has to be modified. The escape probability of a photon is given by
$$P(d)=\mathrm{exp}\left(_0^dN_\mathrm{e}\sigma 𝑑l\right),$$
(12)
where $`N_\mathrm{e}=N(𝒗)d^3v`$ is the electron density, $`N(𝒗)`$ is the electron velocity distribution, $`d`$ is the distance to the cloud boundary along the direction of photon motion, $`𝛀`$, and
$$\sigma =\frac{1}{N_\mathrm{e}}N(𝒗)(1𝒗𝛀/c)\sigma (x)d^3v$$
(13)
is the scattering cross section averaged over $`N(𝒗)`$. Here
$$x=2\frac{h\nu }{m_\mathrm{e}c^2}\gamma (1𝒗𝛀/c)$$
(14)
is the energy of an incoming photon in the electron rest frame, $`\gamma `$ is the Lorentz factor and $`\sigma (x)`$ is the Klein-Nishina cross section. For an isotropic $`N(𝒗)`$ e.g. Maxwell distribution, $`\sigma `$ is a function of photon energy only (for a given $`kT_\mathrm{e}`$).
With the non-zero bulk velocity we cannot use the method presented by Pozdnyakov et al. (1983) to evaluate the integral in Eq. (13), since $`N(𝒗)`$ is no longer isotropic. The presence of the specific direction – the bulk velocity vector $`𝜷`$ – introduces an additional angular dependence of $`\sigma `$. We now have to compute the integral as in Eq. (13) but with $`𝒗`$ in the dot-product $`𝒗𝛀`$ replaced by total electron velocity, $`𝒖`$. Here $`𝒖`$ is the sum (in the sense of Lorentz transformation) of the thermal velocity and the bulk velocity. We compute this 3-D integral numerically. Introducing a coordinate system with the $`z`$-axis along the bulk velocity and the $`x`$-axis along the direction of photon motion we obtain $`𝛀=(\mathrm{sin}\theta ,0,\mathrm{cos}\theta )`$, where $`\theta `$ is the angle between $`𝛀`$ and $`𝜷`$. Applying the Lorentz transformation we obtain the electron velocity in the disc frame,
$$𝒖=(\frac{v_\mathrm{x}}{\gamma \left(1+\beta \frac{v_\mathrm{z}}{c}\right)},\frac{v_\mathrm{y}}{\gamma \left(1+\beta \frac{v_\mathrm{z}}{c}\right)},\frac{v_\mathrm{z}+\beta c}{1+\beta \frac{v_\mathrm{z}}{c}}),$$
(15)
with $`\gamma =(1\beta ^2)^{1/2}`$, and $`v_\mathrm{x},v_\mathrm{y}\mathrm{and}v_\mathrm{z}`$ – components of the electron’s thermal velocity. This enables to compute $`𝒖𝛀`$ and calculate the required integral, which is now additionally a function of photon’s direction of motion, $`\theta `$. We used the procedure quad3d from Press et al. (1992) to evaluate the integral.
The second modification concerns the scattering event. Since this is modeled in the electron rest frame, we introduced a pair of additional Lorentz transformations of a photon’s momentum: from the disc/corona frame to the frame comoving with the bulk velocity before simulating the scattering event, and the reverse transformation after the scattering event.
This paper has been processed by the authors using the Blackwell Scientific Publications style file.
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# GREEN’S FUNCTION MONTE CARLO APPROACH TO SU(3) YANG-MILLS THEORY IN (3+1)D
## 1 Introduction
Classical Monte Carlo simulations provide a very powerful and accurate method for the study of Euclidean lattice gauge theories. In the Hamiltonian formulation, on the other hand, the corresponding quantum Monte Carlo methods have been somewhat neglected. Here we present a study of SU(3) Yang-Mills theory in (3+1) dimensions, using the Green’s Function Monte Carlo approach.
Heys and Stump and Chin et al. pioneered the use of “Green’s Function Monte Carlo” (GFMC) or “Diffusion Monte Carlo” techniques in Hamiltonian LGT, in conjunction with a weak-coupling representation involving continuous gauge field link variables. This was successfully adapted to non-Abelian Yang-Mills theories, with no minus sign problem arising. In this representation, however, one is simulating the wave function in gauge field configuration space by a discrete ensemble or density of random walkers: it is not possible to determine the derivatives of the gauge fields for each configuration, or to enforce Gauss’s law explicitly. the ensemble always relaxes back to the ground state sector. In order to compute the string tension or mass gap, one must measure an appropriate correlation function, and estimate the mass gap as the inverse of the correlation length. We have introduced the ‘forward-walking’ technique, well-known in many-body theory, to measure the expectation values and correlation functions. The technique has been demonstrated for the cases of the transverse Ising model in (1+1)D, and the U(1) LGT in (2+1)D.
Here we apply the technique for the first time to a non-Abelian model, namely SU(3) Yang-Mills theory in (3+1)D. The ground state energy and Wilson loop values are calculated, and approximate values are extracted for the string tension in the weak-coupling regime. Comparisons are made with earlier calculations, where they are available.
## 2 Method
### 2.1 Lattice Hamiltonian
The Green’s Function Monte Carlo formalism has been adapted to SU(2) Yang-Mills theory by Chin, van Roosmalen, Umland and Koonin, and sketched for the SU(3) case by Chin, Long and Robson.
The SU(3) lattice Hamiltonian is given by
$$H=\frac{g^2}{2a}\{\underset{l}{}E_l^aE_l^a\frac{\lambda }{3}\underset{p}{}Tr(U_p+U_p^{})\}$$
(1)
where $`E_l^a`$ is a component of the electric field at link l, $`\lambda =6/g^4`$, the index $`a`$ runs over the 8 generators of SU(3), and $`U_p`$ denotes the product of four link operators around an elementary plaquette. We will work with the dimensionless operator
$$H=\frac{1}{2}\underset{l}{}E_l^aE_l^a\frac{\lambda }{6}\underset{p}{}Tr(U_p+U_p^{})$$
(2)
The link variables are elements of the group SU(3) in the fundamental representation
$$U=\mathrm{exp}(i\frac{1}{2}\lambda ^aA^a)$$
(3)
### 2.2 Green’s Function Monte Carlo method
The Green’s Function Monte Carlo method employs the operator $`\mathrm{exp}(\tau (HE))`$, i.e. the time evolution operator in imaginary time, as a projector onto the ground state $`|\psi _0`$:
$$|\psi _0\underset{\mathrm{\Delta }\tau 0,N\mathrm{\Delta }\tau \mathrm{}}{lim}e^{N\mathrm{\Delta }\tau (HE)}|\mathrm{\Phi }$$
(4)
where $`|\mathrm{\Phi }`$ is any suitable trial state. To procure some variational guidance, one performs a “similarity transformation” with the trial wave function $`\mathrm{\Phi }`$, and evolves the product $`\mathrm{\Phi }|\psi _0`$ in imaginary time. The heart of the procedure is the calculation of the matrix element corrersponding to a single small time step $`\mathrm{\Delta }\tau `$. Chin et al show that
$`𝐱^{}|\mathrm{\Phi }e^{\mathrm{\Delta }\tau (HE)}\mathrm{\Phi }^1|𝐱`$ $`=`$ $`{\displaystyle \underset{l}{}}U_l^{}|N\{\mathrm{exp}({\displaystyle \frac{1}{2}}\mathrm{\Delta }\tau E_l^aE_l^a)\mathrm{exp}[\mathrm{\Delta }\tau E_l^a(E_l^a\mathrm{ln}\mathrm{\Phi })]\}|U_l`$ (5)
$`\mathrm{exp}\{\mathrm{\Delta }\tau [E\mathrm{\Phi }^1H\mathrm{\Phi }(𝐱)]\}+O(\mathrm{\Delta }\tau ^2)`$
$``$ $`p(𝐱^{},𝐱)w(𝐱)+O(\mathrm{\Delta }\tau ^2)`$
where $`𝐱=\{U_l\}`$ denotes an entire lattice configuration of link fields.
The product $`\mathrm{\Phi }|\psi `$ is simulated by the density of an ensemble of random walkers. At the kth. step, the ‘weight’ of each walker at $`𝐱_k`$ is multiplied by $`w(𝐱_k)`$. The effect of $`p(𝐱_{𝐤+\mathrm{𝟏}},𝐱_𝐤)`$ is to alter each link variable $`U_l`$ in $`\{𝐱_𝐤\}`$ to $`U_l^{}`$ by a Gaussian random walk plus a “drift step” guided by the trial wave function:
$$U^{}=\mathrm{\Delta }UU_dU$$
(6)
where $`U_d=\mathrm{exp}[i\frac{1}{2}\lambda ^a(i\mathrm{\Delta }\tau E^a\mathrm{ln}\mathrm{\Phi })]`$ is the drift step, and $`\mathrm{\Delta }U`$ is an SU(3) group element randomly chosen from a Gaussian distribution around the identity, with variance $`\mathrm{\Delta }s^2=8\mathrm{\Delta }\tau `$ (i.e. $`\mathrm{\Delta }\tau `$ for each index a), where
$$\mathrm{\Delta }s^2\underset{a}{}A^aA^a=8\mathrm{\Delta }\tau ,$$
(7)
for small $`A^a`$.
The simulation is carried out for a large number of iterations $`\mathrm{\Delta }\tau `$, until an equilibrium distribution $`\mathrm{\Phi }|\psi _0`$ is reached. The energy E in (5) is adjusted after each iteration so as to maintain the total ensemble weight constant. The average value of E can then be taken as an estimate of $`E_0`$, the ground-state energy.
As time evolves, the weights of some walkers grow larger, while others grow smaller, which would produce an increased statistical error. To avoid this, a “branching” process is employed, whereby a walker with weight larger than some threshold is split into two independent walkers, while others with weights lower than another threshold are amalgamated.
### 2.3 Trial Wave Function
The trial wave function is chosen to be the one-parameter form
$$\mathrm{\Phi }=\mathrm{exp}[\alpha \underset{p}{}Tr(U_p+U_p^{})]$$
(8)
Then the drift step for each linkis
$$U_d=\mathrm{exp}[i\frac{\lambda ^a}{2}A_l^a],$$
(9)
$$A_l^a=i\mathrm{\Delta }\tau \frac{\alpha }{2}\underset{pl}{}Tr[\lambda ^aU_l..U_4^{}h.c.]$$
(10)
Finally, the trial energy factor is
$`\mathrm{\Phi }^1H\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \underset{l}{}}\{{\displaystyle \frac{\alpha ^2}{8}}({\displaystyle \underset{pl}{}}Tr[\lambda ^aU_l..U_4^{}h.c.])^2`$ (12)
$`+({\displaystyle \frac{2\alpha }{3}}{\displaystyle \frac{\lambda }{24}}){\displaystyle \underset{pl}{}}Tr(U_p+U_p^{})\}.`$
### 2.4 Forward Walking estimates
The “forward walking” technique is used to estimate expectation values. Its application to the U(1) lattice gauge theory in (2+1)D was discussed by Hamer et al. It is implemented for an operator Q (assumed diagonal, for simplicity) by recording the value $`Q(𝐱_i)`$ for each “ancestor” walker at the beginning of a measurement; propagating the ensemble as normal for $`J`$ iterations, keeping a record of the “ancestor” of each walker in the current population; and taking the weighted average of the $`Q(𝐱_i)`$ with respect to the weights of the descendants of $`𝐱_i`$ after the $`J`$ iterations, using sufficient iterations $`J`$ that the estimate reaches a ‘plateau’.
## 3 Results
Simulations were carried out for LxLxL lattices up to L=8 sites, using runs of typically 4000 iterations and an ensemble size of 250 to 1000 depending (inversely) on lattice size. Time steps $`\mathrm{\Delta }\tau `$ of 0.01 and 0.05 “seconds” were used, with each iteration consisting of 5 sweeps and 1 sweep through the lattice, respectively, followed by a branching process. The first 400 iterations were discarded to allow for equilibration. The data were block averaged over blocks of up to 256 iterations, to minimize the effect of correlations on the error estimates.
The results taken at $`\mathrm{\Delta }\tau =0.01`$ and $`\mathrm{\Delta }\tau =0.05`$ were extrapolated linearly to $`\mathrm{\Delta }\tau =0`$. The variational parameter $`c`$ was given values as used by Chin et al, obtained from a variational Monte Carlo calculation. We checked that these were approximately the optimum values for small lattices.
Forward-walking measurements were taken over $`J`$ iterations, where $`J`$ ranged from 20 to 100, depending on the coupling $`\lambda `$. Ten separate measurements were taken over this time interval, in order to check whether the value measured by forward-walking had reached equilibrium. A new measurement was started soon after the previous one had finished.
### 3.1 Ground-state Energy
The dependence of the ground-state energy per site on lattice size is illustrated in Figure 1, at two fixed couplings $`\lambda =3.0`$ and $`\lambda =5.0`$. In the “strong-coupling” case, $`\lambda =3.0`$, it can be seen that the results converge exponentially fast in $`L`$, whereas in the “weak-coupling” regime, $`\lambda =5.0`$, the convergence is more like $`1/L^4`$ at these lattice sizes. This behaviour merits some further explanation.
A similar phenomenon occurs in the case of the U(1) theory in (2+1)D. In the strong-coupling regime, where the mass gap is large, the usual exponential convergence occurs. In the weak-coupling regime, however, where the mass gap M is very small, the finite-size scaling behaviour for small lattice sizes is that of a massless theory, and it is only at much larger lattice sizes $`L1/M`$ that a crossover to exponential convergence occurs. An “effective Lagrangian” corresponding to free, massless gluons (non-interacting QCD) should describe the finite-size behaviour in the present case, in line with the idea of asymptotic freedom. By analogy with the (2+1)D case, we expect a $`1/L^4`$ dependence for the corrections to the ground-state energy per site. We hope to pursue this analysis further at a later date.
An anomalous feature in Figure 3b) is that the $`L=8`$ point lies well out of line with the others. This occurs at other couplings also. We suspect that the results for $`L=8`$ are not reliable, and that the trial wave function will have to be further improved to give reliable results for such large lattices.
We have made estimates of the bulk limit, extrapolating mainly from the smaller L values where possible. The estimates for the bulk ground-state energy per site are graphed as a function of coupling in Figure 2, where they are compared with previous estimates obtained by an ‘Exact Linked Cluster Expansion’ (ELCE) procedure, and with the asymptotic weak-coupling series. The Monte Carlo results agree very well with the ELCE estimates, and appear to match nicely onto the expected weak-coupling behaviour for $`\lambda 6`$.
### 3.2 Wilson Loops
The forward-walking method was used to estimate values for the m x n Wilson loops, $`W(m,n)`$. A graph of the ‘mean plaquette’ $`W(1,1)`$ versus the variational parameter $`c`$ is shown in Figure 3. A problem is immediately apparent. The estimate for $`W(1,1)`$ is not independent of $`c`$, in fact it depends linearly on $`c`$ over this range, and the size of the variation is such that the probable systematic error due to the choice of $`c`$ is an order of magnitude larger than the random statistical error in the results. Thus it would be advantageous in future studies to put more effort into improving the trial wave function, rather than merely improving the statistics.
The finite-size behaviour for the Wilson loops is similar to that of the ground-state energy. The estimates for the mean plaquette in the bulk limit are graphed as a function of coupling $`\lambda `$ in Figure 4, and compared with series estimates at strong and weak coupling. The agreement is quite good.
### 3.3 String Tension
Having obtained estimates for the Wilson loop values on the bulk lattice, one can extract estimates for the ‘spacelike’ string tension using the Creutz ratios:
$$Ka^2R_n=\mathrm{ln}\left[\frac{W(n,n)W(n1,n1)}{W(n,n1)^2}\right]$$
(13)
The results are shown in Figure 5. Also shown in Figure 5 are some previous estimates derived from the ‘axial’ string tension, obtained using an ‘Exact Linked Cluster Expansion’ (ELCE) method. The axial string tension $`aT`$ is calculated as an energy per link, and must be converted to a dimensionless, ‘spacelike’ tension by dividing by the ‘speed of light’ c. We have also used the weak-coupling relationship between the scales of Euclidean and Hamiltonian lattice Yang-Mills theory calculated by Hasenfratz et al to plot the results against the Euclidean coupling $`\beta =6/g_E^2`$.
It can be seen that the present GFMC results are in rough agreement with the axial string tension results in the region $`4\beta 5`$, which is also the region where the ‘roughening’ transition occurs in the string tension. For $`\beta >5`$, however, the Creutz ratio $`R_2`$ runs above the ELCE estimate, and shows no sign of the expected crossover to an exponentially decreasing scaling behaviour at $`\beta 6`$. We presume that this is a finite-size effect, and that the Creutz ratios $`R_n`$ for larger $`n`$ will show a substantial decrease in the ‘weak-coupling’ regime $`\beta 6`$. That is certainly the pattern seen in the Euclidean calculations, or in the U$`(1)_{2+1}`$ model. Unfortunately, however, our present results for the larger Wilson loops are not of sufficient accuracy to allow worthwhile estimates of $`R_n`$ for $`n2`$.
## 4 Conclusions
Some significant problems with the GFMC method have emerged from this study. The ‘forward-walking’ technique was introduced specifically to avoid any variational bias from the trial wave function. As it turns out, however, the results for the Wilson loops show a substantial dependence on the trial wave function parameter $`c`$. The systematic error due to this dependence is an order of magnitude larger than the statistical error, so it would pay to put more effort in future studies into improving the trial wave function, rather than simply increasing the statistics. Furthermore, the effective ensemble size decreases during each measurement as the descendants of each ‘ancestor’ state die out, and this produces a substantial loss in statistical accuracy at weak coupling, as well.
It would be preferable if one were able to do away entirely with all the paraphernalia of trial wave function, weights, branching algorithms, etc, and just rely on some sort of Metropolis-style accept/reject algorithm to produce a correct distribution of walkers. Within a quantum Hamiltonian framework, a way is known to do this, namely the Path Integral Monte Carlo (PIMC) approach. We conclude that the PIMC approach may be better suited than GFMC to the study of large and complicated lattice Hamiltonian systems.
## Acknowledgments
This work is supported by the Australian Research Council. Calculations were performed on the SGI Power Challenge Facility at the New South Wales Centre for Parallel Computing and the Fujitsu VPP300 vector machine at the Australian National Universtiy Supercomputing Facility: we are grateful for the use of these facilities.
## References
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# 1 Introduction
## 1 Introduction
It has long been suggested that the description of spacetime based on the usual notion of geometry may not be valid at the Planck scale, and perhaps the spacetime becomes noncommutative or may show non-Archimedean structure at such small length scale. It has therefore been believed that the noncommutative description of spacetime might be relevant to quantum theory of gravity. It is a generic property of a noncommutative space that the notion of a point has no meaning but the lattice structure of spacetime emerges due to the uncertainty in the measurement of the particle position in space. So, in some cases such lattice structure of spacetime at the Planck scale eliminates ultraviolet divergence problem. Recently, there has been active investigation of noncommutative theories after it was found out that noncommutative spacetime emerges naturally in M-theory compactified in the presence of constant background three-form field and in the worldvolume theory of D-brane with nonzero constant NS $`B`$-field .
It is the purpose of this paper to study the noncommutative generalization of the conformal quantum mechanics of de Alfaro, Fubini and Furlan . In Ref. , it is observed that such conformal quantum mechanics model can be realized as a special limit of the mechanics of a massive charged point particle in the near horizon background of the extremal Reissner-Nordström black hole, indicating the possible relevance of the conformal quantum mechanics model to the quantum theory of black holes. Furthermore, the fact that the $`SO(1,2)SU(1,1)`$ isometry symmetry of the AdS$`{}_{2}{}^{}\times S^n`$ near-horizon geometry of the extremal Reissner-Nordström black hole coincides with the $`SL(2,𝐑)SU(1,1)`$ symmetry of the conformal quantum mechanics indicates that the conformal quantum mechanics may have some relevance to the poorly understood AdS$`{}_{2}{}^{}/`$CFT<sub>1</sub> duality.
In the case of one spatial dimension, it is unclear what is meant by noncommutative space (in this paper we assume that the time coordinate is a commutative variable), since a coordinate always commutes with itself. A possible way of introducing noncommutativity for such case is by following the Manin’s proposal , which is based upon a differential calculus on quantum plane , that the Heisenberg algebra can be modified through q-deformation of the phase space. The Manin’s proposal was first applied to some nonrelativistic dynamical system in Refs. . Following this line of approach, we shall construct a q-deformed version of the conformal quantum mechanics model of Ref. and study its properties.
The paper is organized as follows. In section 2, we review the relevant aspect of the conformal quantum mechanics. In section 3, we discuss the q-deformed Heisenberg algebra with a complex deformation parameter and with unitary time evolution of the system, and apply this to construct a q-deformed version of the conformal quantum mechanics. In section 4, we study dynamics of the system and from this we construct differential calculus on the q-deformed phase space.
## 2 Conformal Quantum Mechanics
In the following, we briefly summarize the conformal quantum mechanics studied in Refs. . The Lagrangian density for the system is given by
$$=\frac{1}{2}m\dot{x}^2\frac{g}{2x^2}.$$
(1)
The action is invariant under the following $`SL(2,𝐑)`$ conformal algebra, spanned by the Hamiltonian $``$, the dilatation generator $`𝒟`$ and the special conformal generator $`𝒦`$:
$$[𝒟,]=2i,[𝒟,𝒦]=2i𝒦,[,𝒦]=i𝒟.$$
(2)
Here, the $`SL(2,𝐑)`$ generators are explicitly given by
$$=\frac{p^2}{2m}+\frac{g}{2x^2},𝒟=\frac{1}{2}(px+xp),𝒦=\frac{1}{2}mx^2.$$
(3)
The problem with the above Hamiltonian $``$ is that its eigenspectrum is continuous and bounded from below but without an endpoint or ground state, and its eigenstates are not normalizable. Such problem was circumvented by redefining the Hamiltonian as a linear combination of the above $`SL(2,𝐑)`$ generators with a suitable condition on the coefficients. Particularly, the following choice is found to be convenient :
$$L_0=\frac{1}{2}\left(a+\frac{1}{a}𝒦\right),$$
(4)
where the introduction of the constant $`a`$ leads to breakdown of scale invariance. Then, the potential term in $`L_0`$ has the minimum and the energy eigenstates become discrete and normalizable. Furthermore, along with the following linear combinations:
$$L_{\pm 1}=\frac{1}{2}\left(a\frac{1}{a}𝒦i𝒟\right),$$
(5)
$`L_0`$ satisfies the following $`SL(2,𝐑)`$ algebra in the Virasoro form:
$$[L_1,L_1]=2L_0,[L_0,L_{\pm 1}]=L_{\pm 1}.$$
(6)
## 3 q-Deformation of Conformal Quantum Mechanics
For the ordinary commutative quantum mechanics in two spacetime dimensions, the Heisenberg algebra of observables can be defined as the following quotient:
$$H(I,x,p)=C[I,x,p]/J(I,x,p),$$
(7)
where $`C[I,x,p]`$ is a unital associative algebra freely generated by the identity $`I`$, the position operator $`x`$ and the canonical momentum operator $`p`$, and $`J(I,x,p)`$ is a two-sided ideal in $`C[I,x,p]`$ generated by the following relation corresponding to the Heisenberg rule:
$$xppx=iI,$$
(8)
where we are using the unit in which $`\mathrm{}=1`$. The operators $`x`$ and $`p`$ are assumed have the following property under the antilinear anti-involution operation in $`C[I,x,p]`$:
$$x^{}=x,p^{}=p.$$
(9)
The formalism of Manin’s quantum space can be applied to the above Heisenberg algebra by making use of the q-deformed differential calculus developed in Ref. . Namely, one can deform the above Heisenberg algebra by deforming the usual Heisenberg rule (8) as follows:
$$xpqpx=iI,$$
(10)
where the deformation parameter $`q`$ can be either complex or real. First, if $`q`$ is a complex number, the consistency of the relation (10) along with the Hermiticity condition (9) on $`x`$ and $`p`$ requires that $`|q|=1`$. According to Ref. , which first studied the q-deformed classical and quantum mechanics (with a complex deformation parameter $`q`$) of a particle in one-dimensional space and whose work is later generalized to the relativistic case in Ref. , the parameters of the dynamics such as the inertial mass $`m`$ of the particle do not commute with the generators $`x`$ and $`p`$ of the algebra and there is no unitary time evolution of the system at the quantum level. Later, it is found out that to achieve unitary noncommutative q-dynamics on the quantum level, i.e. for the Heisenberg equation of motion $`\dot{\mathrm{\Omega }}=\frac{i}{\mathrm{}}[,\mathrm{\Omega }]+_t\mathrm{\Omega }`$ to be satisfied after the q-deformation, one has to introduce additional generators into the algebra. Second, if $`q`$ is a real number, $`x`$ and $`p`$ cannot both be Hermitian, as can be seen by applying the involution operation to Eq. (10). So, one has to assume that only one of $`p`$ and $`x`$ is Hermitian and the involution of the other is a separate operator . One can alternatively describe the q-deformed Heisenberg algebra with a real $`q`$ by redefining the generators, say, in terms of the above generators $`p`$, $`x`$ and $`x^{}`$ so that new momentum and position operators can be both hermitian, as was done in Ref. . In such case, an additional generator (expressed in terms of $`p`$, $`x`$ and $`x^{}`$), which approaches $`I`$ as $`q1`$, other than the hermitian position and momentum operators, is introduced into the algebra. Such alternative algebra can be obtained also by making use of the Leibniz rule $`_𝗑𝗑=1+q𝗑_𝗑`$ for the differential calculus in the one-dimensional q-deformed Euclidean space $`𝐑_q^1`$. In the present paper, we shall apply the first approach for studying the q-deformed generalization of the conformal quantum mechanics of Refs. .
The q-deformed Heisenberg algebra of observables with a complex deformation parameter is given by the following quotient:
$$H=A[I,x,p,K,\mathrm{\Lambda }]/J(I,x,p,K,\mathrm{\Lambda }).$$
(11)
In the case of a particle under the influence of non-trivial potential $`V`$, with the assumption of the proper limit of no q-deformation, the two-sided ideal $`J`$ is defined by the following q-deformed Heisenberg relations or the Bethe Ansatz re-ordering rules:
$`xp`$ $`=`$ $`q^2px+iq\mathrm{\Lambda }^2,x\mathrm{\Lambda }=\xi \mathrm{\Lambda }x,p\mathrm{\Lambda }=\xi ^1\mathrm{\Lambda }p,`$ (12)
$`xK`$ $`=`$ $`\xi ^2Kx,pK=Kp,\mathrm{\Lambda }K=\xi ^1K\mathrm{\Lambda }.`$ (13)
Here, the generators $`K`$ and $`\mathrm{\Lambda }`$ are assumed to be invertible and time-independent, and one can consistently (with the above q-deformed Heisenberg relations) impose the following reality conditions on the generators under the involution operation:
$$x^{}=x,p^{}=p,K^{}=K,\mathrm{\Lambda }^{}=\mathrm{\Lambda },$$
(14)
along with $`|q|=1=|\xi |`$.
In the q-deformed quantum phase space described in the above, the Hamiltonian (3) of the conformal quantum mechanics of Ref. is deformed in the following way:
$$_\xi =p^2K^2+\frac{mg}{\xi ^4}x^2K^2\mathrm{\Lambda }^4,$$
(15)
where we obtained this form of Hamiltonian from the requirement of the consistency of the Hamiltonian form of the Heisenberg equations with the q-deformed Heisenberg relations (13), namely the requirement of the unitary time evolution of the system. To further impose the naturalness condition that the velocity $`\dot{x}`$ is linear in the momentum $`p`$ in the Heisenberg equation of motion $`\dot{x}=i[_\xi ,x]`$, one has to further let $`\xi =q`$, which we assume from now on. Note, in the limit of no q-deformation, $`K`$ and $`\mathrm{\Lambda }`$ belong to the center of the algebra. The requirement of irreducibility of the representation level implies that they should be proportional to the identity $`I`$ when $`\xi =q=1`$. We choose $`K=\frac{1}{\sqrt{2m}}I`$ and $`\mathrm{\Lambda }=I`$ when $`\xi =q=1`$ so that the Hamiltonian (15) reduces to the form (3) in the limit of no q-deformation.
The dilatation generator $`𝒟`$ and the special conformal generator $`𝒦`$ in Eq. (3) of the $`SL(2,𝐑)`$ algebra can be q-deformed in such a way that the commutation relations (2) of the $`SL(2,𝐑)`$ algebra continue to be satisfied after the q-deformation. Such q-deformed $`SL(2,𝐑)`$ generators are given by
$`_q`$ $`=`$ $`p^2K^2+{\displaystyle \frac{mg}{q^4}}x^2K^2\mathrm{\Lambda }^4,`$ (16)
$`𝒟_q`$ $`=`$ $`{\displaystyle \frac{1}{2}}(qpx+q^1xp)\mathrm{\Lambda }^2,`$ (17)
$`𝒦_q`$ $`=`$ $`{\displaystyle \frac{1}{4q^4}}x^2K^2\mathrm{\Lambda }^4.`$ (18)
By using the q-deformed Heisenberg relations (13) with $`\xi =q`$, one can show that these q-deformed generators satisfy the following commutation relations:
$$[𝒟_q,_q]=2i_q,[𝒟_q,𝒦_q]=2i𝒦_q,[_q,𝒦_q]=i𝒟_q.$$
(19)
Just as in the case of undeformed conformal quantum mechanics, one can redefine the generators through the linear combinations so that the resulting new generators satisfy the $`SL(2,𝐑)`$ algebra in the Virasoro form. The q-deformed forms of the $`SL(2,𝐑)`$ generators (4) and (5) are given by
$`L_0^q`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(a_q+{\displaystyle \frac{1}{a}}𝒦_q\right),`$ (20)
$`L_{\pm 1}^q`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(a_q{\displaystyle \frac{1}{a}}𝒦_qi𝒟_q\right),`$ (21)
where $`_q`$, $`𝒟_q`$ and $`𝒦_q`$ are defined in Eq. (18). It is straightforward to show that the q-deformed generators (21) still satisfy Eq. (6). It might be possible to construct a q-deformed version of the de Alfaro, Fubini and Furlan conformal quantum mechanics in such a way that the generators instead satisfy the q-deformed commutation relations and thereby the $`SL(2,𝐑)`$ algebra (6) is deformed to the quantum $`SL_q(2,𝐑)`$ algebra. The q-deformed version of the symmetry group, the so-called quantum group, of a dynamical system was originally studied in Refs. within the context of the quantum noncommutative harmonic oscillator with q-deformed creation and annihilation operators.
## 4 Differential Calculus on the q-Deformed Phase Space
Note, the above q-deformed algebra (11) generated by $`I`$, $`x`$, $`p`$, $`K`$ and $`\mathrm{\Lambda }`$, satisfying the commutation relations (13), is the zero-form sector of the q-deformed quantum de Rham complex generated by these generators and their differentials. The quantum de Rham complex contains information about not only algebra of observables but also dynamics of theory. Namely, by relating the velocity vector $`(\dot{x},\dot{p})`$ for a particle in the quantum phase space to the one forms $`dx`$ and $`dp`$ as $`dx=\dot{x}dt`$ and $`dp=\dot{p}dt`$, one can learn about the dynamics of a particle moving on the quantum phase space from the commutation relations among the generators of the algebra and their differentials. Here, the dot denotes the derivative with respect to the time coordinate $`t`$, which we assume to be a commuting parameter. According to Ref. , there are three families <sup>2</sup><sup>2</sup>2In the case of the relativistic motion of a particle in the two-dimensional noncommutative Minkowski spacetime, one of the families is excluded and the remaining two coincide under the condition of the reasonable description of the particle dynamics. of possible differential calculi associated with the Manin’s plane defined by the commutation relations among the generators. In this section, we construct the q-deformed quantum de Rham complex directly from the Heisenberg equations of motion, instead of applying the result of Refs. . In the following, we restore the Planck’s constant $`\mathrm{}`$ in the equations just for the purpose of making it easy to see various limits. In particular, in the $`q1`$ and $`\mathrm{}0`$ limit (i.e., the limit of undeformed classical theory), the formulae obtained in the following reduce to those of the usual commutative classical geometry.
The Heisenberg equations associated with the q-deformed Hamiltonian (15) with $`\xi =q`$ have the following form <sup>3</sup><sup>3</sup>3For more general $`\xi q`$ case, the Heisenberg equations take the following form: $`\dot{x}`$ $`=`$ $`\left[{\displaystyle \frac{i}{\mathrm{}}}(\xi ^4q^4)p^2x+q(q^2+\xi ^2)p\mathrm{\Lambda }^2\right]K^2,`$ (22) $`\dot{p}`$ $`=`$ $`{\displaystyle \frac{img}{\mathrm{}q^4}}\left[\left({\displaystyle \frac{q}{\xi }}\right)^41\right]px^2K^2\mathrm{\Lambda }^4+{\displaystyle \frac{mg}{q^4}}\left[\left({\displaystyle \frac{q}{\xi }}\right)^2+1\right]x^3K^2\mathrm{\Lambda }^6.`$ (23) As mentioned in the previous section, the velocity $`\dot{x}`$ becomes linear in the momentum $`p`$ when $`\xi =q`$.:
$$\dot{x}=\frac{i}{\mathrm{}}[_q,x]=2qpK^2\mathrm{\Lambda }^2,\dot{p}=\frac{i}{\mathrm{}}[_q,p]=\frac{2mg}{q^4}x^3K^2\mathrm{\Lambda }^6.$$
(24)
One can also show by using Eq. (13) that $`\dot{K}=\frac{i}{\mathrm{}}[_q,K]=0`$ and $`\dot{\mathrm{\Lambda }}=\frac{i}{\mathrm{}}[_q,\mathrm{\Lambda }]=0`$. The Aref’eva-Volovich limit is achieved by further letting $`\mathrm{\Lambda }=I`$. In this case, the system does not evolve unitarily with time, as can be seen from the fact that the re-ordering rules (13), which are derived from the condition of unitary time evolution, cannot be consistent when $`\mathrm{\Lambda }=I`$. As expected, in the limit of no q-deformation, the above Heisenberg equations of motion reduce to the Heisenberg equations associated with the Hamiltonian given by Eq. (3).
From the above Heisenberg equations, one can express the differentials of the generators as follows:
$`dx`$ $`=`$ $`\dot{x}dt=2qpK^2\mathrm{\Lambda }^2dt,dp=\dot{p}dt={\displaystyle \frac{2mg}{q^4}}x^3K^2\mathrm{\Lambda }^6dt,`$ (25)
$`dK`$ $`=`$ $`\dot{K}dt=0,d\mathrm{\Lambda }=\dot{\mathrm{\Lambda }}dt=0.`$ (26)
By using the fact that we defined the time coordinate $`t`$ to be a commuting parameter, one can derive the commutation relations among the generators and their differentials. The following commutation relations can be obtained by making use of the relations (26) and the q-deformed Heisenberg relations (13) with $`\xi =q`$:
$`xdx`$ $`=`$ $`dx(x+i\mathrm{}q^1p^1\mathrm{\Lambda }^2),`$ (27)
$`pdx`$ $`=`$ $`q^2dxp,`$ (28)
$`Kdx`$ $`=`$ $`q^2dxK,`$ (29)
$`\mathrm{\Lambda }dx`$ $`=`$ $`q^1dx\mathrm{\Lambda },`$ (30)
$`xdp`$ $`=`$ $`q^2dpx,`$ (31)
$`pdp`$ $`=`$ $`dp(p+3i\mathrm{}qx^1\mathrm{\Lambda }^2),`$ (32)
$`Kdp`$ $`=`$ $`dpK,`$ (33)
$`\mathrm{\Lambda }dp`$ $`=`$ $`qdp\mathrm{\Lambda }.`$ (34)
The first and the sixth relations in Eq. (34) can be rewritten in more symmetric forms as follows:
$$pxdx=q^4dxxp,x^3K^2\mathrm{\Lambda }^4pdp=q^4dppx^3K^2\mathrm{\Lambda }^4.$$
(35)
By assuming the usual Leibniz rule and nilpotency condition for the external differential operator $`d`$, one obtains the following product rules for the differentials:
$`(dx)^2`$ $`=`$ $`{\displaystyle \frac{i\mathrm{}}{2}}q^7p^2dxdp,`$ (36)
$`(dp)^2`$ $`=`$ $`{\displaystyle \frac{3}{2}}i\mathrm{}qdxdp(_xx^1)\mathrm{\Lambda }^2,`$ (37)
$`dxdp`$ $`=`$ $`q^2dpdx,`$ (38)
where $`_x`$ denotes the generalized q- and $`\mathrm{}`$-deformed partial derivative with respect to $`x`$, which we define in the following. We see that the first commutation relation in Eq. (34) is similar to that of the universal calculus over a lattice, except that the ‘lattice spacing’ is an operator. In fact, the q-deformation leads to deformation of continuous phase space to a lattice structure, where the Hilbert space of representations of the q-deformed system has a discrete spectrum, putting physics on a q-lattice.
We further enlarge the algebra by defining the derivatives $`_x`$ and $`_p`$ on the q-deformed phase space in the following way:
$$d=dx_x+dp_p,$$
(39)
along with the assumption of the usual Leibniz rule and the nilpotency condition as above. Note, we have seen in the above that $`dK=0=d\mathrm{\Lambda }`$. The following commutation relations between the partial derivatives and the generators can be obtained by applying the Leibniz rule:
$`_xx`$ $`=`$ $`1+(x+i\mathrm{}q^1p^1\mathrm{\Lambda }^2)_x,`$ (40)
$`_xp`$ $`=`$ $`q^2p_x,`$ (41)
$`_xK`$ $`=`$ $`q^2K_x,`$ (42)
$`_x\mathrm{\Lambda }`$ $`=`$ $`q^1\mathrm{\Lambda }_x,`$ (43)
$`_px`$ $`=`$ $`q^2x_p,`$ (44)
$`_pp`$ $`=`$ $`1+(p+3i\mathrm{}qx^1\mathrm{\Lambda }^2)_p,`$ (45)
$`_pK`$ $`=`$ $`K_p,`$ (46)
$`_p\mathrm{\Lambda }`$ $`=`$ $`q\mathrm{\Lambda }_p,`$ (47)
The commutation relations between the partial derivatives and the differentials can be obtained by demanding their consistency with the q-deformed Heisenberg rules and the product rules obtained in the above. We have not yet been successful in obtaining the commutation relations for the general $`\mathrm{}0`$ case. In the q-deformed classical phase space (the $`\mathrm{}=0`$ case), the commutation relations are given by
$$_xdx=dx_x,_xdp=q^2dp_x,_pdx=q^2dx_p,_pdp=dp_p.$$
(48)
In the $`q1`$ and $`\mathrm{}0`$ limit (i.e., the limit of undeformed classical theory), the formulae obtained in the above reduce to those of the usual commutative classical geometry. Particularly interesting limits are the $`q1`$ limit and the $`\mathrm{}0`$ limit, which respectively correspond to the $`\mathrm{}`$-deformation (or quantization) and the q-deformation of the differential calculus on the “classical” commutative phase space. Note, the q-deformation and the $`\mathrm{}`$-deformation generally do not commute with one another, i.e., the so-called Faddeev’s rectangle is not always commutative. In this paper, we consider the case of the q-deformation of the commutative quantized (or $`\mathrm{}`$-deformed) classical conformal mechanics. Had we first q-deformed the commutative classical conformal quantum mechanics and then quantized it, we might have obtained different noncommutative theory.
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# Spin structure of the octet baryons
## I Introduction
Since the European Muon Collaboration (EMC) measured the first moment of the proton spin structure function $`g_1^p`$ , there has been a great deal of discussion about the spin content of the proton. A series of following experiments confirmed the EMC measurement. In contrast to the result from the naive nonrelativistic quark model, which is reflected in the Ellis-Jaffe sum rule , the strange quark contribution to the nucleon spin deviates from zero. The global fit performed by Ellis and Karliner gives the following value $`\mathrm{\Delta }s=0.11\pm 0.03`$. For more recent analysis, see Refs.. These results, however, are obtained with an assumption of the exact SU(3) symmetry for the baryon semileptonic decays.
Great efforts have been already spent on understanding the spin and flavor content of the proton (see for review and recent papers ). While it is now known that a large fraction of the nucleon spin is provided by gluons and their orbital angular momenta, it is still very important to understand the mechanism of how the quarks carry the nucleon spin. In particular, since the extraction of the flavor content of the nucleon spin relies on the empirical data of the baryon semileptonic decays, it is of great significance to examine the influence of the SU(3) symmetry breaking on the axial properties of the baryons in a consistent way.
One piece of information comes from the first moment of the spin structure function $`g_1^\mathrm{p}(x)`$ of the proton:
$$I_\mathrm{p}=\underset{0}{\overset{1}{}}𝑑xg_1^\mathrm{p}(x)=\frac{1}{18}\left(4\mathrm{\Delta }u_\mathrm{p}+\mathrm{\Delta }d_\mathrm{p}+\mathrm{\Delta }s_\mathrm{p}\right)\left(1\frac{\alpha _\mathrm{s}}{\pi }+\mathrm{}\right).$$
(1)
The analysis of Karliner and Lipkin implies $`I_\mathrm{p}=0.124\pm 0.011`$ which can be translated into:
$$\mathrm{\Gamma }_\mathrm{p}4\mathrm{\Delta }u_\mathrm{p}+\mathrm{\Delta }d_\mathrm{p}+\mathrm{\Delta }s_\mathrm{p}=2.56\pm 0.23.$$
(2)
if $`\alpha _\mathrm{s}(Q^2=3(\mathrm{GeV}/c)^2)=0.4`$ is assumed. Let us for completeness quote also the result for the neutron:
$$\mathrm{\Gamma }_n4\mathrm{\Delta }d_\mathrm{p}+\mathrm{\Delta }u_\mathrm{p}+\mathrm{\Delta }s_\mathrm{p}=0.928\pm 0.186$$
(3)
where the isospin symmetry (Bjorken sum rule) has been assumed.
Another piece of information comes from the semileptonic decays, which in the case of the exact SU(3) symmetry can be parametrized by two reduced matrix elements $`F`$ and $`D`$. Taking for $`F=0.46`$ and for $`D=0.80`$ together with Eq.(2), one gets for the proton: $`\mathrm{\Delta }u_\mathrm{p}=0.79`$, $`\mathrm{\Delta }d_\mathrm{p}=0.47`$ and $`\mathrm{\Delta }s_\mathrm{p}=0.13`$, which implies $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{p}=0.19`$, quite a small number as compared with the naive expectation from the quark model: $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{p}=1`$.
It is important to realize that $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{p}`$ is not directly measured; it is extracted from the data through some theoretical model. The standard way to calculate $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{p}`$ is to assume the SU(3) symmetry for the semileptonic decays. In this case it is enough to take any two decays and $`\mathrm{\Gamma }_\mathrm{p}`$ of Eq.(2) as an input. Normally, as in the example above, one uses neutron beta decay and $`\mathrm{\Sigma }^{}`$ decay as an input. However, if the SU(3) symmetry breaking was not important, any pair out of six known semileptonic decays should give roughly the same number for $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{p}`$. This is, however, not the case. As we shall see in the next Section, $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{p}`$ can be any number between 0.02 and 0.30. These numbers do not take into account the experimental errors, therefore, as shown in Figure 1, the uncertainty of $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{p}`$ due to the SU(3) symmetry breaking in the semileptonic decays is even larger. This is the key observation which motivated this work.
It is almost impossible to analyze the symmetry breaking in weak decays without resorting to some specific model . In this paper, following Ref. , we will implement the symmetry breaking for the semileptonic decays using the Chiral Quark-Soliton Model ($`\chi `$QSM for short) (see Ref. for review) which satisfactorily describes the axial-vector properties of the hyperons . Since the symmetry breaking pattern of the $`\chi `$QSM is identical to the one derived in large $`N_\mathrm{c}`$ QCD , our analysis is in fact much more general than the model itself.
However, since $`g_\mathrm{A}^{(0)}(`$B$`)`$ does not correspond to the SU(3) octet axial-vector current, it is an independent quantity in QCD and it cannot be expressed in terms of $`F`$ and $`D`$ without some further assumptions. The $`\chi `$QSM (as most of the hedgehog models ) has a remarkable virtue of connecting the singlet axial-vector constant with $`g_\mathrm{A}^{(3)}`$ and $`g_\mathrm{A}^{(8)}`$, and the semileptonic decay constants in a direct manner. This connection introduces a model dependence into our analysis. However, as we discussed in our previous paper on the proton spin structure and as will be shown in Section V.A, there is no significant numerical difference between the results obtained with and without this model dependent ingredient. Whether this remains true for other baryons cannot be checked because of the lack of the data which could be additionally used if the model formula for $`g_\mathrm{A}^{(0)}(`$B$`)`$ is abandoned.
In Section II.E we give an additional theoretical argument in favor of the model prediction for $`g_\mathrm{A}^{(0)}(`$B$`)`$ .
In the previous paper we have shown how the symmetry breaking influences the determination of $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{p}`$ from the existing data on the weak semileptonic baryon decays. Here our analysis is extended to all other members of the octet using the same ”model-independent” method where the dynamical quantities, which are in principle calculable within the model , are treated as free parameters. By adjusting them to the experimentally known semileptonic decays we allow not only for maximal phenomenological input but also for minimal model dependence. In Refs. magnetic moments of the octet and decuplet have been studied in this way. Model calculations for the vector-axial properties of baryons have been presented in Ref.. There exist also direct model calculations of the spin polarization function itself .
Although the spin content of the hyperons will be most probably not directly measured (with an exception of $`\mathrm{\Lambda }`$ where spin structure function can be related to the measured fragmentation function ), there is a substantial theoretical interest in the spin properties of the hyperons. We find that despite the fact that the symmetry breaking for the semileptonic decays themselves is not strong, other quantities like $`\mathrm{\Delta }s`$ and $`\mathrm{\Delta }\mathrm{\Sigma }`$ are much more affected. We observe splitting of $`\mathrm{\Delta }\mathrm{\Sigma }`$ for different baryons. Unfortunately our analysis suffers from large errors which are mainly due to the experimental errors of the $`\mathrm{\Xi }^{}`$ decays. It is therefore of utmost importance to measure these two decays with higher precision.
The paper is organized as follows: In Section II we recall the SU(3) symmetry results and discuss various ways of determining $`\mathrm{\Delta }\mathrm{\Sigma }`$ and separately $`\mathrm{\Delta }q`$’s. In Section III, following Ref., we recall the main properties of the $`\chi `$QSM with special emphasis on the mass splittings, which we subsequently use in Section IV to parametrize the SU(3) breaking of the semileptonic weak decays. In Section V numerical analysis is carried out and the conclusions are given in Section VI.
## II SU(3) symmetry at work
Let us first briefly recall how the standard analysis is carried out. Three diagonal axial-vector coupling constants define the integrated polarized quark densities for a given baryon B:
$`g_\mathrm{A}^{(3)}(\mathrm{B})`$ $`=`$ $`\mathrm{\Delta }u{}_{\mathrm{B}}{}^{}\mathrm{\Delta }\mathrm{d}_\mathrm{B},`$ (4)
$`\sqrt{3}g_\mathrm{A}^{(8)}(\mathrm{B})`$ $`=`$ $`\mathrm{\Delta }\mathrm{u}_\mathrm{B}+\mathrm{\Delta }\mathrm{d}_\mathrm{B}2\mathrm{\Delta }\mathrm{s}_\mathrm{B},`$ (5)
$`g_\mathrm{A}^{(0)}(\mathrm{B})`$ $`=`$ $`\mathrm{\Delta }\mathrm{u}_\mathrm{B}+\mathrm{\Delta }\mathrm{d}_\mathrm{B}+\mathrm{\Delta }\mathrm{s}_\mathrm{B}.`$ (6)
Note that in our normalization $`g_\mathrm{A}^{(0)}(`$B$`)=`$ $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{B}`$.
Assuming the SU(3) symmetry, one can calculate $`g_\mathrm{A}^{(3,8)}(`$B$`)`$ in terms of the reduced matrix elements $`F`$ and $`D`$:<sup>*</sup><sup>*</sup>*Note that $`g_\mathrm{A}^{(3)}`$ is proportional to $`I_3`$(third component of the isospin which we assume to take the highest value )
$`g_\mathrm{A}^{(3)}(\mathrm{p})=F+D,`$ $`\sqrt{3}g_A^{(8)}(\mathrm{p})=3FD,`$ (7)
$`g_\mathrm{A}^{(3)}(\mathrm{\Lambda })=0,`$ $`\sqrt{3}g_\mathrm{A}^{(8)}(\mathrm{\Lambda })=2D,`$ (8)
$`g_\mathrm{A}^{(3)}(\mathrm{\Sigma }^+)=2F,`$ $`\sqrt{3}g_\mathrm{A}^{(8)}(\mathrm{\Sigma }^+)=2D,`$ (9)
$`g_\mathrm{A}^{(3)}(\mathrm{\Xi }^0)=FD,`$ $`\sqrt{3}g_\mathrm{A}^{(8)}(\mathrm{\Xi }^0)=3FD.`$ (10)
At this stage $`g_\mathrm{A}^{(0)}=\mathrm{\Delta }\mathrm{\Sigma }`$ is an independent quantity and it is identical for all octet states. These equations together with (6) allow one to express $`\mathrm{\Delta }q`$’s in terms of $`D`$, $`F`$ and $`\mathrm{\Delta }\mathrm{\Sigma }`$:
$`\mathrm{\Delta }u_\mathrm{p}`$ $`=`$ $`1/3\left(D+3F+\mathrm{\Delta }\mathrm{\Sigma }\right),`$ (11)
$`\mathrm{\Delta }d_\mathrm{p}`$ $`=`$ $`1/3\left(2D+\mathrm{\Delta }\mathrm{\Sigma }\right),`$ (12)
$`\mathrm{\Delta }s_\mathrm{p}`$ $`=`$ $`1/3\left(D3F+\mathrm{\Delta }\mathrm{\Sigma }\right),`$ (13)
$`\mathrm{\Delta }u_\mathrm{\Lambda }`$ $`=`$ $`1/3\left(D+\mathrm{\Delta }\mathrm{\Sigma }\right),`$ (14)
$`\mathrm{\Delta }s_\mathrm{\Lambda }`$ $`=`$ $`1/3\left(2D+\mathrm{\Delta }\mathrm{\Sigma }\right),`$ (15)
$`\mathrm{\Delta }u_{\mathrm{\Sigma }^0}`$ $`=`$ $`1/3\left(D+\mathrm{\Delta }\mathrm{\Sigma }\right).`$ (16)
The SU(3) symmetry imposes certain relations between $`\mathrm{\Delta }q`$’s of different flavor for different baryons:
$`\mathrm{\Delta }u_\mathrm{p}`$ $`=`$ $`\mathrm{\Delta }u_{\mathrm{\Sigma }^+}=\mathrm{\Delta }s_{\mathrm{\Xi }^0},`$ (17)
$`\mathrm{\Delta }d_\mathrm{p}`$ $`=`$ $`\mathrm{\Delta }s_{\mathrm{\Sigma }^+}=\mathrm{\Delta }u_{\mathrm{\Xi }^0},`$ (18)
$`\mathrm{\Delta }s_\mathrm{p}`$ $`=`$ $`\mathrm{\Delta }d_{\mathrm{\Sigma }^+}=\mathrm{\Delta }d_{\mathrm{\Xi }^0},`$ (19)
so that $`\mathrm{\Delta }q`$’s given in Eq.(16) are the only independent ones in the SU(3) symmetry limit. In addition we have the isospin relations
$`\mathrm{\Delta }u_\mathrm{p}`$ $`=`$ $`\mathrm{\Delta }d_\mathrm{n},\mathrm{\Delta }d_\mathrm{p}=\mathrm{\Delta }u_\mathrm{n},\mathrm{\Delta }s_\mathrm{p}=\mathrm{\Delta }s_\mathrm{n},`$ (20)
$`\mathrm{\Delta }u_{\mathrm{\Sigma }^+}`$ $`=`$ $`\mathrm{\Delta }d_\mathrm{\Sigma }^{},\mathrm{\Delta }d_{\mathrm{\Sigma }^+}=\mathrm{\Delta }u_\mathrm{\Sigma }^{},\mathrm{\Delta }u_{\mathrm{\Sigma }^0}=\mathrm{\Delta }d_{\mathrm{\Sigma }^0}`$ (21)
$`\mathrm{\Delta }u_\mathrm{\Lambda }`$ $`=`$ $`\mathrm{\Delta }d_\mathrm{\Lambda },\mathrm{\Delta }s_{\mathrm{\Sigma }^+}=\mathrm{\Delta }s_\mathrm{\Sigma }^{}=\mathrm{\Delta }s_{\mathrm{\Sigma }^0},`$ (22)
$`\mathrm{\Delta }u_{\mathrm{\Xi }^0}`$ $`=`$ $`\mathrm{\Delta }d_\mathrm{\Xi }^{},\mathrm{\Delta }d_{\mathrm{\Xi }^0}=\mathrm{\Delta }u_\mathrm{\Xi }^{},\mathrm{\Delta }s_{\mathrm{\Xi }^0}=\mathrm{\Delta }s_\mathrm{\Xi }^{}`$ (23)
which remain still valid after the inclusion of the SU(3) symmetry breaking.
In order to find the numerical values of $`\mathrm{\Delta }q`$’s one considers different scenarios which we shortly discuss in the following.
### A Naive quark model
In the naive quark model there exist two relations between the constants $`F`$ and $`D`$ :
$$F/D=2/3,F+D=5/3F=2/3,D=1.$$
(24)
Moreover, one assumes that the total spin is carried by the quarks, i.e.:
$$\mathrm{\Delta }\mathrm{\Sigma }=1.$$
(25)
With these parameters one gets $`\mathrm{\Delta }s_\mathrm{p}=0`$. Values for all $`\mathrm{\Delta }q`$’s and $`\mathrm{\Gamma }_\mathrm{p}`$ are presented in Table I. The prediction for $`\mathrm{\Gamma }_\mathrm{p}`$ is, however, very bad, about two times the experimental value.
### B Extracting $`F`$ and $`D`$ from the semileptonic weak decays
Certainly these naive quark model values (24) are not realistic. One can do better by extracting $`F`$ and $`D`$ from experiment. For example, assuming the exact SU(3) symmetry, one has
$$A_1=\left(g_1/f_1\right)^{(\mathrm{n}\mathrm{p})}=F+D,A_4=\left(g_1/f_1\right)^{(\mathrm{\Sigma }^{}\mathrm{n})}=FD.$$
(26)
For convenience, we denote the ratios of axial-vector to vector decay constants by $`A_i`$ (see Table III). Taking for these decays the experimental values, one obtains
$$F=0.46\mathrm{and}D=0.80,$$
(27)
as displayed in the column $`(A_1,A_4)`$ in Table I.
One could, however, use any two $`A_i`$’s out of six known weak semileptonic decays to extract $`F`$ and $`D.`$ The number of combinations is fourteen (actually fifteen, but two conditions are linearly dependent). Taking these fourteen combinations into account, one gets:
$$F=0.40÷0.55,D=0.70÷0.89.$$
(28)
These are the uncertainties of the central values due to the theoretical error caused by using the exact SU(3) symmetry to describe the weak semileptonic decays. These uncertainties are further increased by the experimental errors of all individual decays.
Looking at Eq.(28), one might get an impression that a typical error associated with the use of the SU(3) symmetry in analyzing the hyperon decays is of the order of 15 % or so. While this is true for the hyperon decays themselves, the values of $`\mathrm{\Delta }q`$ and $`\mathrm{\Delta }\mathrm{\Sigma }`$ for various baryons might be much more affected by the symmetry breaking. Indeed, since
$$\mathrm{\Delta }\mathrm{\Sigma }=\frac{1}{2}\left(\mathrm{\Gamma }_\mathrm{p}3FD\right)$$
(29)
in the SU(3) symmetry limit we get
$$\mathrm{\Delta }\mathrm{\Sigma }=0.02÷0.30$$
(30)
for $`F`$ and $`D`$ corresponding to Eq.(28) and $`\mathrm{\Gamma }_\mathrm{p}`$ as given by Eq.(2). This large uncertainty of the central value of $`\mathrm{\Delta }\mathrm{\Sigma }`$ is entirely due to the SU(3) symmetry breaking in the hyperon decays. In Fig.1 we plot $`\mathrm{\Delta }\mathrm{\Sigma }`$ together with experimental errors for each pair of the semileptonic decays.
Anticipating the results of Section IV let us mention that there exist two linear combinations of $`A_i`$’s which are free of the linear $`m_\mathrm{s}`$ corrections in the $`\chi `$QSM (and large $`N_\mathrm{c}`$ QCD ), namely:
$`F`$ $`=`$ $`{\displaystyle \frac{1}{12}}(4A_14A_23A_3+3A_4+3A_5+5A_6),`$ (31)
$`D`$ $`=`$ $`{\displaystyle \frac{1}{12}}(4A_2+3A_33A_43A_5+3A_6)`$ (32)
which give numerically
$$F=0.50\pm 0.07\mathrm{and}D=0.77\pm 0.04,$$
(33)
as displayed in Table I in the column ”average”. It is important to note that by adopting this way of extracting $`F`$ and $`D`$ in the symmetry limit, no refitting of $`F`$ and $`D`$ is required when $`m_\mathrm{s}`$ corrections are added.
In what follows we shall use these two sets – Eqs.(27,33) – of values for $`F`$ and $`D`$ while discussing the predictions for $`\mathrm{\Delta }q`$’s.
In order to extract all $`\mathrm{\Delta }q`$’s separately, one needs some additional information. Either another experimental input is needed, or a model which predicts $`g_\mathrm{A}^{(0)}`$ ($`\mathrm{B}`$) in terms of $`F`$ and $`D`$.
### C Conjecture of Ellis and Jaffe
In 1974 Ellis and Jaffe made an assumption, based on the naive quark model that
$$\mathrm{\Delta }s_\mathrm{p}=0.$$
(34)
From our SU(3) formula (16), we see that this amounts to
$$\mathrm{\Delta }\mathrm{\Sigma }=3FD$$
(35)
which indeed gives 1 for the naive quark model values (24). For the experimental values of $`F`$ and $`D`$ discussed in the previous section we get $`\mathrm{\Delta }\mathrm{\Sigma }`$ around 0.6 as displayed in Table I. Unfortunately, the value of $`\mathrm{\Gamma }_\mathrm{p}`$ is much larger than the experimental value.
### D Linking hyperon decays with the high energy data
Instead of using the low energy data alone, one can also use the high energy data on the first moment of the polarized structure function of the proton (1) with $`\mathrm{\Gamma }_\mathrm{p}=2.56`$. The results of such fits for two choices of $`F`$ and $`D`$ constants are presented in columns 5 and 6 of Table I. A striking feature of these fits is that the resulting $`\mathrm{\Delta }\mathrm{\Sigma }`$ is very small. This fact is often referred to as a spin crisis.
### E Chiral Quark Soliton Model
As will be shown in the following, the $`\chi `$QSM predicts in the SU(3) symmetry limit :
$$\mathrm{\Delta }\mathrm{\Sigma }=9F5D$$
(36)
for all octet baryons. This formula has a remarkable feature: It interpolates between the naive quark model and the Skyrme model. Indeed, for (24) $`\mathrm{\Delta }\mathrm{\Sigma }=1`$, whereas in the case of the simplest Skyrme model for which $`F/D=5/9`$, $`\mathrm{\Delta }\mathrm{\Sigma }=0`$, as observed for the first time in Ref..
Here $`\mathrm{\Delta }\mathrm{\Sigma }`$ is very sensitive to small variations of $`F`$ and $`D`$, since it is a difference of the two, with relatively large coefficients. Indeed, for the 14 fits mentioned before Eq.(28) the central value for $`\mathrm{\Delta }\mathrm{\Sigma }`$ varies between $`0.25`$ to approximately 1. Thus, despite the fact that the hyperon semileptonic decays are relatively well described by the model in the SU(3) symmetry limit, the singlet axial-vector constant is basically undetermined. This is a clear signal of the importance of the symmetry breaking for this quantity.
In fact, conclusions similar to ours have been obtained in chiral perturbation theory in Ref..
## III Mass splittings in the $`\chi `$QSM
In this Section we shall briefly recall how the model parameters are fixed. Because of the SU(3) symmetry breaking due to the strange quark mass $`m_\mathrm{s}`$ the collective baryon Hamiltonian is no longer SU(3)-symmetric. Indeed :
$$H=H_0+H^{}$$
(37)
where
$$H_0=M_{\mathrm{sol}}+\frac{1}{2I_1}S(S+1)+\frac{1}{2I_2}(C_2(\mathrm{SU}(3)S(S+1)\frac{N_c^2}{12})$$
(38)
and
$$\widehat{H^{}}=m_\mathrm{s}\left(\alpha D_{88}^{(8)}+\beta \widehat{Y}+\frac{\gamma }{\sqrt{3}}\underset{A=1}{\overset{3}{}}D_{8A}^{(8)}\widehat{S}_A\right).$$
(39)
Here $`\widehat{S}_A`$ denotes baryon spin, $`C_2(`$SU(3)$`)`$ the Casimir operator and $`D_{BS}^{()}`$ are the SU(3) Wigner matrices in representation $``$. Constants $`\alpha `$, $`\beta `$ and $`\gamma `$ are given by Ref.:
$$\alpha =\sigma +\frac{K_2}{I_2},\beta =\frac{K_2}{I_2},\gamma =2\left(\frac{K_1}{I_1}\frac{K_2}{I_2}\right).$$
(40)
Here $`K_i`$ and $`I_i`$ are the “moments of inertia” and $`\sigma `$ is related to the nucleon sigma term: $`3\sigma =\mathrm{\Sigma }/\overline{m}`$, $`\overline{m}`$ being the average mass of the up and down quarks.
The collective splitting Hamiltonian (39) mixes the states in various SU(3) representations. The octet states are mixed with the higher representations such as antidecuplet $`\overline{\mathrm{𝟏𝟎}}`$ and eikosiheptaplet $`\mathrm{𝟐𝟕}`$. In the linear order in $`m_\mathrm{s}`$ the wave function of a state $`B=(Y,I,I_3)`$ of spin $`S_3`$ is given as:
$$\psi _{B,S_3}=()^{\frac{1}{2}S_3}\left(\sqrt{8}D_{BS}^{(8)}+c_B^{(\overline{10})}\sqrt{10}D_{BS}^{(\overline{10})}+c_B^{(27)}\sqrt{27}D_{BS}^{(27)}\right),$$
(41)
where $`S=(1,\frac{1}{2},S_3)`$. Mixing parameters $`c_B^{()}`$ can be found for example in Ref. . They are given as products of $`m_\mathrm{s}`$ (which we assume to be 180 MeV) times a known numerical constant $`N_B^{()}`$ depending on the baryonic state $`B`$ and a dynamical parameter $`c_{}`$. Since $`c_{}`$ depends on the model parameter $`I_2`$, which is responsible for the splitting between the octet and higher exotic multiplets and is not constrained from the data we will take them as free parameters in our fits.
## IV Semileptonic weak decays in the chiral quark-soliton model
The transition matrix elements of the hadronic axial-vector current $`B_2|A_\mu ^X|B_1`$ can be expressed in terms of three independent form factors:
$$B_2|A_\mu ^X|B_1=\overline{u}_{B_2}(p_2)\left[\left\{g_1(q^2)\gamma _\mu \frac{ig_2(q^2)}{M_1}\sigma _{\mu \nu }q^\nu +\frac{g_3(q^2)}{M_1}q_\mu \right\}\gamma _5\right]u_{B_1}(p_1),$$
(42)
where the axial-vector current is defined as
$$A_\mu ^X=\overline{\psi }(x)\gamma _\mu \gamma _5\lambda _X\psi (x)$$
(43)
with $`X=\frac{1}{2}(1\pm i2)`$ for strangeness conserving $`\mathrm{\Delta }S=0`$ currents and $`X=\frac{1}{2}(4\pm i5)`$ for $`|\mathrm{\Delta }S|=1`$.
The $`q^2=Q^2`$ stands for the square of the momentum transfer $`q=p_2p_1`$. The form factors $`g_i`$ are real quantities depending only on the square of the momentum transfer in the case of $`CP`$-invariant processes. We can safely neglect $`g_3`$ for the reason that on account of $`q_\mu `$ its contribution to the decay rate is proportional to the ratio $`\frac{m_l^2}{M_1^2}1`$, where $`m_l`$ represents the mass of the lepton ($`e`$ or $`\mu `$) in the final state and $`M_1`$ that of the baryon in the initial state. Similarly we shall neglect $`g_2`$. In principle this form factor is proportional to $`m_\mathrm{s}`$ and therefore should be included in the consistent analysis of the weak decays data. Unfortunately, such an analysis is still missing and all experimental results on $`g_1`$ assume $`g_20`$.
Another possible small $`m_\mathrm{s}`$ corrections come from the evolution of $`g_1`$ with $`Q^2`$, due to the non-conservation of the axial-vector currents caused by the SU(3) symmetry breaking. These corrections are also neglected in our approach.
It is already well known how to treat hadronic matrix elements such as $`B_2|A_\mu ^X|B_1`$ in the $`\chi `$QSM (see for example and references therein). Taking into account the $`1/N_c`$ rotational and $`m_\mathrm{s}`$ corrections, we can write the resulting axial-vector constants $`g_1^{B_1B_2}(0)`$ in the following form:
$`g_1^{(B_1B_2)}`$ $`=`$ $`a_1B_2|D_{X3}^{(8)}|B_1+a_2d_{pq3}B_2|D_{Xp}^{(8)}\widehat{S}_q|B_1+{\displaystyle \frac{a_3}{\sqrt{3}}}B_2|D_{X8}^{(8)}\widehat{S}_3|B_1`$ (44)
$`+`$ $`m_s[{\displaystyle \frac{a_4}{\sqrt{3}}}d_{pq3}B_2|D_{Xp}^{(8)}D_{8q}^{(8)}|B_1+a_5B_2|(D_{X3}^{(8)}D_{88}^{(8)}+D_{X8}^{(8)}D_{83}^{(8)})|B_1`$ (45)
$`+`$ $`a_6B_2|(D_{X3}^{(8)}D_{88}^{(8)}D_{X8}^{(8)}D_{83}^{(8)})|B_1],`$ (46)
where $`a_i`$ denote parameters depending on the specific dynamics of the chiral soliton model. Their explicit form in the $`\chi `$QSM can be found in Ref. .
Analogously to Eq.(46) one defines the diagonal axial-vector couplings. In that case $`X`$ can take two values: $`X=3`$ and $`X=8`$. For $`X=0`$ (singlet axial-vector current) we have the following expression :
$$g_B^{(0)}\widehat{S}_3=a_3\widehat{S}_3+\sqrt{3}m_\mathrm{s}(a_5a_6)B|D_{83}^{(8)}|B.$$
(47)
A remark concerning constants $`a_i`$ is here in order. Coefficient $`a_1`$ contains the terms which are leading and subleading in the $`1/N_\mathrm{c}`$ expansion. The presence of the subleading terms enhances the numerical value of $`a_1`$ calculated in the $`\chi `$QSM for the self-consistent profile and makes the model predictions remarkably close to the experimental data . This feature, although very important for the model phenomenology, does not concern us here, since our procedure is based on fitting all coefficients $`a_i`$ from the data. Constants $`a_2`$ and $`a_3`$ are both subleading in $`1/N_\mathrm{c}`$ and come from the anomalous part of the effective Euclidean action. In the Skyrme model they are related to the Wess-Zumino term. However in the simplest version of the Skyrme model (which is based on the pseudo-scalar mesons only) $`a_3=0`$ identically . In the case of the $`\chi `$QSM $`a_30`$ and it provides a link between the SU(3) octet of axial-vector currents and the singlet current of Eq.(47). It was shown in Ref. that in the limit of the artificially large soliton, which corresponds to the “Skyrme limit” of the present model, $`a_3/a_10`$ in agreement with . On the contrary, for the small solitons $`g_\mathrm{p}^{(0)}1`$ reproducing the result of the non-relativistic quark model.
Instead of calculating 7 dynamical parameters $`a_i`$ and $`I_2`$ (or $`c_{\overline{10}}`$ and $`c_{27}`$) within the $`\chi `$QSM (which was done in Ref., we shall fit them from the weak semileptonic decay data. It is convenient to introduce the following set of new parameters:
$`r={\displaystyle \frac{1}{30}}\left(a_1{\displaystyle \frac{1}{2}}a_2\right),s={\displaystyle \frac{1}{60}}a_3,x={\displaystyle \frac{1}{540}}m_\mathrm{s}a_4,y={\displaystyle \frac{1}{90}}m_\mathrm{s}a_5,z={\displaystyle \frac{1}{30}}m_\mathrm{s}a_6,`$
$$p=\frac{1}{6}m_\mathrm{s}c_{\overline{10}}\left(a_1+a_2+\frac{1}{2}a_3\right),q=\frac{1}{90}m_\mathrm{s}c_{27}\left(a_1+2a_2\frac{3}{2}a_3\right).$$
(48)
Employing this new set of parameters, we immediately express all possible semileptonic decay constants between the octet baryons:
$`\left(g_1/f_1\right)^{(\mathrm{n}\mathrm{p})}`$ $`=`$ $`14r+2s44x20y4z4p+8q,`$ (49)
$`\left(g_1/f_1\right)^{(\mathrm{\Sigma }^+\mathrm{\Lambda })}`$ $`=`$ $`9r3s42x6y3p+15q,`$ (50)
$`\left(g_1/f_1\right)^{(\mathrm{\Lambda }\mathrm{p})}`$ $`=`$ $`8r+4s+24x2z+2p6q,`$ (51)
$`\left(g_1/f_1\right)^{(\mathrm{\Sigma }^{}\mathrm{n})}`$ $`=`$ $`4r+8s4x4y+2z+4q,`$ (52)
$`\left(g_1/f_1\right)^{(\mathrm{\Xi }^{}\mathrm{\Lambda })}`$ $`=`$ $`2r+6s6x+6y2z+6q,`$ (53)
$`\left(g_1/f_1\right)^{(\mathrm{\Xi }^{}\mathrm{\Sigma }^0)}`$ $`=`$ $`14r+2s+22x+10y+2z+2p4q,`$ (54)
$`\left(g_1/f_1\right)^{(\mathrm{\Sigma }^{}\mathrm{\Lambda })}`$ $`=`$ $`9r3s42x6y3p+15q,`$ (55)
$`\left(g_1/f_1\right)^{(\mathrm{\Sigma }^{}\mathrm{\Sigma }^0)}`$ $`=`$ $`5r+5s18x6y+2z2p,`$ (56)
$`\left(g_1/f_1\right)^{(\mathrm{\Xi }^{}\mathrm{\Xi }^0)}`$ $`=`$ $`4r+8s+8x+8y4z8q,`$ (57)
$`\left(g_1/f_1\right)^{(\mathrm{\Xi }^0\mathrm{\Sigma }^+)}`$ $`=`$ $`14r+2s+22x+10y+2z+2p4q.`$ (58)
The U(3) axial-vector constants $`g_A^{(0,3,8)}`$ can be also expressed in terms of the new set of parameters Eq.(48). For the triplet onesTriplet $`g^{(3)}`$’s are proportional to $`I_3`$, formulae in Eq.(62) correspond to the highest isospin state we have:
$`g_\text{A}^{(3)}(\text{p})`$ $`=`$ $`14r+2s44x20y4z4p+8q,`$ (59)
$`g_\text{A}^{(3)}(\mathrm{\Lambda })`$ $`=`$ $`0,`$ (60)
$`g_\text{A}^{(3)}(\mathrm{\Sigma }^+)`$ $`=`$ $`10r+10s36x12y+4z4p,`$ (61)
$`g_\text{A}^{(3)}(\mathrm{\Xi }^0)`$ $`=`$ $`4r+8s+8x+8y4z8q,`$ (62)
and for the octet ones, we get:
$`g_\text{A}^{(8)}(\text{p})`$ $`=`$ $`\sqrt{3}(2r+6s+12x+4p+24q),`$ (63)
$`g_\text{A}^{(8)}(\mathrm{\Lambda })`$ $`=`$ $`\sqrt{3}(6r+2s36x+36q),`$ (64)
$`g_\text{A}^{(8)}(\mathrm{\Sigma }^+)`$ $`=`$ $`\sqrt{3}(6r2s+20x+8y+4p+16q),`$ (65)
$`g_\text{A}^{(8)}(\mathrm{\Xi }^0)`$ $`=`$ $`\sqrt{3}(8r4s24x12y+24q).`$ (66)
As already explained in the Introduction the model provides a link between the octet currents and the singlet axial current. For the singlet axial-vector constants, we have:
$`g_\text{A}^{(0)}(\text{p})`$ $`=`$ $`60s18y+6z,`$ (67)
$`g_\text{A}^{(0)}(\mathrm{\Lambda })`$ $`=`$ $`60s+54y18z,`$ (68)
$`g_\text{A}^{(0)}(\mathrm{\Sigma })`$ $`=`$ $`60s54y+18z,`$ (69)
$`g_\text{A}^{(0)}(\mathrm{\Xi })`$ $`=`$ $`60s+72y24z,`$ (70)
Let us note that by redefinition of $`q`$ and $`x`$ we can get rid of the variable $`p`$:
$$x^{}=x\frac{1}{9}p,q^{}=q\frac{1}{9}p.$$
(71)
In the chiral limit parameters $`x`$, $`y`$, $`z`$, $`p`$ and $`q`$ vanish and we recover the SU(3) symmetric relations from Section II with
$$D=3s9r,F=5s5r,$$
(72)
from which Eq.(36) follows.
## V The SU(3) symmetry breaking
We fix the newly-defined set of parameters from the experimental data of semileptonic decays. Their numerical values are given in Table III. We do not quote the experimental errors on these parameters, since they are highly correlated and cannot be used directly to calculate the errors of the physical quantities of interest. Instead, we expressed all observables directly in terms of the $`A_i`$’s. This is, however, not enough since, as in the chiral limit, the extra input is needed.
At this point a necessity of a complete description of the symmetry breaking is clearly seen. The strange quark mass causes all SU(3) symmetry relations (19) to break. So in principle one needs one extra experimental input for each isospin multiplet. Let us first discuss the case of the nucleon first.
### A Spin content of the nucleon
We shall repeat here the analysis of Section II, however, with the symmetry breaking taken into account. Again four different choices for an additional input will be considered: 1) $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{p}=1`$, 2) $`\mathrm{\Delta }s_\mathrm{p}=0`$, 3) $`\mathrm{\Gamma }_\mathrm{p}=2.56`$ and 4) the $`\chi `$QSM formulae (70) for $`g_\mathrm{A}^{(0)}`$. The results are summarized in Table II. It can be immediately seen that the first two possibilities are in contradiction with experimental data on $`\mathrm{\Gamma }_\mathrm{p}`$ and $`\mathrm{\Gamma }_\mathrm{n}`$. On the other hand, if we use the experimental value of $`\mathrm{\Gamma }_\mathrm{p}`$ as an additional input (but no model formula (70) for $`g_\mathrm{A}^{(0)}`$), or alternatively the $`\chi `$QSM prediction for $`g_\mathrm{A}^{(0)}`$, the results are almost indistinguishable. This gives a numerical support for the correctness of the $`\chi `$QSM formula for the axial-vector singlet current with the SU(3) symmetry breaking.
Of course the results of Table II have to be taken with a bit of care because of large experimental errors which are not displayed. As we have argued in Ref., one could still accommodate $`\mathrm{\Delta }s_\mathrm{p}=0`$ due to the large errors of $`\mathrm{\Xi }`$ decays. We shall come back to this point in the following.
### B Numerical results
It the present Section we shall present the numerical results of our analysis based on the Chiral Quark Soliton Model with the SU(3) symmetry breaking. Our strategy is very simple: using model parametrization (58) we expressed $`\mathrm{\Delta }q`$’s and $`\mathrm{\Delta }\mathrm{\Sigma }`$’s in terms of the six known weak semileptonic decays. Errors are added in quadrature. The numerical results are summarized in Table IV and in Figures 2 – 9. To guide an eye it is convenient to restore the linear $`m_\mathrm{s}`$ dependence for the quark densities in the following way:
$$\mathrm{\Delta }q=\mathrm{\Delta }q^{(0)}+\frac{m_\mathrm{s}}{180\mathrm{MeV}}\left(\mathrm{\Delta }q\mathrm{\Delta }q^{(0)}\right),$$
(73)
and similarly for $`\mathrm{\Delta }\mathrm{\Sigma }`$. This is possible because our chiral parameters $`r`$ and $`s`$ do not need to be refitted as the symmetry breaking corrections are included. In order to display the errors which come from the experimental errors of the weak decays, at both ends of each figure we also plot the theoretical predictions as black dots together with the error bars.
Let us first comment on the results on $`\mathrm{\Gamma }_\mathrm{p}`$ and $`\mathrm{\Gamma }_\mathrm{n}`$. We see from Table III that the experimental values are quite well reproduced by the model, provided the $`m_\mathrm{s}`$ corrections are included. In the symmetry limit their values are way off from the experimental data.
Next, let us observe that the singlet axial-vector current couplings $`g_\mathrm{A}^{(0)}`$ split when the symmetry breaking is switched on. This is due to the term proportional to $`D_{83}^{(8)}`$ in Eq.(47). This splitting is depicted in Fig.2. We see that $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{p}`$ shows the weakest $`m_\mathrm{s}`$ dependence, whereas $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{\Lambda }`$ and $`\mathrm{\Delta }\mathrm{\Sigma }_\mathrm{\Xi }`$ depend quite strongly on $`m_\mathrm{s}`$. Large error bars for these quantities are due almost entirely to the large errors of $`\mathrm{\Xi }`$ decays $`A_5`$ and $`A_6`$. It is however evident from Fig. 2 that $`\mathrm{\Lambda }`$ and $`\mathrm{\Xi }`$ are much closer to the nonrelativistic limit than p and $`\mathrm{\Sigma }`$.
In Figs.3 – 6 we plot $`\mathrm{\Delta }q`$ for the nucleon, $`\mathrm{\Lambda }`$, $`\mathrm{\Sigma }`$ and $`\mathrm{\Xi }`$ respectively. We see that in all 4 cases $`\mathrm{\Delta }s`$ rises relatively strongly with $`m_\mathrm{s}`$. It is therefore not justified to extract the strange quark polarization assuming the exact SU(3) symmetry. Unfortunately, $`\mathrm{\Delta }s`$’s have also the largest error coming, as in the case of $`\mathrm{\Delta }\mathrm{\Sigma }`$, almost entirely from the errors of $`\mathrm{\Xi }`$ decays.
In Figs.7–9 we examine the breaking of the SU(3) relations given by Eqs.(19). Interestingly we find that there is an approximate equality between $`\mathrm{\Delta }u_\mathrm{p}`$ and $`\mathrm{\Delta }u_{\mathrm{\Sigma }^+}`$ for all values of $`m_\mathrm{s}`$.
## VI Summary and conclusions
In the analysis of the polarized structure function $`g_1`$ of the proton and neutron one has to take an additional input from the low energy hyperon decays. Customarily the SU(3) symmetry for these decays is assumed. However, if one takes all possible combinations of the low energy decays the resulting $`\mathrm{\Delta }\mathrm{\Sigma }`$ can take any value between 0.02 and 0.30. As depicted in Fig.1 this range is further increased if the errors coming from the experimental error bars of the semileptonic decays are properly included. This observation implies that the SU(3) symmetry breaking plays an essential role in extracting $`\mathrm{\Delta }\mathrm{\Sigma }`$ from the experimental data. It was therefore the aim of this paper to study the influence of the symmetry breaking on the determination of $`\mathrm{\Delta }\mathrm{\Sigma }`$ and $`\mathrm{\Delta }s`$ for the octet baryons in a consistent way.
For this purpose we have performed the ”model-independent” analysis based on the algebraic structure of the Chiral Quark Soliton Model. In this approach, one makes merely use of the algebraical structure of the model, treating the dynamical quantities which are in principle calculable in the model as free parameters. Model predictions of the axial-vector properties of the octet baryons have been already calculated elswhere . There are two model ingredients which are of importance. The first one is the model formula for the octet axial-vector currents which have been derived in the linear order in $`m_\mathrm{s}`$ and $`1/N_\mathrm{c}`$. Our formulae here have the same algebraical structure as in the large $`N_\mathrm{c}`$ QCD , and therefore they are more general than the model itself. Secondly, unlike in QCD, the model provides a link between the octet axial-vector currents and the singlet axial-vector current. This connection is a truly model-dependent ingredient, however, we have given the arguments in favor of Eq.(47), based on the fact that apart from the general success of the $`\chi `$QSM in reproducing form factors and parton distributions, in the limit of the small soliton it properly reduces to the Nonrelativistic Quark Model prediction, and in the limit of the large soliton it reproduces the Skyrme Model prediction for $`\mathrm{\Delta }\mathrm{\Sigma }`$. Similarly, in Ref. the argument has been given that Eq.(36) naturally emerges in the limit of the large $`m_\mathrm{s}`$, where the SU(3) flavor symmetry reduces to the SU(2) one. The numerical analysis of Section V.A provides a further support for the model formula for $`\mathrm{\Delta }\mathrm{\Sigma }`$.
We have presented two parametrizations of all available semileptonic decays. The first one is obtained assuming the SU(3) symmetry, however the two reduced matrix elements $`F`$ and $`D`$ were extracted from the combinations of the semileptonic decays which are free of the $`m_\mathrm{s}`$ corrections (32), rather than from the neutron and $`\mathrm{\Sigma }^{}`$ decays alone. The second one is obtained by fitting all 6 measured semileptonic decays in terms of 6 free parameters defined in Eqs.(48,58). The difference between the two fits, as seen from Table III, is rather small, except perhaps for the $`\mathrm{\Sigma }^{}\mathrm{n}`$ decay. Despite the fact that the symmetry breaking for the semileptonic decays themselves is not strong, other quantities like $`\mathrm{\Delta }s`$ and $`\mathrm{\Delta }\mathrm{\Sigma }`$ are much more affected by taking into account the effects of the non-zero strange quark mass. This is clearly shown in Figs.2–9.
Whether this sensitivity is a sign of the breakdown of the perturbative approach to the strangenes, as it was recently suggested in Ref., is hard to say, since our anaysis suffers from large errors which are mainly due to the experimental errors of the $`\mathrm{\Xi }^{}`$ decays. It is therefore of utmost importance to measure these two decays with the precision comparable to the other four decays. One should bare in mind that this is one of a few cases, where the low energy data have an important impact on our understanding of the high energy scattering. Given the theoretical implications of these experiments as far as the role of the axial anomaly and the gloun polarization is concerned , one should make it clear how important the new measurements of the $`\mathrm{\Xi }^{}`$ decays would be. This is perhaps the most important message of our analysis.
## Acknowledgments
The work of H.-Ch.K has been supported by the Korean Physical Society. The work of M.P. has been supported by Polish KBN Grant PB 2 P03B 019 17. The work of K.G. has been supported by the BMBF, the DFG, and the COSY–Project(Jülich).
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# Triple Linking of Surfaces in 4-Space11footnote 1MRCN:57Q45
## 1 Introduction
A surface-link is a closed surface $`F`$ embedded in $`𝐑^4`$ locally flatly. In this paper, we always assume that $`F`$ is oriented, that is, each component of $`F`$ is orientable and given a fixed orientation. For a 3-component surface-link $`F=K_1K_2K_3`$, a linking number was defined in using its projection in $`𝐑^3`$ in a way that is analogous to the linking number in classical knot theory. In that paper, it was introduced as an example of non-triviality of the state-sum invariants of surface-links. In fact, the state-sum invariants in the classical link and surface-link case generalize linking number and Fox’s coloring number.
In the current paper, we give several alternative definitions of the triple linking number and some properties. The reader will find that this invariant is a quite natural generalization of the notion of classical linking number in contrast to a statement in Rolfsen page 136: “There is, however no analogous notion of linking number to help us with codimension two link theory, for example, in higher dimensions”. We note, however, that Rolfsen himself with Massey and with Fenn generalized classical linking numbers to higher dimensions using degrees of maps.
Let $`F=K_1K_2K_3`$ be a $`3`$-component surface-link in $`𝐑^4`$. It is known (see for example) that a projection of $`F`$ into $`𝐑^3`$ can be assumed to have transverse double curves and isolated branch/triple points. At a triple point, three sheets intersect that have distinct relative heights with respect to the projection direction, and we call them top, middle, and bottom sheets, accordingly. If the orientation normals to the top, middle, bottom sheets at a triple point $`\tau `$ matches with this order the fixed orientation of $`𝐑^3`$, then the sign of $`\tau `$ is positive and $`\epsilon (\tau )=1`$. Otherwise the sign is negative and $`\epsilon (\tau )=1`$. (See .) It is also known that any closed oriented embedded surface $`F`$ in $`𝐑^4`$ bounds an oriented compact $`3`$-manifold $`M`$ embedded in $`𝐑^4`$, called a Seifert hypersurface of $`F`$, such that $`M=F`$.
We give six methods for defining an integer (triple linking number) as follows.
* Consider a surface diagram of $`K_1K_2K_3`$ in $`𝐑^3`$. A triple point is of type $`(i,j,k)`$ if the top sheet comes from $`K_i`$, the middle comes from $`K_j`$, and the bottom comes from $`K_k`$. The sum of the signs of all the triple points of type $`(1,2,3)`$ is denoted by $`\text{Tlk}_1(K_1,K_2,K_3)`$. This is the definition given in .
* Let $`M_i`$ be a Seifert hypersurface for $`K_i`$ $`(i=1,3)`$. Assume that $`M_iK_2`$ is a $`1`$-manifold in $`K_2`$ and that $`M_1K_2`$ and $`M_3K_2`$ intersect transversely. Count the intersections between them algebraically and denote the sum by $`\text{Tlk}_2(K_1,K_2,K_3)`$.
* Consider a Seifert hypersurface $`M_1`$ for $`K_1`$. Assume that $`M_1K_2`$ is a $`1`$-manifold, which is disjoint from $`K_3`$. The linking number $`\text{Link}(M_1K_2,K_3)`$ is denoted by $`\text{Tlk}_3(K_1,K_2,K_3)`$.
* Let $`M_i`$ be a Seifert hypersurface for $`K_i`$ $`(i=1,3)`$ such that $`M_1M_3`$ is a $`2`$-manifold which intersects $`K_2`$ transversely. Count the intersections between them algebraically and denote the sum by $`\text{Tlk}_4(K_1,K_2,K_3)`$.
* Let $`M_i`$ be a Seifert hypersurface for $`K_i`$ $`(i=1,2,3)`$ and let $`N_2`$ be a regular neighbourhood of $`K_2`$ in $`𝐑^4`$. We may assume that $`M_iN_2`$ is a $`2`$-manifold in $`N_2`$ and that $`M_1N_2`$, $`M_2N_2`$ and $`M_3N_2`$ intersect transversely in a finite number of points. Count the intersections algebraically and denote the sum by $`\text{Tlk}_5(K_1,K_2,K_3)`$.
* Let $`f:F_1F_2F_3𝐑^4`$ denote an embedding of the disjoint union of oriented surfaces $`F_i`$ representing $`F=K_1K_2K_3`$. Define a map $`L:F_1\times F_2\times F_3S^3\times S^3`$ by
$$L(x_1,x_2,x_3)=(\frac{f(x_1)f(x_2)}{f(x_1)f(x_2)},\frac{f(x_2)f(x_3)}{f(x_2)f(x_3)})$$
for $`x_1F_1`$, $`x_2F_2`$ and $`x_3F_3`$, and denote the degree of $`L`$ by $`\text{Tlk}_6(K_1,K_2,K_3)`$.
###### Theorem 1.1
$`\text{Tlk}_i(K_1,K_2,K_3)=\pm \text{Tlk}_j(K_1,K_2,K_3)`$ for any $`i,j=1,\mathrm{},6.`$
Remark. In general, the triple linking number $`\text{Tlk}(K_i,K_j,K_k)`$ for $`ijk`$ is defined to be the sum of the signs of all the triple points of type $`(i,j,k)`$ on a surface diagram of $`F`$;
$$\text{Tlk}(K_i,K_j,K_k)=\underset{\tau :\text{type}(i,j,k)}{}\epsilon (\tau ).$$
It is proved in that this number is an invariant of the surface-link $`F`$ (independent of a diagram in $`𝐑^3`$) by use of Roseman moves (Reidemeister moves for surface-link diagrams) , and that this invariant vanishes in the case that $`i=k`$; that is, $`\text{Tlk}(K_i,K_j,K_i)=0`$ for $`ij`$. Hence throughout this paper, we always assume that $`i,j,k`$ are all distinct whenever we refer to $`\text{Tlk}(K_i,K_j,K_k)=\text{Tlk}_1(K_i,K_j,K_k)`$.
We prove the following properties of triple linking by using the above interpretations.
###### Theorem 1.2 ()
* $`\text{Tlk}(K_1,K_2,K_3)=\text{Tlk}(K_3,K_2,K_1)`$.
* $`\text{Tlk}(K_1,K_2,K_3)+\text{Tlk}(K_2,K_3,K_1)+\text{Tlk}(K_3,K_1,K_2)=0`$.
###### Theorem 1.3
* If $`K_2`$ is homeomorphic to a $`2`$-sphere, then $`\text{Tlk}(K_1,K_2,K_3)=0`$.
* If both of $`K_1`$ and $`K_3`$ are homeomorphic to a $`2`$-sphere, then $`\text{Tlk}(K_1,K_2,K_3)=0`$.
In the asymmetric linking number $`\text{Alk}(K,K^{})`$ for a two component oriented surface-link $`F=KK^{}`$ was defined to be the non-negative generator of the image of $`H_1(K)H_1(S^4\backslash K^{})𝐙`$.
###### Theorem 1.4
If $`\text{Alk}(K_2,K_3)=0`$, then $`\text{Tlk}(K_1,K_2,K_3)=\text{Tlk}(K_3,K_2,K_1)=0`$.
Two $`n`$-component surface-links $`F=K_1\mathrm{}K_n`$ and $`F^{}=K_1^{}\mathrm{}K_n^{}`$ are link homologous if there is a compact oriented $`3`$-manifold $`W`$ properly embedded in $`𝐑^4\times [0,1]`$ such that $`W`$ has $`n`$ components $`W_1,\mathrm{},W_n`$ with $`W_i=K_i\times \{0\}(K_i^{})\times \{1\}`$. This relation is sometimes called link-cobordism, but that term also denotes the concordance relation. Since link homotopy implies link homology, the following theorem implies that triple linking invariants are link homotopy invariants (this fact is also seen from the sixth definition of Tlk). For related topics, refer to .
###### Theorem 1.5
Triple linking invariants are link homology invariants: If $`F=K_1K_2K_3`$ and $`F^{}=K_1^{}K_2^{}K_3^{}`$ are link homologous, then $`\text{Tlk}(K_i,K_j,K_k)=\text{Tlk}(K_i^{},K_j^{},K_k^{})`$.
See Remark 6.2 for further information about link homology.
By Theorem 1.2, for any 3-component surface-link $`F=K_1K_2K_3`$, there exists a pair of integers $`a`$ and $`b`$ such that
$$()\{\begin{array}{c}\text{Tlk}(K_1,K_2,K_3)=\text{Tlk}(K_3,K_2,K_1)=(a+b),\hfill \\ \text{Tlk}(K_2,K_3,K_1)=\text{Tlk}(K_1,K_3,K_2)=b,\hfill \\ \text{Tlk}(K_3,K_1,K_2)=\text{Tlk}(K_2,K_1,K_3)=a.\hfill \end{array}$$
In , it is shown that for any pair of integers $`a`$ and $`b`$, there exists a surface-link $`F`$ whose triple linking numbers satisfy the above equations. However, that paper does not treat any problem about genera of the components of $`F`$. By Theorem 1.3, we see that
(1) if $`a0`$ and $`b=0`$, then $`g(K_i)1`$ $`(i=1,2)`$, and
(2) if $`a0`$, $`b0`$ and $`a+b0`$, then $`g(K_i)1`$ $`(i=1,2,3)`$,
where $`g(K_i)`$ denotes the genus of $`K_i`$.
###### Proposition 1.6
(i) For any integer $`a0`$, there exists a surface-link $`F=K_1K_2K_3`$ whose triple linking numbers satisfy the above equations $`()`$ with $`b=0`$ and $`g(K_i)=1`$ $`(i=1,2)`$ and $`g(K_3)=0`$.
(ii) For any pair of integers $`a`$ and $`b`$ with $`a0`$, $`b0`$ and $`a+b0`$, there exists a surface-link $`F=K_1K_2K_3`$ whose triple linking numbers satisfy the above equations $`()`$ and $`g(K_i)=1`$ $`(i=1,2,3)`$.
This paper is organized as follows: in Section 2, we interprete $`\text{Tlk}_1`$ in terms of the decker curves of a surface diagram. In Section 3 we give precise definitions of triple linking numbers $`\text{Tlk}_i`$ for $`i=2,\mathrm{},5`$ (in terms of homology) and prove Theorem 1.1. Section 4 is devoted to proving Theorems 1.2–1.5. Proposition 1.6 is proved in Section 5.
Throughout this paper, all the homology and cohomology groups have the Z-coefficient.
## 2 Decker Curves and Triple Linking
Let $`F`$ be a surface-link and $`F^{}`$ a surface diagram of $`F`$ with respect to a projection $`p:𝐑^4𝐑^3`$. Let $`\mathrm{\Gamma }(F^{})`$ denote the double point set of $`F^{}`$;
$$\{p(x)|xF,p(x)=p(y)\text{for some}yF,xy\},$$
which consists of immersed curves, called double curves. A double curve $`C^{}`$ is an immersed circle or an immersed arc in $`𝐑^3`$. If $`C^{}`$ is an immersed circle, then $`(p|_F)^1(C^{})=CC^{}`$ for some pair of immersed circles $`C`$ and $`C^{}`$ in $`F`$. If $`C^{}`$ is an immersed arc, then its endpoints are branch points of $`F^{}`$ and $`(p|_F)^1(C^{})=CC^{}`$ for some pair of immersed arcs $`C`$ and $`C^{}`$ in $`F`$ with $`C=C^{}`$. The curves $`C`$ and $`C^{}`$ are called decker curves over $`C^{}`$: one of them is in higher position than the other with respect to the projection direction, which is called an upper decker curve and the other is called an lower decker curve. We notice that the preimage of a triple point consists of three points of $`F`$ which are intersections of decker curves. See for details.
Double curves and decker curves are oriented as follows: Let $`x`$ be a point of $`F`$ whose image $`x^{}=p(x)`$ is not a branch point. There is a regular neighborhood $`N`$ of $`x`$ in $`F`$ such that $`p|_N`$ is an embedding. An orientation normal $`\stackrel{}{n}`$ to $`N^{}=p(N)`$ in $`𝐑^3`$ at $`x^{}`$ is specified in such a way that $`(\stackrel{}{v}_1,\stackrel{}{v}_2,\stackrel{}{n})`$ matches the orientation of $`𝐑^3`$, where the pair of tangents $`(\stackrel{}{v}_1,\stackrel{}{v}_2)`$ defines the orientation of $`N^{}`$ that is induced from the orientation of $`NF`$. If $`y`$ is a double point on a double curve $`C^{}`$, then $`C^{}`$ is locally an intersection of $`N_1^{}`$ and $`N_2^{}`$, where $`N_1^{}`$ is upper and $`N_2^{}`$ is lower. We assign a tangent vector $`\stackrel{}{v}`$ of $`C^{}`$ at $`y`$ such that $`(\stackrel{}{n}_1,\stackrel{}{n}_2,\stackrel{}{v})`$ matches the orientation of $`𝐑^3`$. This defines an orientation of $`C^{}`$, cf. . We give an orientation to the lower decker curve over $`C^{}`$ such that it inherits the orientation from $`C^{}`$, and give the opposite orientation to the upper decker curve. Note that, if $`C^{}`$ is an arc, then the orientations of $`C`$ and $`C^{}`$ are compatible (i.e., the union $`CC^{}`$ forms an oriented immersed circle in $`F`$).
Let $`F=K_1K_2K_3`$ be a 3-component surface-link. A double curve $`C^{}`$ is of type $`(i,j)`$ if the upper decker curve lies in $`K_i`$ and the lower decker curve lies in $`K_j`$. A decker curve over $`C^{}`$ is of type $`(i,j)`$ if $`C^{}`$ is so.
At a triple point $`\tau `$, if the orientation normals to the top, middle, and bottom sheets at $`\tau `$ matches with this order the fixed orientation of $`𝐑^3`$, then the sign of $`\tau `$ is positive and $`\epsilon (\tau )=+1`$; otherwise the sign is negative and $`\epsilon (\tau )=1`$.
We interprete the triple linking $`\text{Tlk}_1`$ in terms of double decker curves as follows. Let $`D_{12}^{\mathrm{}}`$ (resp. $`D_{23}^u`$) denote the union of lower decker curves of type $`(1,2)`$ (resp. upper decker curves of type $`(2,3)`$). Note that both $`D_{12}^{\mathrm{}}`$ and $`D_{23}^u`$ are contained in $`K_2`$.
###### Lemma 2.1
$`\text{Tlk}_1(K_1,K_2,K_3)=\text{Int}_{K_2}(D_{12}^{\mathrm{}},D_{23}^u),`$ where $`\text{Int}_{K_2}(D_{12}^{\mathrm{}},D_{23}^u)`$ is the intersection number in $`K_2`$.
Proof. Let $`\tau `$ be a triple point of type $`(1,2,3)`$. The preimage of $`\tau `$ consists of three points of $`F`$. Exactly one of them is on $`K_2`$ and that is a double point of $`D_{12}^{\mathrm{}}`$ and $`D_{23}^u`$. Conversely the image of a double point of $`D_{12}^{\mathrm{}}`$ and $`D_{23}^u`$ is a triple point of $`F^{}`$ of type $`(1,2,3)`$. Hence there is a one-to-one correspondence between the set of triple points of type $`(1,2,3)`$ and double points of $`D_{12}^{\mathrm{}}`$ and $`D_{23}^u`$. If the sign of $`\tau `$ is positive (or negative, resp.) then the corresponding intersection of $`D_{12}^{\mathrm{}}`$ and $`D_{23}^u`$ is negative (resp. positive), see Figure 1. Thus we have the result.
Since $`D_{12}^{\mathrm{}}`$ is the union of circles in $`𝐑^4`$ disjoint from $`K_3`$, the linking number $`\text{Link}(D_{12}^{\mathrm{}},K_3)`$ is defined.
###### Lemma 2.2
$`\text{Tlk}_1(K_1,K_2,K_3)=\text{Link}(D_{12}^{\mathrm{}},K_3).`$
Proof. Without loss of generality, we may assume that the projection $`p`$ is given by $`p(w,x,y,z)=(x,y,z)`$. For a real number $`\lambda `$ we denote by $`t_\lambda :𝐑^4𝐑^4`$ the translation with $`t_\lambda (w,x,y,z)=(w+\lambda ,x,y,z)`$. Let $`M_3^{}`$ be a 3-chain in $`𝐑^4`$ with $`M_3^{}=K_3`$. Take a sufficiently large number $`R`$ and consider a $`3`$-chain
$$M_3=_{\lambda [0,R]}t_\lambda (K_3)+t_R(M_3^{})$$
so that $`M_3=K_3`$ and $`D_{12}^{\mathrm{}}M_3=D_{12}^{\mathrm{}}(_{\lambda [0,R]}t_\lambda (K_3))`$. The projection $`p`$ induces a one-to-one correspondence between the geometric intersection $`D_{12}^{\mathrm{}}(_{\lambda [0,R]}t_\lambda (K_3))`$ and the subset of $`D_{12}^{\mathrm{}}K_3^{}=p(D_{12}^{\mathrm{}})p(K_3)`$ consisting of points where $`D_{12}^l`$ is higher than $`K_3^{}`$ (in the over-under information of the surface diagram $`F^{}`$), i.e., the set of triple points of $`F^{}`$ of type $`(1,2,3)`$. Since the orientation of $`D_{12}^{\mathrm{}}`$ is parallel to the orientation of $`D_{12}^{\mathrm{}}`$, the sign of an intersection of $`D_{12}^{\mathrm{}}`$ and $`_{\lambda [0,R]}t_\lambda (K_3)`$ coincides with the sign of the corresponding intersection of $`D_{12}^{\mathrm{}}`$ and $`K_3^{}`$, which is the sign of the triple point (see Figure 2). Thus we have the result.
## 3 Proof of Theorem 1.1
For a compact oriented $`n`$-manifold $`M`$ with $`\{A,B\}=\{M,\mathrm{}\}`$, we denote by
$$_M:H_p(M,A)\times H_q(M,A)H_{p+qn}(M,A)$$
the intersection map, which is defined by
$$x_My=P_M(P_M^1(x)P_M^1(y))$$
where $`P_M:H^{}(M,B)H_n(M,A)`$ is the Poincaré duality isomorphism (see , page 391). We will use $``$ and $`P`$ instead of $`_M`$ and $`P_M`$ when their meanings are obvious in context.
Let $`F=K_1K_2K_3`$ be a $`3`$-component surface-link. For simplicity of argument, we assume that $`F`$ is embedded in the $`4`$-sphere $`S^4=𝐑^4\{\mathrm{}\}`$. For a regular neighborhood $`N_i`$ of $`K_i`$ in $`S^4`$, we put
$$E_i=\mathrm{Cl}(S^4\backslash N_i),E_{ij}=\mathrm{Cl}(S^4\backslash (N_iN_j))\text{ for }ij\text{, and }E=\mathrm{Cl}(S^4\backslash (N_1N_2N_3)),$$
where Cl denotes the closure. We denote by $`M_i`$ a $`3`$-chain in $`S^4`$ with $`M_i=K_i`$ for $`i=1,2,3`$ (the reader may suppose that it is a Seifert hypersurface for $`K_i`$, i.e., a compact oriented $`3`$-manifold embedded in $`S^4`$ with $`M_i=K_i`$). We also denote by $`M_i`$ the homology class in
$$H_3(S^4,K_i)H_3(S^4,N_i)H_3(E_i,E_i)$$
represented by $`M_i`$. By $`u_iH^1(E_i)`$ we denote the Poincaré dual of $`M_iH_3(E_i,E_i)`$, i.e., $`M_i=P(u_i)=u_i[E_i]`$. For a subset $`X`$ of $`E_i`$, we will denote by $`u_i|_XH^1(X)`$ the image of $`u_i`$ by the inclusion-induced homomorphism $`H^1(E_i)H^1(X)`$. Moreover, if $`X`$ is an $`n`$-manifold, we denote by $`M_i|_{(X,X)}`$ (or $`M_i|_X`$ if $`X=\mathrm{}`$) the Poincaré dual $`P_X(u_i|_X)=(u_i|_X)[X]H_{n1}(X,X)`$ of $`u_i|X`$.
For $`i\{1,3\}`$, since $`K_2E_i`$, $`M_i|_{K_2}H_1(K_2)`$ is defined. (When we consider $`M_i`$ as a $`3`$-chain, the intersection of $`M_i`$ and $`K_2`$ (as a $`1`$-cycle in $`K_2`$) represents $`M_i|_{K_2}`$.) Let
$$\text{Tlk}_2(K_1,K_2,K_3)=\epsilon _{K_2}(M_1|_{K_2}M_3|_{K_2}),$$
where $`\epsilon _{K_2}:H_0(K_2)𝐙`$ is the augmentation.
###### Lemma 3.1
$`\text{Tlk}_1(K_1,K_2,K_3)=\text{Tlk}_2(K_1,K_2,K_3).`$
Proof. We assume $`F𝐑^4S^4`$ and continue the situation of the proof of Lemma 2.2. Let $`M_1^{}`$ be a 3-chain in $`𝐑^4(S^4)`$ with $`M_1^{}=K_1`$ and consider a $`3`$-chain $`M_1`$ such that
$$M_1=_{\lambda [R,0]}t_\lambda (K_1)+t_R(M_1^{})$$
with $`M_1=K_1`$. The intersection of $`M_1`$ and $`K_2`$ is equal to that of $`_{\lambda [R,0]}t_\lambda (K_1)`$ and $`K_2`$ which is the 1-chain $`D_{12}^{\mathrm{}}`$ in $`K_2`$, and the intersection of $`M_3`$ and $`K_2`$ is equal to that of $`_{\lambda [0,R]}t_\lambda (K_3)`$ and $`K_2`$ which is the 1-chain $`D_{23}^u`$ in $`K_2`$. Therefore, by Lemma 2.1, we have
$$\begin{array}{cc}\text{Tlk}_1(K_1,K_2,K_3)\hfill & =\text{Int}_{K_2}(D_{12}^{\mathrm{}},D_{23}^u)\hfill \\ & =\epsilon _{K_2}(D_{12}^{\mathrm{}}D_{23}^u)\hfill \\ & =\epsilon _{K_2}(M_1|_{K_2}M_3|_{K_2})\hfill \\ & =\text{Tlk}_2(K_1,K_2,K_3).\text{ }\hfill \end{array}$$
Remark. The argument in Lemmas 2.2 and 3.1 implies that for any Seifert hypersurface $`M_1`$ for $`K_1`$, the intersection of $`M_1`$ and $`K_2`$ (as a 1-cycle in $`K_2`$) is homologous to $`D_{12}^{\mathrm{}}`$, and that for any Seifert hypersurface $`M_3`$ for $`K_3`$, the intersection of $`M_3`$ and $`K_2`$ (as a 1-cycle in $`K_2`$) is homologous to $`D_{23}^u`$.
We denote by $`[M_1K_2]_{E_3}H_1(E_3)`$ the homology class of the intersection $`M_1K_2`$ as a $`1`$-cycle in $`E_3`$. This is equal to the image of $`M_1|_{K_2}H_1(K_2)`$ under the inclusion-induced homomorphism $`H_1(K_2)H_1(E_3)`$ and also equal to the image of $`M_1|_{(E_{13},E_{13})}[K_2]_{E_{13}}H_1(E_{13})`$ under the inclusion-induced homomorphism $`H_1(E_{13})H_1(E_3)`$, where $`[K_2]_{E_{13}}H_2(E_{13})`$ is represented by $`K_2`$. Let
$$\begin{array}{cc}\text{Tlk}_3(K_1,K_2,K_3)\hfill & =\text{Link}([M_1K_2]_{E_3},K_3)\hfill \\ & =\epsilon _{E_3}([M_1K_2]_{E_3}M_3),\hfill \end{array}$$
where $`M_3H_3(E_3,E_3)`$ and $`\epsilon _{E_3}:H_0(E_3)𝐙`$ is the augmentation.
###### Lemma 3.2
$`\text{Tlk}_1(K_1,K_2,K_3)=\text{Tlk}_3(K_1,K_2,K_3).`$
Proof. In the situation of the proof of Lemma 2.2, $`[M_1K_2]_{E_3}H_1(E_3)`$ is represented by the $`1`$-cycle $`D_{12}^{\mathrm{}}`$. Therefore, by Lemma 2.2, we have
$$\begin{array}{cc}\text{Tlk}_3(K_1,K_2,K_3)\hfill & =\text{Link}([M_1K_2]_{E_3},K_3)\hfill \\ & =\text{Link}(D_{12}^{\mathrm{}},K_3)\hfill \\ & =\text{Tlk}_1(K_1,K_2,K_3).\text{ }\hfill \end{array}$$
We denote by $`[M_1M_3]_{(E_{13},E_{13})}H_1(E_{13},E_{13})`$ the class of the intersection $`M_1M_2`$ as a $`2`$-cycle in $`(E_{13},E_{13})`$ when we regard $`M_i`$ as a $`3`$-chain. This is equal to the intersection product $`M_1|_{(E_{13},E_{13})}M_3|_{(E_{13},E_{13})}H_2(E_{13},E_{13})`$. Let
$$\begin{array}{cc}\text{Tlk}_4(K_1,K_2,K_3)\hfill & =\epsilon _{E_{13}}([M_1M_3]_{(E_{13},E_{13})}[K_2]_{E_{13}})\hfill \\ & =\epsilon _{E_{13}}(M_1|_{(E_{13},E_{13})}M_3|_{(E_{13},E_{13})}[K_2]_{E_{13}}).\hfill \end{array}$$
###### Lemma 3.3
$`\text{Tlk}_3(K_1,K_2,K_3)=\text{Tlk}_4(K_1,K_2,K_3).`$
Proof. Let $`i_{}:H_{}(E_{13})H_{}(E_3)`$ and $`i^{}:H^{}(E_3)H^{}(E_{13})`$ be the inclusion-induced homomorphisms. Recall that $`[M_1K_2]_{E_3}=i_{}(M_1|_{(E_{13},E_{13})}[K_2]_{E_{13}})`$. Thus,
$$\begin{array}{cc}\text{Tlk}_4(K_1,K_2,K_3)\hfill & =\epsilon _{E_{13}}(M_1|_{(E_{13},E_{13})}M_3|_{(E_{13},E_{13})}[K_2]_{E_{13}})\hfill \\ & =\epsilon _{E_3}i_{}(M_1|_{(E_{13},E_{13})}M_3|_{(E_{13},E_{13})}[K_2]_{E_{13}})\hfill \\ & =\epsilon _{E_3}i_{}(M_3|_{(E_{13},E_{13})}M_1|_{(E_{13},E_{13})}[K_2]_{E_{13}})\hfill \\ & =\epsilon _{E_3}i_{}(u_3|_{E_{13}}(M_1|_{(E_{13},E_{13})}[K_2]_{E_{13}}))\hfill \\ & =\epsilon _{E_3}i_{}(i^{}(u_3)(M_1|_{(E_{13},E_{13})}[K_2]_{E_{13}}))\hfill \\ & =\epsilon _{E_3}(u_3i_{}(M_1|_{(E_{13},E_{13})}[K_2]_{E_{13}}))\hfill \\ & =\epsilon _{E_3}(M_3[M_1K_2]_{E_3})\hfill \\ & =\epsilon _{E_3}([M_1K_2]_{E_3}M_3)\hfill \\ & =\text{Tlk}_3(K_1,K_2,K_3).\text{ }\hfill \end{array}$$
For $`i\{1,2,3\}`$, since $`N_2E_i`$, $`M_i|_{N_2}H_2(N_2)`$ is defined. Let
$$\begin{array}{cc}\text{Tlk}_5(K_1,K_2,K_3)\hfill & =\epsilon _{N_2}(M_1|_{N_2}M_2|_{N_2}M_3|_{N_2})\hfill \\ & =<u_1|_{N_2}u_2|_{N_2}u_3|_{N_2},[N_2]>.\hfill \end{array}$$
###### Lemma 3.4
$`\text{Tlk}_4(K_1,K_2,K_3)=\text{Tlk}_5(K_1,K_2,K_3).`$
Proof. Let $`i:N_2N_2`$ be the inclusion map. In $`H_0(N_2)`$, we have
$$\begin{array}{cc}i_{}(M_1|_{N_2}M_2|_{N_2}M_3|_{N_2})\hfill & =i_{}(M_1|_{N_2}M_3|_{N_2}M_2|_{N_2})\hfill \\ & =i_{}(_{}(M_1|_{(N_2,N_2)})_{}(M_3|_{(N_2,N_2)})M_2|_{N_2})\hfill \\ & =i_{}(_{}(M_1|_{(N_2,N_2)}M_3|_{(N_2,N_2)})M_2|_{N_2})\hfill \\ & =(M_1|_{(N_2,N_2)}M_3|_{(N_2,N_2)})i_{}(M_2|_{N_2})\hfill \\ & =M_1|_{(N_2,N_2)}M_3|_{(N_2,N_2)}([K_2]_{N_2})\hfill \\ & =M_1|_{(N_2,N_2)}M_3|_{(N_2,N_2)}[K_2]_{N_2}.\hfill \end{array}$$
Thus
$$\text{Tlk}_5(K_1,K_2,K_3)=\epsilon _{N_2}(M_1|_{(N_2,N_2)}M_3|_{(N_2,N_2)}[K_2]_{N_2}).$$
It is obvious that
$$\epsilon _{N_2}(M_1|_{(N_2,N_2)}M_3|_{(N_2,N_2)}[K_2]_{N_2})=\epsilon _{E_{13}}(M_1|_{(E_{13},E_{13})}M_3|_{(E_{13},E_{13})}[K_2]_{E_{13}})$$
and hence we have the result.
###### Lemma 3.5
$`\mathrm{Tlk}_6(K_1,K_2,K_3)=\pm \mathrm{Tlk}_1(K_1,K_2,K_3).`$
Proof. Since $`\mathrm{Tlk}_6`$ is an ambient isotopy invariant, we may assume that the surface-link $`F=f(F_1)f(F_2)f(F_3)`$ is in general position with respect to the projection $`p:𝐑^4𝐑^3`$ with $`p(w,x,y,z)=(x,y,z)`$. The preimage of a particular point $`((1,0,0,0),(1,0,0,0))`$ by $`L`$ consists of triples $`(x_1,x_2,x_3)F_1\times F_2\times F_3`$ such that $`p(f(x_1))=p(f(x_2))=p(f(x_3))`$ and $`f(x_1)`$ is the upper, $`f(x_2)`$ is the middle, $`f(x_3)`$ is the lower lift of the triple point $`p(f(x_1))`$. For each such triple $`(x_1,x_2,x_3)`$, let $`D_T^2`$, $`D_M^2`$, $`D_B^2`$ be regular neighborhoods of them in $`F_1F_2F_3`$, and let $`\epsilon \{+1,1\}`$ be the sign of the triple point $`p(f(x_1))`$. Let $`(x_T,y_T)`$, $`(x_M,\epsilon z_M)`$ and $`(y_B,z_B)`$ be coordinate systems of $`D_T^2`$, $`D_M^2`$ and $`D_B^2`$ around $`x_1,x_2`$ and $`x_3`$, respectively. Modifying $`f`$ up to ambient isotopy, we may assume that the restriction of $`f`$ to $`D_T^2D_M^2D_B^2`$ is given by defined by
$`(x_T,y_T)`$ $``$ $`(0,x_0,y_0,z_0)+(3,x_T,y_T,0)`$
$`(x_M,\epsilon z_M)`$ $``$ $`(0,x_0,y_0,z_0)+(2,x_M,0,z_M)`$
$`(y_B,z_B)`$ $``$ $`(0,x_0,y_0,z_0)+(1,y_B,z_B)`$
where $`(x_0,y_0,z_0)𝐑^3`$ is the triple point $`p(f(x_1))`$. In this situation, the restriction
$$L^{}:D_T^2\times D_M^2\times D_B^2S^3\times S^3$$
is given by the formula
$$(\frac{(1,x_Tx_M,y_T,z_M)}{\sqrt{1+(x_Tx_M)^2+y_T^2+z_M^2}},\frac{(1,x_M,y_B,z_Mz_B)}{\sqrt{1+x_M^2+y_B^2+(z_Mz_B)^2}}).$$
The map $`L^{}`$ is injective and hence it is a homeomorphism onto its image. Its (local) degree is $`+1`$ or $`1`$ which depends only on $`\epsilon `$. Since the degree of $`L`$ is the sum of the (local) degrees of $`L^{}`$ for all triples $`(x_1,x_2,x_3)`$ in the preimage $`L^1((1,0,0,0),(1,0,0,0))`$, this number agrees up to sign with the triple linking number $`\mathrm{Tlk}_1(f(F_1),f(F_2),f(F_3)).`$
By Lemmas 3.1–3.5, we have Theorem 1.1.
## 4 Proof of Theorems 1.2–1.5
To prove Theorem 1.2, it is useful to change $`N_2`$ in the definition of $`\text{Tlk}_5`$ for $`E_2`$.
###### Lemma 4.1
$`\text{Tlk}_5(K_1,K_2,K_3)=\epsilon _{E_2}(M_1|_{E_2}M_2|_{E_2}M_3|_{E_2})`$, where the intersections are taken in $`E_2`$.
Proof. Since $`N_2`$ and $`E_2`$ are the same $`3`$-submanifold of $`S^4`$ with opposite orientations, $`[N_2]=[E_2]`$ in $`H_3(N_2)=H_3(E_2)`$. Thus, in $`H_0(N_2)=H_0(E_2)`$,
$$\begin{array}{cc}M_1|_{N_2}M_2|_{N_2}M_3|_{N_2}\hfill & =(u_1|_{N_2}u_2|_{N_2}u_3|_{N_2})[N_2]\hfill \\ & =(u_1|_{E_2}u_2|_{E_2}u_3|_{E_2})([E_2])\hfill \\ & =M_1|_{E_2}M_2|_{E_2}M_3|_{E_2}.\text{ }\hfill \end{array}$$
Proof of Theorem $`1.2`$ (i)
$$\begin{array}{cc}\text{Tlk}_2(K_1,K_2,K_3)\hfill & =\epsilon _{K_2}(M_1|_{K_2}M_3|_{K_2})\hfill \\ & =\epsilon _{K_2}(M_3|_{K_2}M_1|_{K_2})\hfill \\ & =\text{Tlk}_2(K_3,K_2,K_1).\hfill \end{array}$$
(ii) Note that $`M_i|_{(E,E)}=P_E(u_i|_E)H_3(E,E)`$ is the image of $`M_i`$ under
$$H_3(E_i,E_i)H_3(E_i,E_iN_jN_k)H_3(E,E),$$
and $`M_i|_{E_2}H_2(E_2)`$ is the image of $`M_i|_{(E,E)}`$ under
$$H_3(E,E)H_2(E)H_2(E_1)H_2(E_2)H_2(E_3)H_2(E_2),$$
the boundary operator followed by the projection to $`H_2(E_2)`$. We denote by $`(M_i|_{E_2})_EH_2(E)`$ the image of $`M_i|_{E_2}`$ under the inclusion-induced homomorphism $`H_2(E_2)H_2(E)`$. By Lemma 4.1,
$$\text{Tlk}_5(K_1,K_2,K_3)=\epsilon _E((M_1|_{E_2})_E(M_2|_{E_2})_E(M_3|_{E_2})_E).$$
Thus, we have
$$\begin{array}{cc}\hfill \text{Tlk}_5(K_1,K_2,K_3)+& \text{Tlk}_5(K_2,K_3,K_1)+\text{Tlk}_5(K_3,K_1,K_2)\hfill \\ \hfill =& \epsilon _E((M_1|_{E_2})_E(M_2|_{E_2})_E(M_3|_{E_2})_E\hfill \\ & +(M_2|_{E_3})_E(M_3|_{E_3})_E(M_1|_{E_3})_E\hfill \\ & +(M_3|_{E_1})_E(M_1|_{E_1})_E(M_2|_{E_1})_E)\hfill \\ \hfill =& \epsilon _E((M_1|_{E_2})_E(M_2|_{E_2})_E(M_3|_{E_2})_E\hfill \\ & +(M_1|_{E_3})_E(M_2|_{E_3})_E(M_3|_{E_3})_E\hfill \\ & +(M_1|_{E_1})_E(M_2|_{E_1})_E(M_3|_{E_1})_E)\hfill \\ \hfill =& \epsilon _E(M_1|_EM_2|_EM_3|_E)\hfill \\ \hfill =& \epsilon _E(_{}(M_1|_{(E,E)})_{}(M_2|_{(E,E)})_{}(M_3|_{(E,E)}))\hfill \\ \hfill =& \epsilon _E(_{}(M_1|_{(E,E)}M_2|_{(E,E)}M_3|_{(E,E)}))\hfill \\ \hfill =& 0.\text{ }\hfill \end{array}$$
Proof of Theorem $`1.3`$ (i) The intersection number between two oriented curves on a 2-sphere vanishes. By Lemma 2.1, we have $`\text{Tlk}(K_1,K_2,K_3)=0`$.
(ii) This is an immediate consequence of (i) and Theorem 1.2(ii).
Proof of Theorem $`1.4`$ If $`\text{Alk}(K_2,K_3)=0`$, then $`\text{Tlk}_3(K_1,K_2,K_3)=\text{Link}([M_1K_2]_{E_3},K_3)=0`$, for $`[M_1K_2]_{E_3}H_1(E_3)`$ is the image of $`M_1|_{K_2}H_1(K_2)=0`$. By Theorem 1.2, we have $`\text{Tlk}(K_3,K_2,K_1)=0`$.
We consider surface-links $`F=K_1K_2K_3`$ in which each $`K_i`$ is not necessarily connected. Such a surface-link is called a $`3`$-partitioned surface-link. The definition of the triple linking of $`F=K_1K_2K_3`$ is generalized directly for 3-partitioned surface-links, and all results and proofs in Sections 2 and 3 are valid for 3-partitioned surface-links. Theorem 1.5 is a special case of the following:
###### Theorem 4.2
If two $`3`$-partitioned surface-links $`F=K_1K_2K_3`$ and $`F^{}=K_1^{}K_2^{}K_3^{}`$ are link homologous, then $`\text{Tlk}(K_i,K_j,K_k)=\text{Tlk}(K_i^{},K_j^{},K_k^{})`$.
Proof. It is sufficient to prove $`\text{Tlk}(K_1,K_2,K_3)=\text{Tlk}(K_1^{},K_2^{},K_3^{})`$ in a special case that $`K_i=K_i^{}`$, $`K_j=K_j^{}`$ and $`K_k`$ is homologous to $`K_k^{}`$ in $`S^4\backslash (K_iK_j)`$, where $`\{i,j,k\}=\{1,2,3\}`$. If $`k=2`$, then $`\text{Tlk}_4(K_1,K_2,K_3)=\text{Tlk}_4(K_1,K_2^{},K_3)`$ by definition. If $`k=1`$, then $`\text{Tlk}_3(K_1,K_2,K_3)=\text{Tlk}_3(K_1^{},K_2,K_3)`$. (This is seen as follows: Let $`M_1`$ be a 3-chain with $`M_1=K_1`$. Since $`K_1^{}`$ is homologous to $`K_1`$ in $`S^4\backslash (K_2K_3)`$, there is a 3-chain $`B`$ in $`S^4\backslash (K_2K_3)`$ with $`B=K_1^{}K_1`$. Let $`M_1^{}=M_1+B`$, which is a 3-chain with $`M_1^{}=K_1^{}`$. Then $`[M_1K_2]_{E_3}=[M_1^{}K_2]_{E_3}`$ in $`H_1(E_3)H_1(S^4\backslash K_3)`$. ) The case $`k=3`$ is reduced to the previous case ($`k=1`$) by use of Theorem 1.2(i).
## 5 Proof of Proposition 1.6
(1) Let $`\mathrm{}=k_1k_2`$ be a $`(2,2a)`$-torus link in a $`3`$-disk $`D^3`$ with $`\text{Link}(k_1,k_2)=a`$. Let $`\gamma `$ be a simple loop in $`𝐑^4`$ which intersects a $`3`$-disk $`B_0`$ in $`𝐑^4`$ transversely at a single interior point of $`B_0`$ in the positive direction. Identify $`D^3\times S^1`$ with a regular neighborhood $`N(\gamma )`$ of $`\gamma `$ in $`𝐑^4`$ and let $`T_1T_2`$ be the image of $`\mathrm{}\times S^1=k_1\times S^1k_2\times S^1`$ in $`𝐑^4`$. Let $`F=K_1K_2K_3`$ be a surface-link with $`K_1=T_1`$, $`K_2=T_2`$ and $`K_3=B_0`$. Then $`F`$ is the desired link.
(2) Let $`\mathrm{}=k_1k_2k_3`$ be a pretzel link of type $`(2a,2b)`$ in a 3-disk $`D^3`$ so that $`\text{Link}(k_1,k_2)=a`$, $`\text{Link}(k_2,k_3)=b`$ and $`\text{Link}(k_1,k_3)=0`$. Let $`B_1,B_2,B_3`$ be mutually disjoint 3-disks embedded in $`𝐑^4`$ and let $`\gamma `$ be a simple loop in $`𝐑^4`$ which intersects $`B_i`$ ($`i=1,2,3`$) transversely at a single interior point of $`B_i`$ in the positive direction. Identify $`D^3\times S^1`$ with a regular neighborhood $`N(\gamma )`$ of $`\gamma `$ in $`𝐑^4`$ and let $`T_1T_2T_3`$ be the image of $`\mathrm{}\times S^1=k_1\times S^1k_2\times S^1k_3\times S^1`$ in $`𝐑^4`$. Let $`F=K_1K_2K_3`$ be a surface-link obtained from $`(T_1T_2T_3)(B_1B_2B_3)`$ by piping such that $`F`$ has a projection as in Figure 3. Then $`F`$ is the desired link.
## 6 Remarks
Remark 6.1 The definition of $`\text{Tlk}_6`$ can be seen as a direct analogue of (6) given in page 133. See also . This generalizes the triple linking to all link maps, instead of embeddings. Moreover, it easily generalized to all dimensions. Let $`M_i`$ denote a closed connected $`n`$-manifold for $`i=1,\mathrm{},n+1`$. Let an embedding $`f:_{i=1}^{n+1}M_i𝐑^{n+2}`$ be given. Define $`L:_{i=1}^{n+1}M_i_{j=1}^nS_j^{n+1}`$ as follows. Let $`x_iM_i`$; for $`i=1,\mathrm{},n`$, let $`\mathrm{\Delta }_i=f(x_i)f(x_{i+1})/f(x_i)f(x_{i+1})`$. Then
$$L(x_1,\mathrm{}x_{n+1})=(\mathrm{\Delta }_1,\mathrm{\Delta }_2,\mathrm{},\mathrm{\Delta }_n).$$
The general $`(n+1)`$-fold linking number, Glk, is defined by
$$\text{Glk}(f(M_1),\mathrm{},f(M_{n+1}))=\mathrm{deg}(L).$$
We can generalize the notion of $`\text{Tlk}_1`$ to a diagram in $`𝐑^{n+1}`$ of an $`(n+1)`$-component $`n`$-manifold-link $`M_1\mathrm{}M_{n+1}`$ in $`𝐑^{n+2}`$; namely, a diagram has generic $`(n+1)`$-tuple points and we count the number of times $`M_1`$ is over $`M_2`$ is over … is over $`M_{n+1}`$ with signs. It is difficult to show that this value is an invariant of the $`n`$-manifold-link in $`𝐑^{n+2}`$ directly, since we do not know Reidemeister moves for higher dimensions $`(n3)`$. However, the proof of Lemma 3.5 goes through to show that Glk is the same as this count (up to sign). Thus we have that this number (generalization of $`\text{Tlk}_1`$) is an invariant of an $`(n+1)`$-component $`n`$-manifold-link.
Remark 6.2 In classical link theory, the linking number determines the link homology classes completely. However, the triple linking of surface-links is not a complete invariant of the surface-link homology; there exists a pair of surface-links with the same triple linking invariants which are not link homologous. A classification of surface-link homology classes is discussed in a forthcoming paper.
Acknowledgments The authors would like to thank Akio Kawauchi for helpful suggestions. JSC is being supported by NSF grant DMS-9988107. SK and SS are being supported by Fellowships from the Japan Society for the Promotion of Science. MS is being supported by NSF grant DMS-9988101.
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# Thermodynamic and structural aspects of the potential energy surface of simulated water
## I Introduction
In recent years, numerical study of model liquids in supercooled states has been helpful to clarify the physics of the glass transition . The availability of long trajectories in phase space offers the possibility of closely examining the changes in supercooled states that are responsible for slowing down the dynamics by 15 decades in a narrow temperature range approaching the glass transition. Although current computational studies are limited to times shorter than $`1`$ $`\mu `$s (as opposed to real liquids, for which the dynamics can be studied up to $`10^2`$ s), a coherent picture of the glass transition phenomenon is beginning to emerge.
In addition to characterizing changes in the dynamics, recent studies have demonstrated the utility of examining the underlying potential energy surface (PES) as an aid to understanding the properties of supercooled liquids connecting the dynamics to the thermodynamics and the topology of configuration space . This connection is most obvious at low temperatures, where the motion of the system can be partitioned into motion confined within a single potential energy basin with infrequent inter-basin transitions. It has been shown that at sufficiently high temperature (at constant volume), the system explores always the same distribution of basins, and that the average basin energy is nearly temperature independent. Below a crossover temperature—which is coincident with the onset of a two-step relaxation in the decay of density fluctuations —the system starts to populate basins of progressively lower energy, but which are less numerous.
The thermodynamic approach based on the analysis of the PES, following the formalism proposed by Stillinger and Weber, has become a powerful formalism for the interpretation of numerical data, both in equilibrium and in out-of-equilibrium conditions . The degeneracy of the energy minima, i.e., the number of basins with a selected minimum energy, has been quantified for several model systems and used to calculate the configurational entropy $`S_{\text{conf}}`$, from which an “ideal” glass transition (in the sense of Kauzmann, Adam, Gibbs, and DiMarzio ) has been estimated .
In this paper we present a detailed investigation of the properties of local potential energy minima, or “inherent structures” (IS), sampled by the extended simple point change (SPC/E) model of water . Previous studies on the PES for models of water have shown the relevance of this approach to the deepening of our understanding of both structural and dynamical properties of liquid water . Specifically, we calculate the temperature and density dependence of the basin energy over a wide range of temperatures and densities. We also study in detail the shape of the basins in configuration space and estimate their degeneracy. The information presented provides a detailed characterization of the PES and furnishes all information required to explicitly write the liquid Helmholtz free energy in a wide temperature and density range. Finally, we study the geometrical arrangement of the molecules at IS minima to better understand the changes that take place in the liquid on cooling. Focusing on the IS allows us to eliminate thermal effects which complicate the temperature dependence. In particular, we focus on the fraction of water molecules that are four-coordinated.
## II Simulations
The majority of the state points studied here are from the molecular dynamics simulations of the SPC/E model performed in Ref. . The simulation methods are discussed in Ref. . We have performed additional simulations of ice $`I_h`$, so that we can compare the IS properties of the crystal with those of the liquid. The simulations of ice consist of a periodic box containing 432 molecules with dimensions $`2.634\text{nm}\times 2.281\text{nm}\times 2.151`$ nm for density $`\rho =1.0`$ g/cm<sup>3</sup>. The dimensions are uniformly scaled in order to obtain other densities. The box dimensions have been optimized to generate the lowest energy configuration at density $`1.0`$ g/cm<sup>3</sup>. Proton disorder in the initial configuration is generated by identifying closed hydrogen bond loops, and exchanging hydrogens between molecules, as described in Ref. .
The ice simulations have been performed at $`\rho =0.90`$, 0.95, 1.00, and 1.05 g/cm<sup>3</sup> and temperature $`T=194`$ K. The thermodynamic properties are summarized in Table I. The equilibration time for these sample is far less than that of the liquid at the same temperature since only the vibrational degrees of freedom need to be relaxed.
We have performed conjugate gradient minimizations to locate local minima on the PES closest to any given instantaneous configuration. We use a tolerance of $`10^{15}`$ kJ/mol in the total energy for the minimization. For each state point, we quench at least 100 configurations taken from two independent trajectories. While each configuration is not necessarily separated by the typical relaxation time of the system, the set of points quenched typically spans $`25`$ times the relaxation time of the intermediate scattering function.
## III Inherent Structure Properties of SPC/E
### A Thermodynamics in the Inherent Structure Formalism
Stillinger and Weber formalized the concept of a basin in the potential energy surface by introducing the inherent structure formalism . The set of points that map to the same minimum via steepest descent are those which constitute a basin, and the minimum of a basin is the IS. This approach is particularly well suited to simulated liquids, since it is possible to explicitly calculate the steepest descent trajectory from an equilibrium state point. Moreover, the partition function $`Z`$ can be explicitly written in terms of the basins. In the isochoric-isothermal (NVT) ensemble, for a system of $`N`$ rigid molecules
$$Z=\lambda ^{6N}\mathrm{exp}(V/k_BT)d^N𝐫$$
(1)
which can be written as a sum over all basins in configurational space, i.e.
$`Z`$ $`=`$ $`\lambda ^{6N}{\displaystyle \underset{\text{basins}}{}}\mathrm{exp}(e_{IS}/k_BT)\times `$ (3)
$`{\displaystyle _{R_{\text{basin}}}}\mathrm{exp}((Ve_{IS})/k_BT)d^N𝐫.`$
Here $`\lambda h(2\pi mk_BT)^{1/2}`$ is the de Broglie wavelength, $`VV(𝐫^N)`$ is the potential energy as a function of the atomic coordinates, $`e_{IS}`$ is the energy of the IS and $`R_{basin}`$ is the configuration space associated to a specific basin. The model system we consider here, namely SPC/E water, has six degrees of freedom for each molecule. It is natural to introduce $`\mathrm{\Omega }(e_{IS})`$, the number of minima with energy $`e_{IS}`$, and the free energy of a basin with basin energy $`e_{IS}`$ $`f(T,e_{is})`$ (the “basin free energy”).
$`f(T,e_{IS})`$ $``$ $`k_BT\mathrm{ln}({\displaystyle \frac{1}{\mathrm{\Omega }(e_{IS})}}\lambda ^{6N}\times `$ (5)
$`{\displaystyle \underset{\text{basins}}{}^{}}{\displaystyle _{R_{\text{basin}}}}\mathrm{exp}[(Ve_{IS})/k_BT]d^N𝐫),`$
The asterisk denotes the fact that the sum is constrained to basins of energy $`e_{IS}`$. Eq. (5) accounts for both the basin structure surrounding the minimum and the kinetic degrees of freedom. The complete partition function can be written as the sum over all possible $`e_{IS}`$ values,
$$Z=\underset{e_{IS}}{}\mathrm{\Omega }(e_{IS})\mathrm{exp}\left(\frac{e_{IS}+f(T,e_{IS})}{k_BT}\right)$$
(6)
or
$$Z=\underset{e_{IS}}{}\mathrm{exp}\left(\frac{TS_{\text{conf}}(e_{IS})+e_{IS}+f(T,e_{IS})}{k_BT}\right)$$
(7)
where the configurational entropy $`S_{\text{conf}}(e_{IS})k_B\mathrm{ln}(\mathrm{\Omega }(e_{IS}))`$. In the thermodynamic limit, the corresponding Helmholtz free energy $`F(T,V)`$ is given by
$$F(V,T)=E_{IS}(T)+f(T,E_{IS}(T))TS_{\text{conf}}(E_{IS}(T))$$
(8)
where $`E_{IS}(T)`$ is the thermodynamic average of $`e_{IS}`$ and solves
$`F(V,T)/e_{IS}`$ $`=`$ $`1+f(T,e_{IS})/e_{IS}`$ (10)
$`TS_{\text{conf}}(T,e_{IS})/e_{IS}=0.`$
$`E_{IS}(T)`$ can be numerically calculated by estimating the $`IS`$ which are populated by a system in equilibrium at temperature $`T`$ and fixed volume $`V`$. Hence, if a good model for $`f(T,e_{IS})`$ is available, then the $`e_{IS}`$ dependence of $`S_{\text{conf}}`$ along isochores can be estimated by integrating Eq. (10). Note that Eq. (10) shows that, if the basin free energy does not depend on $`e_{IS}`$, then the configurational entropy is the only quantity controlling the $`T`$-dependence of $`E_{IS}`$. In other words, the statistical mechanics of the basins completely decouples from the vibrational dynamics .
$`S_{\text{conf}}(E_{IS})`$ can also be calculated by studying the probability distribution $`P(E_{IS},T)`$, i.e. the probability that the liquid—in equilibrium at temperature $`T`$—populates the inherent structure $`E_{IS}`$. Indeed, from Eq. (7)
$`S_{\text{conf}}(e_{IS})`$ $`=`$ $`k_B\mathrm{ln}P(e_{IS},T)+e_{IS}/T+`$ (12)
$`f(e_{IS},T)+k_B\mathrm{ln}Z(T).`$
Hence, Eq. 12 gives $`S_{\text{conf}}`$ up to an unknown constant $`k_B\mathrm{ln}Z(T)`$. This “histogram technique” has been recently used to estimate the configurational entropy for a binary-mixture Lennard-Jones system .
### B Inherent Structure Energy as a function of $`\rho `$ and $`T`$
The $`\rho `$\- and $`T`$-dependence of $`E_{IS}`$ for the SPC/E potential for a more limited range of $`T`$ was recently reported in Ref. . Our results for the IS energies are shown in Fig. 1(a) as a function of $`\rho `$ and in Fig. 1(b) as a function of $`T`$.
The IS energy of the liquid along isotherms shows hints of negative curvature for $`T230`$ K (Fig. 1(a)). The presence of this negative curvature can also be observed in the instantaneous equilibrium configurations . This curvature yields a negative contribution to compressibility $`K_T^1=V[(^2U/V^2]_TT(^2S/V^2)_T]`$, which might be related to a low-temperature critical point in SPC/E . Fig. 1 also shows the IS energy for ice $`I_h`$ (which coincides with the ground state energy). At the lowest $`T`$ studied, $`E_{IS}`$ of the disordered liquid is still significantly greater than that of ice . Note that ice $`I_h`$ is not the thermodynamically stable crystalline form for SPC/E , and the stable SPC/E crystal would most likely have a lower ground state energy. However, the stable SPC/E crystal does not correspond to any of the experimentally-known forms of ice . Fig 1(a) also show the Kauzmann energy $`E_K`$, which we discuss in sec. III D.
Fig. 1(b) shows that at high $`T`$, $`E_{IS}`$ is nearly $`T`$ independent. For $`T350400`$K, there is a rapid decrease of $`E_{IS}`$ with a weak density-dependence. At low $`T`$, $`E_{IS}`$ depends linearly on $`1/T`$, as shown in Fig. 1(c). The $`1/T`$ dependence of $`E_{IS}`$ can be derived from the partition function provided $`\mathrm{\Omega }(e_{IS})`$ is Gaussian, and that the basin free energy does not depend on $`e_{IS}`$ . Furthermore, as found for the Lennard-Jones liquids , the $`T`$ at which the IS energy starts to decrease correlates with the temperature at which a two-step relaxation starts to be observed in all characteristic correlation functions (see, e.g., Fig. 13 of Ref. )
### C The Basin Free Energy: Density of States
We next focus on the shape of the IS basin with the aim of developing a model for the basin free energy. In the harmonic approximation, the basin free energy is given by
$$F(E_{IS},T)=k_BT\underset{i=1}{\overset{6N3}{}}\left[\mathrm{ln}(\mathrm{}\omega _i/k_BT)\right],$$
(13)
the free energy of a harmonic oscillator with frequency spectrum $`\omega _i`$. The values $`\omega _i^2`$ are the eigenvalues of the Hessian matrix, defined by the second derivative of the potential energy with respect to the molecular degrees of freedom at the basin minimum. The mass and moments of inertia of a molecule have also been absorbed in the definition of $`\omega _i`$. The distribution of $`\omega _i`$, called the density of states (DOS), is shown in Fig. 2(a) for three different state points at $`T=210`$ K. The pronounced minimum at $`\omega 400`$ cm<sup>-1</sup> separates the translational modes (at lower frequencies) from the rotational modes (at higher frequencies). At larger $`\rho `$, the peaks in the DOS broaden due to the disruption of the H-bond network, which we will discuss in Sec. IV. For comparison, we also show the DOS for ice $`I_h`$, where we see a clear separation between the translational and rotational modes.
The normal mode spectrum contributes to the basin free energy via the term $`6\mathrm{ln}(\mathrm{}\omega )`$ (per molecule), where the brackets denote an average over the DOS and over different configurations. The dependencies on both $`T`$ and $`\rho `$ of $`\mathrm{ln}(\mathrm{}\omega )`$ are shown in Fig. 2(b). The dependence of $`6k_B\mathrm{ln}(\mathrm{}\omega )`$ on $`E_{IS}`$ is shown in Fig. 2(c). The average basin frequency is larger in deeper basins, showing that the basins become increasingly “sharp” on cooling. This is in contrast with the Lennard Jones case, where the basins become broader on cooling . The average curvature of the IS basin at high density has a weaker $`E_{IS}`$-dependence (and hence $`T`$-dependence) than those at low density, but are generally larger than the curvature at low density.
Fig. 3 shows the harmonic free energy estimate of Eq. (13) as a function of $`\rho `$. The range of values of the vibrational free energy of Fig. 3 (6 to 9 kJ/mol) is not very different from the range of values of $`E_{IS}`$ of Fig. 1(a) (from $`55`$ to $`60`$ kJ/mol); thus both make a significant contribution to the free energy of Eq. (8).
The basin free energy estimated using the harmonic approximation of Eq. (13) can be used as a starting point for estimating the true basin free energy. For a more precise quantification of the basin free energy, we must consider anharmonic contributions to the free energy. Indeed the basins of the SPC/E are anharmonic. This can easily be seen by considering the difference $`U_{\text{vib}}UE_{IS}`$; for a molecular system in the harmonic approximation, $`U_{\text{vib}}=(6/2)k_BT`$. Fig. 4 shows a marked deviation from harmonic behavior. However, as we discuss next, the anharmonicity does not have a strong $`E_{IS}`$ dependence. Development of techniques for the estimation of the anharmonic contribution to the basin free energy would be very useful.
### D Basin Degeneracy
A key element in the description of the configuration space in the IS thermodynamic formalism is $`\mathrm{\Omega }(e_{IS})`$, the number of basins with energy $`e_{IS}`$. The corresponding configurational entropy $`S_{\text{conf}}`$—i.e. the logarithm of $`\mathrm{\Omega }(e_{IS})`$—can be calculated by integrating Eq. (10). Using the harmonic approximation of Eq. (13) for the basin free energy, we obtain
$`S_{\text{conf}}(E_{IS})`$ $`=`$ $`S_{\text{conf}}(E_0)+{\displaystyle _{E_0}^{E_{IS}}}{\displaystyle \frac{dE_{IS}}{T}}+`$ (15)
$`k_B\mathrm{ln}(\omega (E_{IS})/\omega ((E_0)).`$
Fig. 5(a) separately shows the contribution from $`𝑑E_{IS}/T`$, and the contribution associated with $`\mathrm{ln}(\omega (E_{IS})/\omega ((E_0))`$. The harmonic contribution is not negligible. To obtain the $`S_{\text{conf}}`$ in absolute scale, the value of $`S_{\text{conf}}`$ at the reference point $`E_0`$ is needed. We have used reference values obtained in Ref. by independently calculating the absolute value of the entropy via thermodynamic integration from the known ideal gas reference point.
The complete results for all densities are shown in Fig. 5(b). In the same graph we show a fit to $`S_{\text{conf}}(E_{IS})`$ using the form
$$S_{\text{conf}}(E_{IS})=A(E_{IS}E_K)^2+B(E_{IS}E_K)$$
(16)
where $`E_K`$, the Kauzmann energy, is the $`E_{IS}`$ value at which the configurational entropy vanishes. This functional form is equivalent to the Gaussian distribution of $`\mathrm{\Omega }(e_{IS})`$ . The resulting $`E_{IS}`$ values, which provide an indication of the $`\rho `$ dependence of $`E_K`$, are reported in Fig. 1(a). The $`E_K`$ values suggest that disordered states with energies comparable to the ordered crystalline states may be available to this system , provided the $`E_K`$ estimates are not significantly affected by unknown errors in the extrapolation procedure or in the reference value for $`S_{\text{conf}}(E_0)`$ used.
To confirm independently the validity of the approach followed to estimate $`E_{IS}`$ dependence of $`S_{\text{conf}}`$, we compare the value obtained from Eq. (15) with the value obtained using Eq. (12) in Fig. 5(c). The agreement between the curves is remarkable. Moreover, the overlap for different $`P(e_{IS},T)`$ distributions after the harmonic $`E_{IS}`$ dependence is taken into account, suggests that there are no other systematic $`E_{IS}`$-dependent contributions. Any remaining $`T`$-dependent contributions are absorbed in the unknown $`Z(T)`$ function.
The results reported in this section provide a detailed analysis of the inherent structures and of their basins. This information can be used to develop a detailed free energy expression for the SPC/E . Calculation of the $`T`$ dependence of $`S_{\text{conf}}`$ was presented in Ref. , to probe the relation between the configurational contributions to thermodynamic quantities and the liquid dynamics in supercooled states.
## IV Structure
On lowering $`T`$ the dramatic changes in the IS energies are known to be accompanied by equally dramatic changes in the dynamic properties of the instantaneous configurations. However, examination of simple measures of structure of the instantaneous configurations (such as the pair distribution function) do not reveal such obvious changes; rather, there is a very gradual change, with the structure becoming slowly more well defined. A more careful analysis of structure is is required to see significant changes . In the following, we focus our attention on the structural changes that can be observed by studying the IS to try to obtain a more clear picture of the structural evolution of the system on cooling. The results we present are complementary to the results recently reported for the same system along similar lines of thought and expand on previous work .
### A Pair distribution function
Fig. 6 shows the oxygen-oxygen pair correlation function for both the equilibrated liquid and for the inherent structures at various $`T`$ and $`\rho `$. We first focus our attention on the $`T`$ dependence along the $`\rho =1.00`$ g/cm<sup>3</sup> isochore. For equilibrated configurations as well as inherent structures, the first and second peaks of the pair correlation function become better defined upon decreasing the temperature. Further, there is a systematic reduction of the intensity between the first and second neighbor peaks. At higher density ($`\rho =1.40`$ g/cm<sup>3</sup>), the behavior is somewhat different; for the instantaneous configurations, there is a clear $`T`$ dependence, while for the IS, the $`T`$ dependence is nearly negligible. For comparison, the model atomic liquid studied in , shows virtually no $`T`$ dependence of $`g(r)`$ in the IS (Fig 7). Hence the behavior at large $`\rho `$ is more like that expected for a simple liquid, consistent with the disappearance at large $`\rho `$ of many of the anomalies that differentiate water at ambient density.
We also show the behavior of $`g(r)`$ for various $`\rho `$ at $`T=210`$ K in the bottom panels of Fig. 6. On increasing $`\rho `$, both the instantaneous and inherent structure configurations show an increase in intensity at $`r0.32`$ nm, a trait already known to develop due to distortion (and eventual interpenetration) of the hydrogen bond network. For water, it is known that the preferred nearest-neighbor geometry is tetrahedral. This tetrahedral ordering is obvious at low densities, where peaks at $`0.28`$ nm and $`0.45`$ nm are the expected distances for a perfect tetrahedral lattice with first neighbor separation of $`0.28`$ nm. To see more clearly the tetrahedral nature of the liquid at all these temperatures and densities, we focus on the neighbor statistics at the various state points.
### B Neighbor Changes
We first consider the the average number of neighbors that a molecule has within a sphere of radius $`r`$, which is obtained from integration of $`g(r)`$:
$$n(r)=4\pi \rho _0^rr^2g(r^{})𝑑r^{}.$$
(17)
We show $`n(r)`$ in Fig. 8 as a function of $`\rho `$ at the lowest $`T`$ studied. We see a plateau in $`n(r)`$ almost exactly equal to four for all densities. Therefore, even at large $`\rho `$, the liquid has short range tetrahedral order. However, the rapid growth of $`n(r)`$ at large $`\rho `$ highlights the distortion and interpenetration.
To quantify the $`T`$ dependence of the structural changes we calculate the distribution of the number of neighbors a molecule has within a distance $`0.31`$ nm (arrow in Fig. 8), roughly corresponding to the first minimum in $`g(r)`$ at low density, shown in Fig. 6. We choose the first minimum in $`g(r)`$ at low density to emphasize the tetrahedrality of the liquid; for ice, all molecules would have 4 neighbors, while in high density liquid configurations, the distortion leads to a significant number of molecules having more or less than 4 neighbors in the first shell. Fig. 9 shows the histograms of the fraction of molecules with a given coordination number. At low density, the histogram for the instantaneous configurations changes from a rather broad one to one that is peaked around the value $`4`$ as the temperature decreases, as we expect. The same trend is visible for the inherent structures, although even at high temperatures, the distribution is quite narrowly peaked around the value $`4`$. Such a comparison permits us to make a separation between deviations from four-coordination arising from thermal agitation, and that arising from configurational change. This less marked $`T`$ dependence of the IS can be seen for all $`\rho `$. At larger $`\rho `$, the distribution is still peaked around 4, but there is significant fraction of molecules with are not four-bonded — these are molecules are most likely participating in the so called bifurcated bonds .
To quantify the changes in the tetrahedrality as a function of $`T`$ for each $`\rho `$, we plot the fraction of four-bonded molecules $`f_4`$ as a function of $`T`$ in Fig. 10. Those molecules which are not four bonded represent the set of bifurcated bonds . At low density, it is interesting to note that for both the equilibrated configurations and inherent structures, this $`f_4`$ is close to $`1`$ at the lowest $`T`$ simulated. Indeed, for $`T`$ somewhat lower that we can currently simulate, it appears that all molecules would be four bonded, and hence form a perfect random tetrahedral network. A simple extrapolation of $`f_4`$ for the lowest densities, displayed in Fig. 11, shows that $`T(f_4=1)`$ appears slightly lower than the mode-coupling transition temperature $`T_{MCT}`$ but well above $`T_K`$. Additionally, it appears that both the instantaneous and IS configurations appear to reach a random tetrahedral network at roughly the same temperature. The close correspondence of $`T(f_4=1)`$ and $`T_{MCT}`$ suggests that this crossover in structural change may be the controlling mechanism for the crossover in dynamic properties at $`T_{MCT}`$ . However, the extrapolation does not allow us to unambiguously associate these temperatures. Since the tetrahedral geometry is the “ideal” configuration at these low densities, there are unlikely to be any significant structural changes in the liquid for $`T<T(f_4=1)`$; to the extent that there is further structural change for $`T<T(f_4=1)`$, the rate of change must be significantly different than for $`T>T(f_4=1)`$. Analysis of the neighbor statistics in the IS at lower $`T`$ might shed some light on the hypothesized change in the dynamics from that of a fragile liquid to a strong liquid, a widely debated topic .
At larger densities increases, it is apparent that a random tetrahedral network is never reached. This is apparent in the fact that there are far fewer four-bonded molecules in the first neighbor shell at the higher densities. Hence, no sharp change in the structural development as a function of $`T`$ is expected. Consequently any crossover from fragile to strong behavior must become less pronounced.
## V Conclusions
We have characterized the properties of the PES basins in the configuration space of the SPC/E model over a range of densities in the supercooled regime. We have shown that the $`T`$ dependence of $`E_{IS}`$ is qualitatively very similar to the previous observations of simple liquids. We have also shown the importance of accounting for the shape of the basins when characterizing the thermodynamic properties of the IS subsystem. In particular, we found the there are significant $`T`$ dependent anharmonic effects. This detailed information on the basins should be useful for any future studies that focus on the dynamics of exploring the PES, or the thermodynamics at low $`T`$. Such a classification of basin properties would also be useful for other model glass formers, and might help to highlight the differences between fragile (such as OTP) and strong (such as SiO<sub>2</sub>) glass formers.
In addition, we presented results for the structural changes of the IS; these changes are more pronounced than what can be seen from equilibrium configurations. In particular, the progression of the structure toward a random tetrahedral network at low $`T`$ holds promise for a more physical and intuitive understanding of the glass transition in water. Below the temperature of crossover to a tetrahedral network, no further structural arrangement is expected, and the configurational entropy of the liquid may be nearly “frozen in” at the value that corresponds to the random tetrahedral network (plus the residual defects that would be present at concentrations varying with density). Thus the rate of change in entropy of the liquid may be expected change substantially near the crossover temperature, resulting in a significantly lower value of $`T_K`$, compared to the value that may be expected if the rate of change above the crossover temperature persisted. Similarly, because of the significant temperature dependence of the fraction of bifurcated bonds – which facilitate structural rearrangement – above the cross-over temperature, and the relative constancy below, the temperature dependence of the dynamical properties may show a corresponding crossover.
Finally, we call attention on the possibility of studying physical aging in this model-system, starting from the thermodynamic description of this system. We note that, at $`\rho =1.40`$ g/cm<sup>3</sup>, there is almost no variation of the basin curvature (Fig. 2) on $`E_{IS}`$, nor does the structure change significantly (Fig. 6). Hence these high density state points may offer an ideal opportunity to check if the out-of-equilibrium dynamics at very low $`T`$ can be still related to equilibrium states of the system .
## VI Acknowledgments
We thank P.G. Debenedetti, S.G. Glotzer, and F.H. Stillinger for helpful discussions. F.W.S. is supported by the National Research Council. F.S. is supported by INFM-PRA-HOP and Iniziativa Calcolo Parallelo and from MURST PRIN 98. This work was also supported by the NSF.
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# WKB Wave Functions with the Induced Gravity Theory
## I Introduction
It is known that the Hartle-Hawking quantum cosmology is based on Einstein’s general relativity. Although the latter is extremely successful at describing the observable universe, it doesn’t fully incorporate Mach’s principle, which demands that spacetime is determined entirely by background matter fields and physical laws. Hence the other type of theories, notably the Brans-Dicke theory and the induced gravity theory, in which the gravitational and cosmological “constants” can result from a scalar field. In such a gravity theory, both the the gravitational and cosmological “constants” are dynamical and time-dependent quantities. Observationally, there exist a number of experimental constraints on the time variation of the gravitational “constant”, $`G`$, of which the tightest bound, $`|\dot{G}/G|=(0.6\pm 2.0)\times 10^{12}\mathrm{yr}^1`$, was found by Thorsett using Bayesian statistical techniques on the measurements of the masses of young and old neutron stars in pulsar binaries. The upper limits on the time variation of the cosmological “constant”, $`\mathrm{\Lambda }`$, can be deduced from number counts of faint galaxies, statistical properties of gravitational lensing, structure formation and other ways. All results indicate an almost constant $`\mathrm{\Lambda }`$. Moreover, many attempts have been made in order to develop a plausible model in which the cosmological constant $`\mathrm{\Lambda }`$ is set to be precisely zero. So a critical problem is how the universe acquires almost constant values for $`G`$ and $`\mathrm{\Lambda }`$, and especially, a vanishing value for the latter. Because quantum cosmology (for an elegant review, see Ref. ) could, with no more than physical laws, provide a scheme which explains the present universe is what it is, we should endeavor to solve the problems mentioned above within the framework of such theory. In this context, We have considered the Brans-Dicke theory in previous papers, here we consider the induced gravity theory.
The first consideration of the quantum cosmology based on the induced gravity theory was given by Mo and Fang using the path integral technique. In this letter, we try to tackle the problem using the canonical quantization method, concentrating particularly on the time variation of $`G`$ and $`\mathrm{\Lambda }`$. The Wheeler-Dewitt equation (WDWE) is constructed in the minisuperspace approximation and the wave functions of the universe are obtained using three kinds of boundary condition, that are proposed by Hartle & Hawking, Linde and Vilenkin respectively. We shall show that no matter how $`G`$ and $`\mathrm{\Lambda }`$ vary in the classical models, they will acquire constant values when the universe comes from quantum creation. Moreover, the amplitude of the resulting wave function under the Vilenkin or Linde boundary condition sharply peaks around the classical trajectory only for a vanishing cosmological constant.
## II The Wheeler-Dewitt equation for the Induced <br>Gravity Theory
The action of induced gravity is
$$S=d^4x\sqrt{g}\left[\frac{1}{2}ϵ\phi ^2R\frac{1}{2}g^{\mu \nu }\phi _{,\mu }\phi _{,\nu }V(\phi )\right]+\mathrm{}.$$
(1)
So the scalar field $`\phi `$ induces a universe where the gravitational and cosmological “constants” are given simutaneously by
$$\begin{array}{cc}\hfill (16\pi G_{\mathrm{ind}})^1& =\frac{1}{2}ϵ\phi ^2,\hfill \\ & \\ \hfill \mathrm{\Lambda }_{\mathrm{ind}}& =\frac{V(\phi )}{ϵ\phi ^2}.\hfill \end{array}$$
(2)
We will not assume a specific form for $`V(\phi )`$, except that it is of the induced gravity type, for example, $`V(\phi )=\frac{\lambda }{8}(\phi ^2\upsilon ^2)^2`$. Note that $`ϵ,\lambda ,`$ and $`\upsilon `$ are all small constants. For quantum cosmology, gravitation is always dominant, and we may neglect terms representing other fields in the action, Eq.1. Under the minisuperspace approximation, the metric of spacetime is given by
$$ds^2=N(t)^2dt^2+a(t)^2d\mathrm{\Omega }_3^2,$$
(3)
where $`d\mathrm{\Omega }_3^2`$ is the line element of the three dimensional unit sphere, $`N(t)`$ is the lapse function. The scalar gravitational field $`\phi `$ depends on $`t`$ only. The total action can thus be written as
$$S=𝑑t\mathrm{\hspace{0.17em}2}\pi ^2\left[3ϵ\frac{a}{N}\phi ^2\dot{a}^26ϵ\frac{a}{N}a\phi \dot{a}\dot{\phi }+\frac{1}{2}\frac{a}{N}a^2\dot{\phi }^2+3ϵ\frac{N}{a}a^2\phi ^2\frac{N}{a}a^4V(\phi )\right],$$
(4)
where the dot stands for derivatives with respect to $`t`$. The momenta conjugate to $`a`$ and $`\phi `$ are defined in usual way and are respectively given by
$$\begin{array}{cc}\mathrm{\Pi }_a\hfill & \frac{\delta S}{\delta \dot{a}}=2\pi ^2\left(6ϵ\frac{a}{N}\phi ^2\dot{a}6ϵ\frac{a}{N}a\phi \dot{\phi }\right),\hfill \\ & \\ \mathrm{\Pi }_\phi \hfill & \frac{\delta S}{\delta \dot{\phi }}=2\pi ^2\left(6ϵ\frac{a}{N}a\phi \dot{a}+\frac{a}{N}a^2\dot{\phi }\right).\hfill \end{array}$$
(5)
Then we have the following relations,
$$\begin{array}{cc}a\dot{\phi }=\hfill & \frac{N}{2\pi ^2a}\frac{\phi \mathrm{\Pi }_\phi a\mathrm{\Pi }_a}{(6ϵ+1)a\phi },\hfill \\ & \\ \dot{a}\phi =\hfill & \frac{N}{2\pi ^2a}\frac{6ϵ\phi \mathrm{\Pi }_\phi +a\mathrm{\Pi }_a}{6ϵ(6ϵ+1)a\phi }.\hfill \end{array}$$
(6)
Taking the variation of the action, Eq.4, with respect to the lapse function $`N`$ and combining with Eq.6, we obtain the Hamiltonian constraint for the induced gravity theory,
$$H=\left(12ϵ(6ϵ+1)a\phi ^2\right)^1\left(\mathrm{\Pi }_{a}^{}{}_{}{}^{2}6ϵ\frac{\phi ^2}{a^2}\mathrm{\Pi }_{\phi }^{}{}_{}{}^{2}+12ϵ\frac{\phi }{a}\mathrm{\Pi }_a\mathrm{\Pi }_\phi \right)\left(2\pi ^2\right)^2a[3ϵ\phi ^2a^2V(\phi )]=0.$$
(7)
By introducing the canonical quantization into the above Hamiltonian constraint, $`\mathrm{\Pi }_a\frac{1}{i}\frac{}{a},\mathrm{\Pi }_\phi \frac{1}{i}\frac{}{\phi }`$, we obtain the WDWE for the induced gravity theory,
$$\left\{\frac{^2}{a^2}6ϵ\frac{\phi ^2}{a^2}\frac{^2}{\phi ^2}+12ϵ\frac{\phi }{a}\frac{}{a}\frac{}{\phi }+48\pi ^4ϵ(6ϵ+1)a^2\phi ^2[a^2V(\phi )3ϵ\phi ^2]\right\}\mathrm{\Psi }(a,\phi )=0.$$
(8)
In constructing the WDWE (Eq.8), we have ignored the operator order problem which is not important in the following discussion.
## III WKB Wave Functions of the Universe
In order to make predictions, we should solve the WDWE Eq.8, which proves to be a very difficult task. A regulous solution seems difficult, so we look for some simple solutions that do not depend sensitively on $`\phi `$. We neglect the second and third terms in Eq.8 to get cosmic wave functions. The problem is then simplified to a standard one-dimensional WKB problem for the scale factor $`a`$ with a potential
$$U(a)=48\pi ^4ϵ(6ϵ+1)a^2\phi ^2[3ϵ\phi ^2a^2V(\phi )],$$
(9)
For the classicall allowed (oscillatory) region $`aa__H\left[3ϵ\phi ^2/V(\phi )\right]^{1/2}`$, where the scale factor is large, there are WKB solutions of the form
$$\begin{array}{cc}\mathrm{\Psi }_\pm (a,\phi )\hfill & =[p(a)]^{1/2}\mathrm{exp}\left[\pm i_{a__H}^ap(a^{})𝑑a^{}i\pi /4\right]\hfill \\ & =\left[4\pi ^2\sqrt{3ϵ(6ϵ+1)}a\phi \sqrt{a^2V(\phi )3ϵ\phi ^2}\right]^{1/2}\mathrm{exp}\left[\pm i\frac{4\pi ^2\sqrt{3ϵ(6ϵ+1)}\phi }{3V(\phi )}\left[a^2V(\phi )3ϵ\phi ^2\right]^{3/2}i\pi /4\right],aa__H,\hfill \end{array}$$
(10)
where $`p(a)=[U(a)]^{1/2}`$. For the classically forbiden (exponential) region $`a<a__H`$, where the scale factor is small, there are WKB solutions of the form
$$\begin{array}{cc}\stackrel{~}{\mathrm{\Psi }}_\pm (a,\phi )\hfill & =|p(a)|^{1/2}\mathrm{exp}\left[\pm _a^{a__H}|p(a^{})|𝑑a^{}\right]\hfill \\ & =\left[4\pi ^2\sqrt{3ϵ(6ϵ+1)}a\phi \sqrt{(3ϵ\phi ^2a^2V(\phi ))}\right]^{1/2}\mathrm{exp}\left[\pm \frac{4\pi ^2\sqrt{3ϵ(6ϵ+1)}\phi }{3V(\phi )}\left[3ϵ\phi ^2a^2V(\phi )\right]^{3/2}\right],aa__H.\hfill \end{array}$$
(11)
We can impose the boundary condition in either the classically allowed region or the classically forbidden region, and then match the solutions in the two regions by the WKB standard matching procedure to specify the WKB wave function.
There are several comprehensive and well studied boundary condition proposals in the literature (for a recent review, see Ref. ). Now we seek for the specified WKB wave function with different boundary conditions following Vilenkin. Hartle & Hawking proposed that the specified wave function should be given by the Euclidean path integral, $`\mathrm{\Psi }_{HH}=e^{S_E}`$, which is taken over compact Euclidean geometries and matter fields with a specified field configuration at the boundary. Note that $`S_E`$ is the Euclidean action. The Hartle-Hawking wave function is specified by requiring that it is given by $`\mathrm{exp}(S_E)`$ in the Euclidean regime. This gives
$$\begin{array}{cc}\mathrm{\Psi }_{HH}(a<a__H,\phi )=\hfill & \stackrel{~}{\mathrm{\Psi }}_{}(a,\phi ),\hfill \\ \mathrm{\Psi }_{HH}(a>a__H,\phi )=\hfill & \mathrm{\Psi }_+(a,\phi )\mathrm{\Psi }_{}(a,\phi ).\hfill \end{array}$$
(12)
However, Linde argued that the wave function should be given by, $`\mathrm{\Psi }_L=e^{+S_E}`$, which requires a reverse sign of the exponential in the Euclidean regime. Together with the continuation to the classically allowed range of $`a`$, one can get the Linde wave function as
$$\begin{array}{cc}\mathrm{\Psi }_L(a<a__H,\phi )=\hfill & \stackrel{~}{\mathrm{\Psi }}_+(a,\phi ),\hfill \\ \mathrm{\Psi }_L(a>a__H,\phi )=\hfill & \frac{1}{2}[\mathrm{\Psi }_+(a,\phi )+\mathrm{\Psi }_{}(a,\phi )].\hfill \end{array}$$
(13)
In addition, Vilenkin suggested that the wave function of the universe should be specified either by the tunneling boundary condition or by a Lorentzian path integral, $`\mathrm{\Psi }_V=e^{iS}`$. Let us give a more precise statement of the Vilenkin’s tunneling boundary condition as follows: $`\mathrm{\Psi }_V`$ is the solution to the WDWE that is everywhere bounded and only consists of outgoing modes at singular boundaries of superspace. The superspace for our model is two-dimensional space with coordinates ($`a,\phi `$), where, $`0<a<\mathrm{},\mathrm{}<\phi <\mathrm{}`$. The unique non-singular part of the boundary is $`a=0`$ with $`\phi `$ being finite. The rest are singular and consist of configuration with one or both $`a`$ and $`\phi `$ being infinite. Note that as $`a`$ approaches zero, the coefficients of the second and third terms in Eq.8 blow up. As the boundary condition requires, we are to get a regular solution, it seems reasonable to neglect the second and third terms in Eq.8 to get the wave function of the universe. It is easy to check that $`\mathrm{\Psi }_{}(a,\phi )`$ and $`\mathrm{\Psi }_+(a,\phi )`$ describe an expanding and a contracting universe, respectively. Vilenkin’s boundary condition requires that only the expanding component $`\mathrm{\Psi }_{}(a,\phi )`$ should be present at large $`a`$. The wave function within the quantum barrier can be found from the WKB connection formula. Finally, we get the Vilenkin wave function,
$$\begin{array}{cc}\mathrm{\Psi }_V(a<a__H,\phi )=\hfill & \stackrel{~}{\mathrm{\Psi }}_+(a,\phi )\frac{i}{2}\stackrel{~}{\mathrm{\Psi }}_{}(a,\phi ),\hfill \\ \mathrm{\Psi }_V(a>a__H,\phi )=\hfill & \mathrm{\Psi }_{}(a,\phi ).\hfill \end{array}$$
(14)
This completes the calculation of the wave function of the universe under different types of boundary condition.
## IV Result and Discussion
Unfortunately, there has been no consensus so far on the interpretation of the wave function of the universe among the quantum cosmology community, except the statement that when $`\mathrm{\Psi }\mathrm{exp}[iS/\mathrm{}]`$, classical behavior should be recovered. In order to understand the physical meaning of the resulting wave functions, we employ the so-called “peak interpretation”, in which a prediction is said to be made when the wave function is sharply peaked in a certain region and almost zero elsewhere. It is worthwhile to point out that, the causal or the the Bohm-de Broglie interpretation will lead to the same result.
For the classical region, $`aa__H`$, where spacetime has the classical meaning and the classical solutions are valid, all three kinds of wave function are essentially oscillating, so the probability distribution does not depend sensitively on $`a`$ or $`\phi `$. This independence implies that there is approximatively equal probability for each point on the classical trajectories. Therefore, the properties of the wave function of the universe in the quantum era is crucial for the subsequent evolution of the universe. Within the quantum barrier, both Linde’s and Vilenkin’s wave functions are dominated by the decaying exponential $`\stackrel{~}{\mathrm{\Psi }}_+(a,\phi )`$, which piles up at $`V(\phi )=0`$, i.e., $`\phi ^2=\upsilon ^2`$, in the case of $`V(\phi )=\frac{\lambda }{8}(\phi ^2\upsilon ^2)^2`$. Since $`ϵ`$ is small, the distribution will be concentrated in the narrow region around $`\phi ^2=\upsilon ^2`$. Therefore the gravitational and cosmological “constants”, $`G_{\mathrm{ind}}`$, $`\mathrm{\Lambda }_{\mathrm{ind}}`$, acquire constant values, and especially, the latter is equal to zero. It implies that no matter how the cosmological constant can vary in the classical models, it will posses zero value when the universe comes from quantum creation. Nevertheless, Hartle-Hawking wave function within the barrier contain only the growing exponential $`\stackrel{~}{\mathrm{\Psi }}_{}(a,\phi )`$, which piles up at $`V(\phi )=V_{max}`$, where $`V_{max}`$ is the maximum value of $`V(\phi )`$. Hence the universe would prefer a large cosmological constant.
In summary, we have investigated a quantum cosmological model with the induced gravity theory. After the WDWE was constructed in the minisuperspace approximation, we have solved it using three kinds of boundary condition. We have shown that the amplitude of the resulting tunneling wave function sharply peaks around the classical trajectory only for a vanishing cosmological constant.
###### Acknowledgements.
We thank the anonymous refree for valuable comments, Prof. L. M. Krauss for helpful suggestions and Prof. T. Kiang of Dunsink Observatory, Ireland, for polishing up the English. This work was supported by the National Natural Science Foundation of China, under Grant No. 19903002.
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# First Observation of the Σ_𝑐^{∗+} Baryon and a New Measurement of the Σ_𝑐⁺ Mass
## Abstract
Using data recorded with the CLEO II and CLEO II.V detector configurations at the Cornell Electron Storage Ring, we report the first observation and mass measurement of the $`\mathrm{\Sigma }_c^+`$ charmed baryon, and an updated measurement of the mass of the $`\mathrm{\Sigma }_c^+`$ baryon. We find $`M(\mathrm{\Sigma }_c^+)M(\mathrm{\Lambda }_c^+)=(231.0\pm 1.1\pm 2.0)\mathrm{MeV}`$, and $`M(\mathrm{\Sigma }_c^+)M(\mathrm{\Lambda }_c^+)=(166.4\pm 0.2\pm 0.3)\mathrm{MeV}`$, where the errors are statistical and systematic respectively.
preprint: CLNS 00/1681 CLEO 00-13
R. Ammar,<sup>1</sup> A. Bean,<sup>1</sup> D. Besson,<sup>1</sup> R. Davis,<sup>1</sup> N. Kwak,<sup>1</sup> X. Zhao,<sup>1</sup> S. Anderson,<sup>2</sup> V. V. Frolov,<sup>2</sup> Y. Kubota,<sup>2</sup> S. J. Lee,<sup>2</sup> R. Mahapatra,<sup>2</sup> J. J. O’Neill,<sup>2</sup> R. Poling,<sup>2</sup> T. Riehle,<sup>2</sup> A. Smith,<sup>2</sup> C. J. Stepaniak,<sup>2</sup> J. Urheim,<sup>2</sup> S. Ahmed,<sup>3</sup> M. S. Alam,<sup>3</sup> S. B. Athar,<sup>3</sup> L. Jian,<sup>3</sup> L. Ling,<sup>3</sup> M. Saleem,<sup>3</sup> S. Timm,<sup>3</sup> F. Wappler,<sup>3</sup> A. Anastassov,<sup>4</sup> J. E. Duboscq,<sup>4</sup> E. Eckhart,<sup>4</sup> K. K. Gan,<sup>4</sup> C. Gwon,<sup>4</sup> T. Hart,<sup>4</sup> K. Honscheid,<sup>4</sup> D. Hufnagel,<sup>4</sup> H. Kagan,<sup>4</sup> R. Kass,<sup>4</sup> T. K. Pedlar,<sup>4</sup> H. Schwarthoff,<sup>4</sup> J. B. Thayer,<sup>4</sup> E. von Toerne,<sup>4</sup> M. M. Zoeller,<sup>4</sup> S. J. Richichi,<sup>5</sup> H. Severini,<sup>5</sup> P. Skubic,<sup>5</sup> A. Undrus,<sup>5</sup> S. Chen,<sup>6</sup> J. Fast,<sup>6</sup> J. W. Hinson,<sup>6</sup> J. Lee,<sup>6</sup> D. H. Miller,<sup>6</sup> E. I. Shibata,<sup>6</sup> I. P. J. Shipsey,<sup>6</sup> V. Pavlunin,<sup>6</sup> D. Cronin-Hennessy,<sup>7</sup> A.L. Lyon,<sup>7</sup> E. H. Thorndike,<sup>7</sup> C. P. Jessop,<sup>8</sup> H. Marsiske,<sup>8</sup> M. L. Perl,<sup>8</sup> V. Savinov,<sup>8</sup> X. Zhou,<sup>8</sup> T. E. Coan,<sup>9</sup> V. Fadeyev,<sup>9</sup> Y. Maravin,<sup>9</sup> I. Narsky,<sup>9</sup> R. Stroynowski,<sup>9</sup> J. Ye,<sup>9</sup> T. Wlodek,<sup>9</sup> M. Artuso,<sup>10</sup> R. Ayad,<sup>10</sup> C. Boulahouache,<sup>10</sup> K. Bukin,<sup>10</sup> E. Dambasuren,<sup>10</sup> S. Karamov,<sup>10</sup> G. Majumder,<sup>10</sup> G. C. Moneti,<sup>10</sup> R. Mountain,<sup>10</sup> S. Schuh,<sup>10</sup> T. Skwarnicki,<sup>10</sup> S. Stone,<sup>10</sup> G. Viehhauser,<sup>10</sup> J.C. Wang,<sup>10</sup> A. Wolf,<sup>10</sup> J. Wu,<sup>10</sup> S. Kopp,<sup>11</sup> A. H. Mahmood,<sup>12</sup> S. E. Csorna,<sup>13</sup> I. Danko,<sup>13</sup> K. W. McLean,<sup>13</sup> Sz. Márka,<sup>13</sup> Z. Xu,<sup>13</sup> R. Godang,<sup>14</sup> K. Kinoshita,<sup>14,</sup><sup>*</sup><sup>*</sup>*Permanent address: University of Cincinnati, Cincinnati, OH 45221 I. C. Lai,<sup>14</sup> S. Schrenk,<sup>14</sup> G. Bonvicini,<sup>15</sup> D. Cinabro,<sup>15</sup> S. McGee,<sup>15</sup> L. P. Perera,<sup>15</sup> G. J. Zhou,<sup>15</sup> E. Lipeles,<sup>16</sup> S. P. Pappas,<sup>16</sup> M. Schmidtler,<sup>16</sup> A. Shapiro,<sup>16</sup> W. M. Sun,<sup>16</sup> A. J. Weinstein,<sup>16</sup> F. Würthwein,<sup>16,</sup>Permanent address: Massachusetts Institute of Technology, Cambridge, MA 02139. D. E. Jaffe,<sup>17</sup> G. Masek,<sup>17</sup> H. P. Paar,<sup>17</sup> E. M. Potter,<sup>17</sup> S. Prell,<sup>17</sup> V. Sharma,<sup>17</sup> D. M. Asner,<sup>18</sup> A. Eppich,<sup>18</sup> T. S. Hill,<sup>18</sup> R. J. Morrison,<sup>18</sup> R. A. Briere,<sup>19</sup> G. P. Chen,<sup>19</sup> B. H. Behrens,<sup>20</sup> W. T. Ford,<sup>20</sup> A. Gritsan,<sup>20</sup> J. Roy,<sup>20</sup> J. G. Smith,<sup>20</sup> J. P. Alexander,<sup>21</sup> R. Baker,<sup>21</sup> C. Bebek,<sup>21</sup> B. E. Berger,<sup>21</sup> K. Berkelman,<sup>21</sup> F. Blanc,<sup>21</sup> V. Boisvert,<sup>21</sup> D. G. Cassel,<sup>21</sup> M. Dickson,<sup>21</sup> P. S. Drell,<sup>21</sup> K. M. Ecklund,<sup>21</sup> R. Ehrlich,<sup>21</sup> A. D. Foland,<sup>21</sup> P. Gaidarev,<sup>21</sup> R. S. Galik,<sup>21</sup> L. Gibbons,<sup>21</sup> B. Gittelman,<sup>21</sup> S. W. Gray,<sup>21</sup> D. L. Hartill,<sup>21</sup> B. K. Heltsley,<sup>21</sup> P. I. Hopman,<sup>21</sup> C. D. Jones,<sup>21</sup> D. L. Kreinick,<sup>21</sup> M. Lohner,<sup>21</sup> A. Magerkurth,<sup>21</sup> T. O. Meyer,<sup>21</sup> N. B. Mistry,<sup>21</sup> E. Nordberg,<sup>21</sup> J. R. Patterson,<sup>21</sup> D. Peterson,<sup>21</sup> D. Riley,<sup>21</sup> J. G. Thayer,<sup>21</sup> D. Urner,<sup>21</sup> B. Valant-Spaight,<sup>21</sup> A. Warburton,<sup>21</sup> P. Avery,<sup>22</sup> C. Prescott,<sup>22</sup> A. I. Rubiera,<sup>22</sup> J. Yelton,<sup>22</sup> J. Zheng,<sup>22</sup> G. Brandenburg,<sup>23</sup> A. Ershov,<sup>23</sup> Y. S. Gao,<sup>23</sup> D. Y.-J. Kim,<sup>23</sup> R. Wilson,<sup>23</sup> T. E. Browder,<sup>24</sup> Y. Li,<sup>24</sup> J. L. Rodriguez,<sup>24</sup> H. Yamamoto,<sup>24</sup> T. Bergfeld,<sup>25</sup> B. I. Eisenstein,<sup>25</sup> J. Ernst,<sup>25</sup> G. E. Gladding,<sup>25</sup> G. D. Gollin,<sup>25</sup> R. M. Hans,<sup>25</sup> E. Johnson,<sup>25</sup> I. Karliner,<sup>25</sup> M. A. Marsh,<sup>25</sup> M. Palmer,<sup>25</sup> C. Plager,<sup>25</sup> C. Sedlack,<sup>25</sup> M. Selen,<sup>25</sup> J. J. Thaler,<sup>25</sup> J. Williams,<sup>25</sup> K. W. Edwards,<sup>26</sup> R. Janicek,<sup>27</sup> P. M. Patel,<sup>27</sup> and A. J. Sadoff<sup>28</sup>
<sup>1</sup>University of Kansas, Lawrence, Kansas 66045
<sup>2</sup>University of Minnesota, Minneapolis, Minnesota 55455
<sup>3</sup>State University of New York at Albany, Albany, New York 12222
<sup>4</sup>Ohio State University, Columbus, Ohio 43210
<sup>5</sup>University of Oklahoma, Norman, Oklahoma 73019
<sup>6</sup>Purdue University, West Lafayette, Indiana 47907
<sup>7</sup>University of Rochester, Rochester, New York 14627
<sup>8</sup>Stanford Linear Accelerator Center, Stanford University, Stanford, California 94309
<sup>9</sup>Southern Methodist University, Dallas, Texas 75275
<sup>10</sup>Syracuse University, Syracuse, New York 13244
<sup>11</sup>University of Texas, Austin, TX 78712
<sup>12</sup>University of Texas - Pan American, Edinburg, TX 78539
<sup>13</sup>Vanderbilt University, Nashville, Tennessee 37235
<sup>14</sup>Virginia Polytechnic Institute and State University, Blacksburg, Virginia 24061
<sup>15</sup>Wayne State University, Detroit, Michigan 48202
<sup>16</sup>California Institute of Technology, Pasadena, California 91125
<sup>17</sup>University of California, San Diego, La Jolla, California 92093
<sup>18</sup>University of California, Santa Barbara, California 93106
<sup>19</sup>Carnegie Mellon University, Pittsburgh, Pennsylvania 15213
<sup>20</sup>University of Colorado, Boulder, Colorado 80309-0390
<sup>21</sup>Cornell University, Ithaca, New York 14853
<sup>22</sup>University of Florida, Gainesville, Florida 32611
<sup>23</sup>Harvard University, Cambridge, Massachusetts 02138
<sup>24</sup>University of Hawaii at Manoa, Honolulu, Hawaii 96822
<sup>25</sup>University of Illinois, Urbana-Champaign, Illinois 61801
<sup>26</sup>Carleton University, Ottawa, Ontario, Canada K1S 5B6
and the Institute of Particle Physics, Canada
<sup>27</sup>McGill University, Montréal, Québec, Canada H3A 2T8
and the Institute of Particle Physics, Canada
<sup>28</sup>Ithaca College, Ithaca, New York 14850
The $`\mathrm{\Sigma }_c`$ states consist of a charmed quark and two light ($`u`$ or $`d`$) quarks, in an isospin one configuration. The $`J^P`$=$`\frac{1}{2}^+`$ $`\mathrm{\Sigma }_c^0`$ and $`\mathrm{\Sigma }_c^{++}`$ have been observed for many years. Their isospin partner, the $`\mathrm{\Sigma }_c^+`$, is more difficult to detect as it decays to the $`\mathrm{\Lambda }_c^+`$ with the emission of a neutral, as opposed to charged, pion. Neutral pion detection is typically prone to higher backgrounds and poorer momentum resolution than charged pion detection. The $`\mathrm{\Sigma }_c^+`$ was reported in one event in 1980, and then in a peak of 111 events by the CLEO collaboration in 1993. This analysis updates the earlier CLEO measurement with a much larger data sample. This permits a more accurate comparison of the isospin splitting of the $`\mathrm{\Sigma }_c`$ states.
The $`J^P=\frac{3}{2}^+`$ $`\mathrm{\Sigma }_c^{}`$ states are more difficult to observe than the $`J^P=\frac{1}{2}^+`$ states because of the larger natural width, which leads to a poorer signal to noise ratio. The $`\mathrm{\Sigma }_c^{++}`$ and $`\mathrm{\Sigma }_c^0`$ have now been identified in $`\mathrm{\Lambda }_c^+\pi ^\pm `$ final states, and their masses and widths measured. This analysis shows the first observation of their isospin partner, the $`\mathrm{\Sigma }_c^+`$, observed by its decay to $`\mathrm{\Lambda }_c^+\pi ^0`$. This observation completes the spectroscopy of the seven L=0 $`\mathrm{\Lambda }_c`$ and $`\mathrm{\Sigma }_c`$ baryons predicted by the quark model.
The data presented here were taken with the CLEO II and CLEO II.V detector configurations operating at the Cornell Electron Storage Ring (CESR). The data sample used in this analysis corresponds to an integrated luminosity of 13.7 $`fb^1`$ taken on the $`\mathrm{{\rm Y}}(4S)`$ resonance and in the continuum at energies just below the $`\mathrm{{\rm Y}}(4S)`$. Of this data, 4.7 $`fb^1`$ was taken with the CLEO II configuration . We detected charged tracks with a cylindrical drift chamber system inside a 1.4T solenoidal magnet, and we detected photons using an electromagnetic calorimeter consisting of 7800 cesium iodide crystals. The remainder of the data was taken with the CLEO II.V configuration, which has upgraded charged particle measurement capabilities, but the same same cesium iodide array to observe photons.
In order to obtain large statistics we reconstructed the $`\mathrm{\Lambda }_c^+`$ baryons using 15 different decay modes <sup>*</sup><sup>*</sup>*Charge conjugate modes are implicit throughout.. Measurements of the branching fractions into these modes have previously been presented by the CLEO collaboration, and the general procedures for finding those decay modes can be found in those references. For this search and data set, the exact analysis used has been optimized for high efficiency and low background. Briefly, particle identification of $`p,K^{}`$, and $`\pi `$ candidates was performed using specific ionization measurements in the drift chamber, and, when present, time-of-flight measurements. Hyperons were found by detecting their decay points separated from the main event vertex.
We reduce the combinatorial background, which is highest for charmed baryon candidates with low momentum, by applying a cut on $`x_p`$, where $`x_p=p/p_{max}`$, $`p`$ is the momentum of the charmed baryon candidate, $`p_{max}=\sqrt{E_{beam}^2M^2},`$ and $`E_{beam}`$ is the beam energy, and $`M`$ is the reconstructed mass of the candidate. Using a cut of $`x_p>0.5`$ (charmed baryons produced from decays of $`B`$ mesons near the $`B\overline{B}`$ threshold are kinematically limited to $`x_p<0.4`$), we fit the invariant mass distributions for these modes to a sum of a Gaussian signal and a low-order polynomial background. Combinations within $`1.6\sigma `$ of the mass of the $`\mathrm{\Lambda }_c^+`$ in each decay mode are taken as $`\mathrm{\Lambda }_c^+`$ candidates, where the resolution of each decay mode is taken from a Monte Carlo simulation (for the CLEO II and CLEO II.V datasets separately). In this $`x_p`$ region, we find a total yield of $`\mathrm{\Lambda }_c^+`$ signal of $``$ 58,000 combinations, and a signal-to-background ratio $`1:1.2`$.
Photons were detected by their energy deposition in the crystal calorimeter. Each photon candidate was required to be well isolated from charged particles, and to have an energy profile consistent with being due to a single photon. To ensure good signal to noise ratio, the transition $`\pi ^0`$ candidates were made from the combination of two photons each from the central part of the detector ($`\theta <0.7)`$, which has the best energy resolution. The calculated invariant mass of the photon pair was required to be within 2.5 standard deviations of the known $`\pi ^0`$ mass, and the momentum of the $`\pi ^0`$ candidate was required to be greater than 150 MeV/c. This momentum cut was optimized to maximize the signal to noise ratio of a resonance in the expected $`\mathrm{\Sigma }_c^+`$ mass range using a Monte Carlo simulation. The $`\pi ^0`$ candidates were then kinematically fit to the $`\pi ^0`$ mass, a procedure that improves the mass resolution of the $`\mathrm{\Sigma }_c^+`$ by around twenty percent.
The $`\mathrm{\Lambda }_c^+`$ candidates were combined with each $`\pi ^0`$ candidate in the event and the mass difference $`M(\mathrm{\Lambda }_c^+\pi ^0)M(\mathrm{\Lambda }_c^+)`$ was calculated. Our requirement on the fractional momentum, $`x_p>0.6`$, is placed on the $`\mathrm{\Lambda }_c^+\pi ^0`$ combination, not on the $`\mathrm{\Lambda }_c^+`$ itself. Given the energetics of the decays to $`\mathrm{\Lambda }_c^+\pi ^0`$, such a criterion corresponds roughly to $`x_p>0.5`$ for the $`\mathrm{\Lambda }_c^+`$ daughters. The mass difference spectrum, shown in Figure 1, shows two clear peaks. The first, near 167 $`\mathrm{MeV}`$, is due to $`\mathrm{\Sigma }_c^+`$ decays. The second, near $`230\mathrm{MeV}`$, we identify as the $`\mathrm{\Sigma }_c^+`$. If we fit this distribution to the sum of a third-order Chebychev polynomial distribution and two Gaussian signals, we obtain a yield of $`661_{60}^{+63}`$ events and a width of $`\sigma =`$ ($`2.84_{0.28}^{+0.31})`$MeV for the $`\mathrm{\Sigma }_c^+`$, and a yield of $`(327_{73}^{+78})`$ events and $`\sigma `$ = $`(5.6\pm 1.4)`$MeV for the second peak. The widths of these Gaussian signals are greater than the detector resolution, calculated from a GEANT-based Monte Carlo simulation program, of 1.90 and 3.55 MeV, respectively, in the relevant mass regions, indicating the likelihood that the particles have non-negligible intrinsic widths. If we fit the distribution instead to a sum of two p-wave Breit-Wigner functions convoluted with Gaussian resolution functions, we obtain values of the intrinsic width, $`\mathrm{\Gamma }`$, of $`(3.1_{0.8}^{+0.9})`$ MeV, and $`(7_5^{+6})`$ MeV respectively, for which the errors are statistical only. The pole masses obtained from this fit are $`M(\mathrm{\Sigma }_c^+)M(\mathrm{\Lambda }_c^+)=(166.44\pm 0.24)`$ MeV and $`M(\mathrm{\Sigma }_c^+)M(\mathrm{\Lambda }_c^+)`$= $`(231.0\pm 1.1)`$ MeV, where again the quoted errors are from the statistical errors in the fit. It is this second fit, which has a $`\chi ^2`$ of 73.3 for 93 degrees of freedom, which is shown in Figure 1. If the $`\mathrm{\Sigma }_c^+`$ signal were not included in the fit, it would have a $`\chi ^2`$ of 123 for 96 degrees of freedom. To obtain an estimate of the relative cross sections for $`\mathrm{\Lambda }_c^+`$, $`\mathrm{\Sigma }_c^+`$ and $`\mathrm{\Sigma }_c^+`$ baryons, we find the yield each of the three states with an $`x_p`$ cut on each candidate of 0.6. After correcting for the efficiency of the transition $`\pi ^0`$, we find the ratio $`N(\mathrm{\Sigma }_c^+`$):$`N(\mathrm{\Lambda }_c^+`$)= $`0.116_{0.014}^{+0.016}`$$`\pm `$$`0.022`$ and $`N(\mathrm{\Sigma }_c^+)`$:$`N(\mathrm{\Lambda }_c^+)`$=$`0.043_{0.012}^{+0.016}`$$`\pm `$$`0.007`$, where the errors are statistical and systematic respectively. The systematic uncertainty includes the uncertainty in the $`\pi ^0`$ reconstruction efficiency and differences in the yield obtained with different signal shapes. We note that we are not calculating the production ratios of these states, as we are unable to measure their full momentum spectra.
We have considered many different possible sources of systematic uncertainty in the measurements of the masses and widths of these resonances. We have checked the consistency of the results obtained with each of the two detector configurations separately, as well as with a variety of different background and signal shapes, different criteria on the $`\pi ^0`$ momenta, and different $`\mathrm{\Lambda }_c^+`$ decay modes. We find the dominating systematic uncertainties in the mass measurement of the $`\mathrm{\Sigma }_c^+`$ to be due to signal shape (0.2 MeV) and the uncertainty in the $`\pi ^0`$ momentum measurement (0.2 MeV). These combine to give a total systematic uncertainty in the measurement of $`M(\mathrm{\Sigma }_c^+)`$ of 0.3 MeV. In the case of the $`\mathrm{\Sigma }_c^+`$, the mass measurement is sensitive to both the shape of the signal and also to the shape of the background function used, and we estimate a total systematic uncertainty of 2 MeV in the measurement of the pole mass. Although the intrinsic width measurement of the $`\mathrm{\Sigma }_c^+`$ is statistically nearly four standard deviations from 0, there should also be added a systematic uncertainty which we estimate to be 0.8 MeV, due mostly to uncertainties in the energy resolution of the transition pion. The combination of statistical and systematic uncertainties lead us to set an upper limit of 4.6 MeV (at the 90% confidence level) on $`\mathrm{\Gamma }(\mathrm{\Sigma }_c^+)`$. The width of the $`\mathrm{\Sigma }_c^+`$ is particularly sensitive to the parameterization of the background shape, and we estimate a systematic uncertainty of 5 MeV in the measurement of $`\mathrm{\Gamma }(\mathrm{\Sigma }_c^+)`$ mostly from this source. This, combined with the statistical error, leads to a 90% confidence level limit of $`\mathrm{\Gamma }<17`$ MeV.
Our result for the mass of the $`\mathrm{\Sigma }_c^+`$ is rather lower than the previous CLEO measurement which was based upon a small subset of these data, and lower than the measured masses of the $`\mathrm{\Sigma }_c^{++}`$ and $`\mathrm{\Sigma }_c^0`$, for which more experimental data is available. This is consistent with the theoretical expectations for this isospin splitting . The mass of the $`\mathrm{\Sigma }_c^+`$ is also lower than that of its isospin partners, but the experimental errors are too large for this splitting to be significant.
In conclusion, we have made a new measurement of the mass of the $`\mathrm{\Sigma }_c^+`$ and find $`M(\mathrm{\Sigma }_c^+)M(\mathrm{\Lambda }_c^+)`$ = $`(166.4\pm 0.2\pm 0.3)`$ MeV. We report the first observation of the $`\mathrm{\Sigma }_c^+`$ and find $`M(\mathrm{\Sigma }_c^+)M(\mathrm{\Lambda }_c^+)`$ = $`(231.0\pm 1.1\pm 2.0)`$ MeV. These measurements are consistent with expectations based upon the previously observed isospin partners of these two particles.
We gratefully acknowledge the effort of the CESR staff in providing us with excellent luminosity and running conditions. This work was supported by the National Science Foundation, the U.S. Department of Energy, the Research Corporation, the Natural Sciences and Engineering Research Council of Canada, the A.P. Sloan Foundation, the Swiss National Science Foundation, the Texas Advanced Research Program, and the Alexander von Humboldt Stiftung.
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# 1 Examples for types of fermionic two-loop diagrams contributing to muon decay.
## Acknowledgements
We thank S. Bauberger, M. Böhm, S. Heinemeyer and A. Stremplat for collaboration in early stages of this work. We also thank P. Gambino and D. Stöckinger for useful discussions and communications. We are grateful to M. Awramik and M. Czakon for detailed comparisons with their results , which helped to identify an error in our computation.
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# Domain Walls in SU(5)
## I Introduction
Topological defects can be produced at a symmetry breaking phase transition and would be long-lived relics of the symmetric phase. If topological defects were produced during a phase transition in the very early universe, they could survive until the present epoch and thus provide a window to the very early universe. The lack of observable defects in the present universe helps place strong constraints on particle physics model building and early universe cosmology.
A prototype symmetry breaking relevant for Grand Unified particle physics is
$$\mathrm{SU}(5)[\mathrm{SU}(3)\times \mathrm{SU}(2)\times \mathrm{U}(1)]/\mathrm{Z}_6.$$
The corresponding phase transition would produce magnetic monopoles. If the only factors affecting the evolution of the monopoles are the sub-luminal expansion of the universe and monopole-antimonopole Coulombic interactions, the monopole abundance grossly violates the observed absence of monopoles in the present universe. The monopole over-abundance problem is solved by invoking superluminal universal expansion (i.e. inflation ) or by extending the particle physics model so that the U(1) symmetry gets broken for a short duration leading to confining forces between monopoles and antimonopoles and thus enhancing their annihilation rate<sup>*</sup><sup>*</sup>*There is another possibility - that the Grand Unified phase transition never occurred and hence there never was a monopole over-abundance problem .. Recently we have investigated the possibility that the Grand Unified phase transition may also have produced a network of domain walls together with the magnetic monopoles. These walls would interact with the monopoles and sweep them away, reducing their abundance to an acceptably low level. It is to pursue this scenario in greater detail that we now study the structure of domain walls in the SU(5) model.
The bosonic sector of the SU(5) model is:
$$L=Tr(D_\mu \mathrm{\Phi })^2V(\mathrm{\Phi })$$
(1)
where $`\mathrm{\Phi }`$ is an adjoint of SU(5), $`D_\mu \mathrm{\Phi }=_\mu \mathrm{\Phi }ig[X_\mu ,\mathrm{\Phi }]`$ $`X_\mu `$ are the gauge fields, and the potential is given by:
$`V(\mathrm{\Phi })=m^2Tr(\mathrm{\Phi }^2)`$ $`+`$ $`h[Tr(\mathrm{\Phi }^2)]^2`$ (2)
$`+`$ $`\lambda Tr(\mathrm{\Phi }^4)+\gamma Tr(\mathrm{\Phi }^3)V_0,`$ (3)
where, $`V_0`$ is a constant that we will choose below. The SU(5) symmetry is broken to $`[SU(3)\times SU(2)\times U(1)]/Z_6`$ if the Higgs acquires a vacuum expectation value (VEV) equal to
$$\mathrm{\Phi }_0=\frac{\eta }{2\sqrt{15}}\mathrm{diag}(2,2,2,3,3),$$
(4)
where
$$\eta =\frac{m}{\sqrt{\lambda ^{}}},$$
(5)
$$\lambda ^{}h+\frac{7}{30}\lambda .$$
(6)
For the potential to have its global minimum at $`\mathrm{\Phi }=\mathrm{\Phi }_0`$, the parameters are constrained to satisfy:
$$\lambda 0,\lambda ^{}0.$$
(7)
For the global minimum to have $`V(\mathrm{\Phi }_0)=0`$, in eq. (3) we set
$$V_0=\frac{\lambda ^{}}{4}\eta ^4.$$
(8)
The model in eq. (1) does not have any topological domain walls because the vacua related by $`\mathrm{\Phi }\mathrm{\Phi }`$ are not degenerate. However if $`\gamma `$ is small, there are walls connecting the two kinds of vacua that are almost topological. In our analysis we will set $`\gamma =0`$, in which case the symmetry of the model is SU(5)$`\times `$Z<sub>2</sub> and an expectation of $`\mathrm{\Phi }`$ breaks the Z<sub>2</sub> symmetry leading to topological domain walls in addition to the magnetic monopoles arising from the SU(5) breaking.
In this paper we will study the domain walls present in the SU(5)$`\times `$Z<sub>2</sub> model. The simplest kind of domain wall is the kink that has been studied in a single scalar field model with Z$`{}_{2}{}^{}1`$ (eg. ). In we studied the interaction of the SU(5) kink with magnetic monopoles and found that the monopoles spread out along the kink on collision and never pass through. This confirmed the conjecture in Ref. that kinks could sweep away magnetic monopoles. However, the investigations of this paper show that the kink solution of the SU(5)$`\times `$Z<sub>2</sub> model is unstable to perturbations. The model contains another domain wall solution that is lighter than the kink and is perturbatively stable. The adjoint field does not vanish in the core of these new domain wall solutions and hence only a subgroup of the SU(5) is restored at the center. For this reason, the interactions of these domain walls with magnetic monopoles is expected to be much more complex (as compared to the kink), depending on the particular group orientation of the monopole relative to the wall.
We will begin our analysis in Sec. II by constructing the kink and performing the stability analysis. Then in Sec. III we will proceed to construct the domain wall in the model, prove that it is lighter than the kink, and show that it is perturbatively stable for a range of parameters. In Sec. IV we will consider the interaction of monopoles and domain walls and show that a monopole whose orientation in group space is aligned with a colliding domain wall, will pass through and not get swept away. We further conjecture that monopoles that are misaligned with the domain wall will be swept away but have not yet shown this. We draw an analogy of the sweeping out process with that of a polarization filter that “sweeps out” orthogonally polarized light and only lets through a certain polarization.
## II SU(5) kink: solution and stability
The kink solution is the Z<sub>2</sub> kink along the $`\mathrm{\Phi }_0`$ direction (see eq. (4)). Therefore:
$$\mathrm{\Phi }_k=\mathrm{tanh}(\sigma z)\mathrm{\Phi }_0$$
(9)
with $`\sigma m/\sqrt{2}`$ (see eq. (6)), and all the gauge fields vanish. It is straightforward to check that $`\mathrm{\Phi }_k`$ solves the equations of motion with the boundary conditions $`\mathrm{\Phi }(z=\pm \mathrm{})=\pm \mathrm{\Phi }_0`$.
As is well-known , the mass (per unit area) of the kink is:
$$M_k=\frac{2\sqrt{2}}{3}\frac{m^3}{\lambda ^{}}.$$
(10)
Here we will examine the stability of the kink under general perturbations. So we write:
$$\mathrm{\Phi }=\mathrm{\Phi }_k+\mathrm{\Psi }$$
(11)
Since the kink solution is invariant under translations and rotations in the $`xy`$plane, it is easy to show that the perturbations that might cause an instability arise from perturbations of the scalar field and can only depend on $`z`$. Therefore we may set the gauge fields to zero and take $`\mathrm{\Psi }=\mathrm{\Psi }(t,z)`$.
The Z<sub>2</sub> kink is stable and hence we can restrict the scalar perturbations to be orthogonal to $`\mathrm{\Phi }_k`$. Furthermore, since the stability of the kink to diagonal perturbations has already been studied in Ref. , we only have to consider perturbations that cannot be diagonalized by a global SU(3)$`\times `$SU(2)$`\times `$U(1) transformation that leaves the kink invariant. Therefore we can write:
$$\mathrm{\Psi }=\underset{i=1}{\overset{12}{}}\psi ^iT^i,$$
(12)
where $`T^i`$ are all generators of SU(5) that do not commute with $`\mathrm{\Phi }_0`$.
Next we analyze the linearized Schrodinger equation for small excitations $`\psi ^i=\psi _0^i(z)exp(i\omega t)`$ in the background of the kink:
$$[_z^2m^2+\varphi _k^2(z)(h+\lambda r_i)]\psi _0^i=\omega _i^2\psi _0^i,$$
(13)
where $`\varphi _k\mathrm{tanh}(\sigma z)`$ and $`r_i=7/30`$. Since this equation is identical for excitations along any of the $`12`$ directions, we can drop the index $`i`$. The kink is unstable if there is a solution to eq. (13) with a negative $`\omega ^2`$. Substituting eq. (9) into eq. (13) yields:
$$[_z^2+m^2(\mathrm{tanh}^2(\sigma z)1)]\psi _0=\omega ^2\psi _0.$$
(14)
This equation has a bound state solution $`\psi _0`$ sech$`(\sigma z)`$ with the eigenvalue $`\omega ^2=m^2/2`$. Since this result is independent of the parameters in the potential, we conclude that the kink in SU(5) is always unstable.
## III Domain wall
The domain wall solution is obtained if we choose the gauge fields to vanish at infinity and the scalar field to satisfy the boundary conditions:
$`\mathrm{\Phi }(z=\mathrm{})=\mathrm{\Phi }^{}{\displaystyle \frac{\eta }{2\sqrt{15}}}\mathrm{diag}(2,3,2,2,3)`$ (15)
$`=\eta \sqrt{{\displaystyle \frac{5}{12}}}(\lambda _3+\tau _3)+{\displaystyle \frac{\eta }{6}}(Y\sqrt{5}\lambda _8)`$ (16)
and
$`\mathrm{\Phi }(z=+\mathrm{})=\mathrm{\Phi }^+{\displaystyle \frac{\eta }{2\sqrt{15}}}\mathrm{diag}(3,2,2,3,2)`$ (17)
$`=\eta \sqrt{{\displaystyle \frac{5}{12}}}(\lambda _3+\tau _3){\displaystyle \frac{\eta }{6}}(Y\sqrt{5}\lambda _8).`$ (18)
Here $`\lambda _3`$, $`\lambda _8`$, $`\tau _3`$ and $`Y`$ are the diagonal generators of SU(5):
$$\lambda _3=\frac{1}{2}\mathrm{diag}(1,1,0,0,0),$$
(19)
$$\lambda _8=\frac{1}{2\sqrt{3}}\mathrm{diag}(1,1,2,0,0),$$
(20)
$$\tau _3=\frac{1}{2}\mathrm{diag}(0,0,0,1,1),$$
(21)
$$Y=\frac{1}{2\sqrt{15}}\mathrm{diag}(2,2,2,3,3).$$
(22)
Note that a local SU(5) transformations can be used to rotate $`\mathrm{\Phi }^+`$ into $`\mathrm{\Phi }^{}`$ so that the boundary conditions are like those of the kink with $`\mathrm{\Phi }(z=+\mathrm{})=\mathrm{\Phi }(z=\mathrm{})`$. However, then the solution for the domain wall will not be diagonal at all $`z`$. We prefer to use the above boundary conditions so that the solution is diagonal throughout.
The domain wall solution can be written as
$$\mathrm{\Phi }_{DW}(z)=a(z)\lambda _3+b(z)\lambda _8+c(z)\tau _3+d(z)Y.$$
(23)
The functions $`a`$, $`b`$, $`c`$, and $`d`$ must satisfy the static equations of motion:
$`a^{\prime \prime }=[m^2+(h+{\displaystyle \frac{2\lambda }{5}})d^2`$ $`+`$ $`(h+{\displaystyle \frac{\lambda }{2}})(a^2+b^2)+hc^2]a`$ (24)
$`+`$ $`{\displaystyle \frac{2\lambda abd}{\sqrt{5}}}`$ (25)
$`b^{\prime \prime }=[m^2+(h+{\displaystyle \frac{2\lambda }{5}})d^2`$ $`+`$ $`(h+{\displaystyle \frac{\lambda }{2}})(a^2+b^2)+hc^2]b`$ (26)
$`+`$ $`{\displaystyle \frac{\lambda d}{\sqrt{5}}}(a^2b^2)`$ (27)
$`c^{\prime \prime }=[m^2+(h+{\displaystyle \frac{9\lambda }{10}})d^2+(h`$ $`+`$ $`{\displaystyle \frac{\lambda }{2}})c^2`$ (28)
$`+`$ $`h(a^2+b^2)]c`$ (29)
$`d^{\prime \prime }=[m^2`$ $`+`$ $`(h+{\displaystyle \frac{7\lambda }{30}})d^2+(h+{\displaystyle \frac{2\lambda }{5}})(a^2+b^2)`$ (30)
$`+`$ $`(h+{\displaystyle \frac{9\lambda }{10}})c^2]d+{\displaystyle \frac{\lambda b}{\sqrt{5}}}(a^2{\displaystyle \frac{b^2}{3}}),`$ (31)
where primes refer to derivatives with respect to $`z`$. For reference, the kink solution (eq. (9)) corresponds to $`a(z)=0=b(z)=c(z)`$ and $`d(z)=\eta \mathrm{tanh}(\sigma z)`$.
The equations of motion for $`b`$ and $`c`$ and can be solved quite easily:
$$b(z)=\sqrt{5}d(z),c(z)=a(z).$$
(32)
This is consistent with the boundary conditions in eqs. (16) and (18). In addition, we require
$$a(z=\pm \mathrm{})=+\eta \sqrt{\frac{5}{12}},d(z=\pm \mathrm{})=\frac{\eta }{6}.$$
(33)
Then the remaining equations we need to solve are:
$$a^{\prime \prime }=\left[m^2+\left(6h+\frac{9}{10}\lambda \right)d^2+\left(2h+\frac{\lambda }{2}\right)a^2\right]a$$
(34)
$$d^{\prime \prime }=\left[m^2+\left(6h+\frac{39}{10}\lambda \right)d^2+\left(2h+\frac{3\lambda }{10}\right)a^2\right]d$$
(35)
These equations can be written in a cleaner form by rescaling:
$$A(z)=\sqrt{\frac{12}{5}}\frac{a}{\eta },D(z)=6\frac{d}{\eta },Z=mz.$$
(36)
Then
$$A^{\prime \prime }=\left[1+\frac{(1p)}{5}D^2+\frac{(4+p)}{5}A^2\right]A$$
(37)
$$D^{\prime \prime }=\left[1+pD^2+(1p)A^2\right]D$$
(38)
where primes on $`A`$ and $`D`$ denote differentiation with respect to $`Z`$, and
$$p=\frac{1}{6}\left[1+\frac{5\lambda }{12\lambda ^{}}\right].$$
(39)
Note that $`p[1/6,\mathrm{})`$ because of the constraints in eq. (7). The boundary conditions now are:
$$A(z=\pm \mathrm{})=+1,D(z=\pm \mathrm{})=\pm 1.$$
(40)
This system of equations has been solved by numerical relaxation and a sample solution is shown in Fig. 1. To find an approximate analytical solution, assume that $`|A^{\prime \prime }/A|<<1`$ is small everywhere. This assumption will be true for a certain range of the parameter $`p`$ which we can later determine. Then the square bracket on the right-hand side of eq. (37) is very small. This gives:
$$A\left[\frac{5}{4+p}\left\{1\frac{(1p)}{5}D^2\right\}\right]^{1/2}$$
(41)
We insert this expression for $`A`$ in eq. (38) and obtain the kink-type differential equation:
$$D^{\prime \prime }=q[1+D^2]D,$$
(42)
where
$$q=\frac{6p1}{p+4}=\frac{6\lambda }{\lambda +60\lambda ^{}}$$
(43)
and the solution is:
$$D(Z)\mathrm{tanh}\left(\sqrt{\frac{q}{2}}Z\right)$$
(44)
The parameter $`q`$ lies in the interval $`[0,6]`$. For $`q=1`$ (i.e. $`p=1`$) it is easy to check that this analytical solution is exact.
We can now check that our assumption $`|A^{\prime \prime }/A|<<1`$ is self-consistent provided $`p`$ is not much larger than a few.
The energy density for the fields $`A`$ and $`D`$ can be found from the Lagrangian in eq. (1) together with the ansatz in eq. (23), the solution for $`b`$ and $`c`$ in eq. (32) and the rescalings in eq. (36). The resulting expression for the energy per unit area of the domain wall is:
$$M_{DW}=\frac{m^3}{12\lambda ^{}}𝑑Z[5A_{}^{}{}_{}{}^{2}+D_{}^{}{}_{}{}^{2}+V(A,D)]$$
(45)
where,
$`V(A,D)=5A^2`$ $``$ $`D^2+{\displaystyle \frac{(p+4)}{2}}A^4`$ (46)
$`+`$ $`{\displaystyle \frac{p}{2}}D^4+(1p)A^2D^2+3.`$ (47)
The energy can be found numerically. However, here we will find an approximate analytic result. We can insert the approximate solution given above in eq. (45) but this leads to an expression that is not transparent. Instead it is more useful to consider another approximation for $`A`$ and $`D`$:
$$A1,D\mathrm{tanh}\left(\sqrt{\frac{p}{2}}Z\right).$$
(48)
(This approximation is exact for $`p=1`$.) A straightforward evaluation then gives:
$$M_{DWapprox}=M_k\frac{\sqrt{p}}{6}$$
(49)
where, $`M_k`$ is given in eq. (10).
We can now compare the domain wall energy to the kink energy. If the domain wall is the least energy solution for the given boundary conditions, the energy of the exact solution for the domain wall will be bounded above by the energy of the approximate solution. Note that this will be true even if the approximation used to find the analytical solution is not good. Hence this simple argument shows that the domain wall is lighter than the kink for $`p<36`$, or for $`h/\lambda >6.94/30`$. A numerical analysis shows that the domain wall is lighter than the kink even in the range $`6.94/30h/\lambda >7/30`$. Therefore the domain wall is lighter than the kink for the full range of parameters specified in eq. (7).
Next we study the stability of the domain wall solution. It is easy to show that the solution is stable to diagonal perturbations, so here we focus on off-diagonal perturbations. We write:
$$\mathrm{\Phi }=\mathrm{\Phi }_{DW}(z)+\underset{a=1}{\overset{20}{}}\psi ^a(z)N^a,$$
(50)
where $`N^a`$ are the non-diagonal generators of $`SU(5)`$ and $`\psi ^a`$ are small perturbations satisfying the boundary conditions $`\psi ^a(\pm \mathrm{})=0`$. Let us first consider the contribution to the energy density due to fields $`\psi ^a`$. To second order in perturbations the contributions from different modes, labeled by index $`a`$, do not couple. A more detailed analysis shows that the mode corresponding to
$$N_1\frac{1}{2}\left(\begin{array}{ccccc}0& 1& 0& 0& 0\\ 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0\end{array}\right)$$
(51)
is one of the $`8`$ modes that are most unstable. Let $`\psi `$ be any of these 8 fields. The contribution to the energy density due to $`\psi `$ is:
$`E_\psi ={\displaystyle \frac{1}{2}}(\psi ^{})^2{\displaystyle \frac{m^2}{2}}\psi ^2+{\displaystyle \frac{h}{4}}(\psi ^2+2a^2+6d^2)^2`$ (52)
$`+{\displaystyle \frac{\lambda }{4}}\psi ^2(a^2+{\displaystyle \frac{9}{5}}d^2)+\mathrm{higher}\mathrm{order}\mathrm{terms},`$ (53)
where $`a`$ and $`d`$ are defined by eq. (23). As in the case of the kink, we are interested in the linearized Schrodinger equation for small excitations $`\psi =\psi _0(z)exp(i\omega t)`$ in the background of the diagonal domain wall solution. Eq. (53) leads to
$$[_z^2m^2+(6h+\frac{9}{10}\lambda )d^2+(2h+\frac{\lambda }{2})a^2]\psi _0=\omega ^2\psi _0.$$
(54)
Comparing this with eq. (34) allows us to write:
$$\psi _0^{\prime \prime }+\frac{a^{\prime \prime }}{a}\psi _0=\omega ^2\psi _0.$$
(55)
If $`a^{\prime \prime }/a=0`$, as happens when $`p=1`$, then there are no non-trivial solutions to eq. (55) that satisfy the correct boundary conditions. Therefore, the diagonal domain wall solution is stable for at least one choice of parameters in the potential, namely, for $`p=1`$. By continuity there is a range of parameters around $`p=1`$ for which the domain wall is perturbatively stable.
## IV Interaction with monopoles and discussion
To understand the interaction of the domain wall with magnetic monopoles, it is first useful to understand the structure of the domain wall core. Since $`a(z)`$ is non-zero inside the domain wall, $`\mathrm{\Phi }_{DW}(z=0)=a(0)(\lambda _3+\tau _3)\mathrm{diag}(1,1,0,1,1)`$. Therefore, the symmetry inside the core is K$``$SU(2)$`\times `$SU(2)$`\times `$U(1)$`\times `$U(1). The first SU(2) factor arises due to rotations in the 2$`\times `$2 block with the entries equal to 1 in $`\mathrm{\Phi }_{DW}`$ (first and fourth rows and columns). The second SU(2) factor is due to the block with entries equal to -1 (second and fifth rows and columns). The two U(1) factors arise since there are two diagonal generators of SU(5) aside from those already accounted for in the two SU(2) factors, that commute with $`\mathrm{\Phi }_{DW}(z=0)`$. (We are ignoring any discrete factors that might be present.) The symmetry group K within the domain wall is to be contrasted with the full SU(5) symmetry which is restored within the kink. The fact that only a subgroup of the SU(5) symmetry is restored in the core of the wall means that the interaction of the monopole will now depend on the particular embedding of the monopole in SU(5) and its orientation in internal space relative to the domain wall.
Consider a magnetic monopole whose winding lies in the fourth and fifth rows and columns of $`\mathrm{\Phi }`$. Staying close to the notation of we write the scalar field of such a monopole as:
$$\mathrm{\Phi }_M(r)=P(r)\underset{a=1}{\overset{3}{}}x^a\tau _a+M(r)\lambda _8^{}+N(r)Y,$$
(56)
where $`\{\tau _a\}`$ are the SU(2) generators (see eq. (21)) and
$$\lambda _8^{}\frac{1}{2\sqrt{3}}\mathrm{diag}(1,2,1,0,0)=\frac{\sqrt{3}}{2}\lambda _3\frac{1}{2}\lambda _8.$$
The non-zero gauge fields are:
$`X_i={\displaystyle \underset{a=1}{\overset{3}{}}}X_i^a\tau _a,`$ (57)
$`X_i^a=ϵ_{ij}^a{\displaystyle \frac{x^j}{er^2}}(1K(r)).`$ (58)
The monopole profile functions, $`P(r)`$, $`M(r)`$, $`N(r)`$ and $`K(r)`$, are solutions of the static equations of motion with boundary conditions:
$`P(\mathrm{})=\eta \sqrt{{\displaystyle \frac{5}{12}}},M(\mathrm{})=\eta {\displaystyle \frac{\sqrt{5}}{3}},`$ (59)
$`N(\mathrm{})={\displaystyle \frac{\eta }{6}},K(\mathrm{})=1.`$ (60)
When the monopole and the wall are very far from each other, the combined field configuration can be described by the following ansatze:
$`\mathrm{\Phi }_{M+DW}=P(r){\displaystyle \frac{c(z^{})}{c(\mathrm{})}}{\displaystyle \underset{a=1}{\overset{3}{}}}x^a\tau _a+N(r){\displaystyle \frac{d(z^{})}{d(\mathrm{})}}Y`$ (61)
$`+M(r)\left[{\displaystyle \frac{\sqrt{3}}{2}}{\displaystyle \frac{a(z^{})}{a(\mathrm{})}}\lambda _3{\displaystyle \frac{1}{2}}{\displaystyle \frac{b(z^{})}{b(\mathrm{})}}\lambda _8\right],`$ (62)
where $`z^{}=zz_0`$ and $`z_0`$ is the initial position of the domain wall. When $`r`$ is small $`\mathrm{\Phi }_{M+DW}\mathrm{\Phi }_M`$ (eq. (56)) and when $`z^{}`$ is small $`\mathrm{\Phi }_{M+DW}\mathrm{\Phi }_{DW}`$ (eq. (23)) along the $`z`$-direction. The gauge fields are the same as for the monopole alone. We have purposely chosen the embedding of the monopole so that all interesting dynamics of the monopole-wall interaction is restricted to the fourth and fifth rows and columns of $`\mathrm{\Phi }_{M+DW}`$. This follows from the equations of motion and the commutation properties of the generators appearing in the ansatze (62). Let us then only consider the relevant part of the $`\mathrm{\Phi }_{M+DW}`$ matrix:
$`\mathrm{\Phi }_{2\times 2}{\displaystyle \frac{1}{2}}P(r){\displaystyle \frac{a(z^{})}{a(\mathrm{})}}\left(\begin{array}{cc}z& xiy\\ x+iy& z\end{array}\right)`$ (63)
$`{\displaystyle \frac{3}{2\sqrt{15}}}N(r){\displaystyle \frac{d(z^{})}{d(\mathrm{})}}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).`$ (64)
The form of $`\mathrm{\Phi }_{2\times 2}`$ suggests that the only field that is going to be considerably affected by the domain wall is $`N(r)`$ because $`a`$ is roughly constant in space. There is no angular dependence in the term with $`N(r)`$ in eq. (56) and hence $`N(r)`$ does not contribute to the winding of the monopole. Therefore, we do not expect the wall to affect the winding. Essentially the reason is that the SU(2) subgroup in which the monopole winding is located is not restored on the wall. We have checked that the monopole passes right through the wall explicitly in this case by numerically colliding the monopole and the wall.
If the magnetic monopole winding lies in the first and fourth block, it will experience a region of restored SU(2) symmetry inside the domain wall and hence we conjecture that such monopoles will unwind on the wall. If our conjecture is correct, the domain walls behave similarly to optical polarization filters, allowing monopoles with certain internal space polarizations to pass through and annihilating other polarizations. The detailed analysis of all possible monopole embeddings is a challenging project, both numerically and analytically, since one cannot avoid dealing with a large number of the fields present in SU(5).
There is another possibility that is worth pointing out. If a domain wall and a magnetic monopole are misaligned in internal space, it may not be possible to superpose the two solutions so as to get a monopole and a domain wall together. (Such a situation is known to occur when attempting to construct multimonopole or multistring solutions.) Then it is likely that there will be a long range force between the domain wall and misaligned monopole that will bring them together. On coming together the monopole could get annihilated on the wall, or else in some cases, it may get aligned and then pass through the wall.
Our considerations point to a very complicated aftermath of the GUT phase transition. Domain walls and magnetic monopoles would both be produced and would start interacting. The outcome of an interaction would depend on the internal space orientations of the monopole relative to the domain wall. Any given domain wall would be transparent to some monopoles but not to others. The relaxation of the system would depend on whether a monopole encounters a sufficient number of randomly oriented (in internal space) domain walls, at least one of which might sweep it away. It remains to be seen if domain walls can provide a means to solve the cosmological monopole over-abundance problem.
###### Acknowledgements.
We would like to thank Mark Trodden for many useful discussion. This work was supported by the DOE.
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# PDK-752TUHEP-00-0517 July 2000 New results on atmospheric neutrinos from Soudan 2 Presented at NU2000, the XIXth Int. Conference on Neutrino Physics and Astrophysics, June 16 - 21, 2000, Sudbury, Canada.
## 1 DETECTOR; DATA EXPOSURE
The Soudan-2 experiment is currently taking data using its fine-grained iron tracking calorimeter of total mass 963 tons. This detector images non-relativistic as well as relativistic charged particles produced in atmospheric neutrino reactions. It is operating underground at a depth of 2100 meters-water-equivalent on level 27 of the Soudan Mine State Park in northern Minnesota (northwest of Sudbury). The calorimeter’s modular design enabled data-taking to commence in April 1989 when the detector was one quarter of its full size; assembly of the detector was completed during 1993. Data-taking has continued with 85% live time, even though dynamite blasting has been underway nearby for the MINOS cavern excavation since Summer 1999. The data exposure as of this Conference is 5.40 fiducial kiloton-years (kty). Results presented here are based upon a 5.1 kty exposure.
The tracking calorimeter operates as a slow-drift (0.6 cm/$`\mu `$s) time projection chamber. Its tracking elements are meter-long plastic drift tubes which are placed into the corrugations of steel sheets. The sheets are stacked to form a tracking lattice of honeycomb geometry. A stack is packaged as a calorimeter module and the detector is assembled building-block fashion using these modules . The calorimeter is surrounded on all sides by a cavern-liner active shield array of two or three layers of proportional tubes .
Topologies for contained events in Soudan 2 include single track and single shower events (mostly $`\nu _\mu `$ and $`\nu _e`$ quasi-elastic reactions) and multiprong events. Flavor-tagging proceeds straightforwardly: An event having a leading, non-scattering track with ionization $`dE/dx`$ compatible with muon mass is a candidate charged current (CC) event of $`\nu _\mu `$ flavor; an event having a prompt, relatively energetic shower prong is a candidate $`\nu _e`$ CC event. Recoil protons of momenta greater than approx. 350 MeV/$`c`$ are imaged by the calorimeter, allowing a much-improved measurement of the incident neutrino direction, especially for sub-GeV quasi-elastic reactions.
## 2 ATMOSPHERIC $`\nu `$ FLAVOR RATIO
We measure the atmospheric neutrino $`\nu _\mu `$/$`\nu _e`$ flavor ratio-of-ratios $`R`$ using single track and single shower events which are fully contained within the calorimeter (all hits more than 20 cm from the nearest surface). These samples contain mostly quasi-elastic neutrino reactions, but include a background of photon and neutron reactions originating in cosmic ray muon interactions in the surrounding cavern rock. The latter “rock events” are mostly tagged by coincident hits in the active shield, however some are unaccompanied by shield hits and constitute a background. The amount of zero-shield-hit rock background in a neutrino event sample is estimated by fitting event vertex-depth distributions to a combination of tagged-rock and $`\nu `$ Monte Carlo distributions. As expected, the fits show the background to be mostly confined to outer regions of the calorimeter. Details of our analysis procedures for quasi-elastic events can be found in Refs. .
The track and shower event samples for our 5.1 kty exposure are summarized in Table 1. Our full detector Monte Carlo simulation of atmospheric neutrino interactions is based on the 1996 Bartol flux for the Soudan site .
After correction for cosmic ray muon induced background, the number of single track events observed in data is less than the number of single shower events, whereas the null oscillation Monte Carlo predicts the relative rates to be other-way-round. Consequently the flavor ratio-of-ratios obtained is less than 1.0 and is anomalous:
$$R=0.68\pm 0.11(stat)\pm 0.06(sys).$$
This value is equal to the $`R`$ value obtained last summer using 4.6 kty exposure .
## 3 SAMPLE FOR $`L/E`$ MEASUREMENT
The phenomenology for $`\nu _\mu `$ to $`\nu _\tau `$ oscillations is quite specific; neutrinos of muon flavor can metamorphose and thereby “disappear” according to the equation
$`P(\nu _\mu \nu _\tau )=`$ (1)
$`\mathrm{sin}^2(2\theta )\mathrm{sin}^2\left[{\displaystyle \frac{1.27\text{ }\mathrm{\Delta }m^2[\text{eV}^2]L\text{[km] }}{E_\nu \text{[GeV] }}}\right]`$
Consequently it is optimal to analyze for neutrino oscillations using the variable $`L/E_\nu `$. With the Soudan-2 calorimeter, measurement of event energy for charged current reactions is straightforward; we do this with resolution $`\mathrm{\Delta }E/E`$ which is 20% for $`\nu _\mu `$ CC’s and 23% for $`\nu _e`$ CC’s. To determine the neutrino path length $`L`$ however, the zenith angle $`\theta _z`$ of the incident neutrino must be reconstructed with accuracy. The path length can be calculated from the zenith angle according to
$`L(\theta _z)=`$ (2)
$`\sqrt{(Rd)^2\mathrm{cos}^2\theta _z+(d+h)(2Rd+h)}`$
$`(Rd)\mathrm{cos}\theta _z`$
where $`R`$ is the Earth’s radius, $`d`$ is the depth of the detector, and $`h`$ is the mean neutrino production height. The latter is a function of $`\nu `$ flavor, $`\nu `$ energy, and $`\theta _z`$ .
We select from our data an event sample suited to this measurement. We use a quasi-elastic track or shower event provided that the recoil proton is measured and that $`P_{lept}`$ exceeds 150 MeV/$`c`$; otherwise, if the recoil nucleon is not visible, we require the single lepton to have $`E_{vis}`$ great than 600 MeV. We also select multiprong events, provided they are energetic ($`E_{vis}`$ greater than 700 MeV) and have vector sum of $`P_{vis}`$ exceeding 450 MeV/$`c`$ (to ensure clear directionality). Additionally, the final state lepton momenta are required to exceed 250 MeV/$`c`$. For the selected sample, flavor tagging is estimated to be correct for more than 92% of events. The resolution for recovering the incident neutrino direction is evaluated using the mean angular separation between “true” versus reconstructed neutrino direction in Monte Carlo events. The mean separations are $`33.2^{}`$ for $`\nu _\mu `$ CC’s and $`21.3^{}`$ for $`\nu _e`$ CC’s . The resolution in $`\mathrm{log}`$ $`L/E_\nu `$ ($`L`$ in kilometers, $`E_\nu `$ in GeV) is better than 0.5 for the selected sample. Hereafter we refer to these events as “HiRes events”.
The zero-shield-hit rock background, as estimated by the fits to event vertex depth distributions, comprises 6.8% (5.1%) of the $`\nu _\mu `$ ($`\nu _e`$) flavor sample of HiRes events.
Table 2 shows the HiRes event populations. After background subtraction there are 106.3 data events of $`\nu _\mu `$ flavor and 132.8 events of $`\nu _e`$-flavor. Using these events, whose mean energy is higher than that of our track and shower events, the ratio-of-ratios is $`R=0.67\pm 0.12`$, which is also significantly less than 1.0.
The atmospheric Monte Carlo (MC) sample represents 28.2 kiloton years of exposure. The MC event rates displayed in Table 2 have been normalized to the $`\nu _e`$ data sample. This normalization is equivalent to a reduction of the Bartol neutrino fluxes by 21%. The assumption implicit with this adjustment is that the $`\nu _e`$ sample is devoid of oscillation effects. Figs. 12 and 3 show HiRes distributions with this normalization in place.
Fig. 1 shows the distributions of these samples in cosine of the zenith angle. For $`\nu _e`$ events, the shape of the distribution for data (Fig. 1a, crosses) coincides with that predicted by the Monte Carlo (dashed histogram) for null oscillations. The distribution of $`\nu _\mu `$ data however, falls below the MC prediction in all bins (Fig. 1b) with the relative dearth being more pronounced for $`\nu _\mu `$’s incident from below horizon. Distributions in $`\mathrm{log}\left(L/E_\nu \right)`$ for HiRes events are shown in Fig. 2. For null oscillations this variable distributes according to a ‘phase space’ which reflects the neutrino points-of-origin throughout the spherical shell volume of the Earth’s atmosphere. That is, down-going $`\nu `$’s populate the peak at lower $`\mathrm{log}\left(L/E_\nu \right)`$ from 0.0 to 2.0. Neutrinos incident from/near the horizon occur within the dip region extending from 2.0 to 2.6, while upward-going neutrinos populate the peak at higher values. Fig. 2a shows that, allowing for statistical fluctuations, the $`\nu _e`$ data follows the shape of the null oscillation MC distribution. In contrast, the $`\nu _\mu `$ data (Fig. 2b) falls below the null ocillation MC for all but the most vertically down-going flux.
## 4 ($`\mathrm{sin}^22\theta `$, $`\mathrm{\Delta }m^2`$) ALLOWED REGION
To convert results of our atmospheric neutrino simulation generated under the no-oscillation hypothesis into simulated neutrino oscillation data, we apply to every MC event an $`L/E_\nu `$-dependent weight representing the probability of $`\nu _\mu `$ flavor survival for a given $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^22\theta `$.
An exploratory matchup of $`\nu _\mu `$ data with trial oscillation scenarios is shown in Fig. 3. For the mixing angle $`\mathrm{sin}^22\theta `$ fixed at unity, we plot the MC distribution weighted for $`\nu _\mu `$ to $`\nu _\tau `$ oscillations with differing $`\mathrm{\Delta }m^2`$ settings. With $`\mathrm{\Delta }m^2=10^4\text{eV}^2`$ (Fig. 3a), the oscillation prediction lies above the data in most $`L/E_\nu `$ bins. With $`\mathrm{\Delta }m^2=10^3\text{eV}^2`$ (Fig. 3b), the upgoing $`\nu `$ flux is now better described by the MC, although the expectation for horizontal and down-going neutrinos remains somewhat high. With $`\mathrm{\Delta }m^2=7\times 10^3\text{eV}^2`$, a rough agreement overall is achieved (Fig. 3c). Going to higher $`\mathrm{\Delta }m^2`$, we find that for $`\mathrm{\Delta }m^2=10^1\text{eV}^2`$ the oscillation-weighted MC falls below the data in most bins (Fig. 3d). These hints concerning the parameters regime preferred by the data are borne out by a more considered analysis, as we now show.
To determine the neutrino oscillation parameters $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^22\theta `$ from our data, we construct a $`\chi ^2`$ function over the plane-of-parameters. For points $`(i,j)`$ in the physical region of the $`(\mathrm{sin}^22\theta _i,\mathrm{log}\mathrm{\Delta }m_j^2)`$ plane, we fit the MC expection to our data at each point. The MC flux normalization, $`f_\nu `$ as well as $`\mathrm{sin}^22\theta _i`$ and $`\mathrm{log}\mathrm{\Delta }m_j^2`$, is a free parameter:
$`(\chi _{data}^2)_{ij}=\chi ^2(\mathrm{sin}^22\theta _i,\mathrm{\Delta }m_j^2,f_\nu )`$
$`={\displaystyle \underset{k=1}{\overset{8}{}}}{\displaystyle \frac{\left(N_k(databkgd)f_\nu N_k(MC)\right)^2}{\sigma _k^2}}.`$ (3)
We assume that the oscillation affecting our data is purely $`\nu _\mu `$ into $`\nu _\tau `$ and that the $`\nu _e`$ data is unaffected.
The $`\chi ^2`$ is summed over data bins containing our selected (HiRes) $`\nu _\mu `$ and $`\nu _e`$ samples, where $`k=17`$ are $`\nu _\mu `$ $`\mathrm{log}\left(L/E_\nu \right)`$ bins, with $`k=8`$ containing all the $`\nu _e`$ events. The denominator $`\sigma _k^2`$ accounts for finite statistics in the neutrino Monte Carlo and for uncertainty in the rock background in the $`\nu `$ data. Not yet included are error terms which address systematic errors in the analysis, however preliminary examination shows statistical errors to be the dominant error source in the analysis. The MC counts $`N_k(MC)`$ for the $`k^{th}`$ bin are constructed using oscillation weight factors.
We find the location of minimum $`\chi _{data}^2`$, and plot $`(\mathrm{\Delta }\chi _{data}^2)_{ij}`$ which is $`(\chi _{data}^2)_{ij}(\chi _{data}^2)_{min}`$. The $`\mathrm{\Delta }\chi ^2`$ surface thereby obtained is shown in Fig. 4. A crater region of low $`\chi ^2`$ values is clearly discerned, at the bottom of which is a relatively flat basin. The lowest point $`\chi _{min}^2`$ occurs at values $`\mathrm{sin}^22\theta `$ = 0.90, $`\mathrm{\Delta }m^2=7.9\times 10^3\text{eV}^2`$, with flux normalization $`f_\nu `$ = 0.78.
An additional structure is the $`\mathrm{\Delta }\chi ^2`$ ridge which occurs at large mixing angle and for $`\mathrm{\Delta }m^2`$ above 10<sup>-2</sup> eV<sup>2</sup> . For oscillation solutions in this regime, depletion in the downward-going $`\nu _\mu `$ neutrino flux with sub-GeV energies is predicted for $`\nu _\mu `$ to $`\nu _\tau `$ oscillations by equation (1) arising from the first oscillation minimum. Our HiRes events have sufficient resolution to show such an effect if it would be present. However, no pronounced depletion is observed, and so the $`\chi ^2`$ has a high value there.
To find the region allowed for the oscillation parameters by our data at 90% confidence level (CL), we use the method of Feldman and Cousins . At each of 2500 points $`(i,j)=(\mathrm{sin}^22\theta _i,\mathrm{\Delta }m_j^2)`$ on a grid spanning the physical region of the plane parameters, we run 1000 simulated experiments. For each of the simulated sets, we find $`(\mathrm{\Delta }\chi _{90}^2)_{ij}`$ such that $`(\mathrm{\Delta }\chi _{sim}^2)_{ij}`$ is less than $`(\mathrm{\Delta }\chi _{90}^2)_{ij}`$ for 90% of the simulated experiments at $`(i,j)`$. The surface defined by local $`\mathrm{\Delta }\chi _{90}^2`$ over the oscillation parameters plane is shown in Fig. 5. Note that the surface is not a plane at $`\mathrm{\Delta }\chi _{90}^2`$ = 4.61, but rather has a concave shape. The central shaded portion is approximately $`\mathrm{\Delta }\chi ^2`$ = 4.6, however the outlying regions have $`\mathrm{\Delta }\chi ^2`$ values which are lower. At each point over the physical region, if $`(\mathrm{\Delta }\chi _{data}^2)_{ij}`$ is less than $`(\mathrm{\Delta }\chi _{90}^2)_{ij}`$, then $`(i,j)`$ belongs to the allowed region of the 90% CL contour.
The region allowed by our data at 90% CL is shown by the shaded area in Fig. 6. Although $`\chi _{min}^2`$ occurs at the location depicted by the solid circle, the relatively flat basin of our $`\mathrm{\Delta }\chi ^2`$ surface extends to lower $`\mathrm{\Delta }m^2`$ values. SuperK has reported their best fit $`\mathrm{\Delta }m^2`$ value to be $`3.2\times 10^3\text{eV}^2`$ ; our data is compatible with that as well as with somewhat higher $`\mathrm{\Delta }m^2`$ values.
## 5 PARTIALLY CONTAINED EVENTS
We plan to include more data in the above analysis. An additional data sample is comprised of $`\nu _\mu `$ flavor events which are partially contained. For each event of this category, the primary vertex is required to be $``$ 80 cm (one hadronic interaction length) from exterior surfaces of the calorimeter, and the final state must contain a non-scattering, exiting track with ionization compatible with a muon mass. This sample is useful because the assignment of $`\nu _\mu `$ flavor is reliable to better than 98%, and because the events are relatively energetic and consequently “point well” to the incident $`\nu `$ direction. The mean energy for neutrinos which initiate PCEs is estimated to be 4.7 GeV, to be compared to a mean energy of 1.3 GeV for $`\nu _\mu `$ HiRes events. The mean angular deviation of the reconstructed $`\nu _\mu `$ direction versus the true direction (in Monte Carlo) is $`14^{}`$. Unfortunately the number of $`\nu _\mu `$ PCEs is low, less than one-third the population of our $`\nu _\mu `$ HiRes sample.
In order to isolate PCE two-prong and multiprong topologies, the data events (with MC events interspersed throughout) are processed through a software filter; those which pass are scanned. The filter is designed to eliminate downward-stopping muons which have endpoint decays. In Soudan-2, an electron from a muon decay near rangeout can give rise to a small shower of $``$ 10 hits from the end of a muon track; the topology is roughly akin to that of a neutrino-induced two-prong. Consequently care must be taken to avoid remnant up-down asymmetry in PCEs introduced by the filter. The problem is avoided by requiring that there be $``$ 20 hits from a PCE vertex which are additional to the muon track.
To mitigate against cosmic-ray induced backgrounds we require that any hit in the active shield which is coincident with a PCE, must be clearly associated with the exiting muon track. Occasionally it happens that a charged pion ejected from the cavern rock is incident upon the calorimeter. If the pion penetrates by more than an interaction length and then scatters inelastically, it can mimic a $`\nu `$ PCE topology. Background events of this type are removed by requiring the net momentum of the hadronic system of a PCE to lie within the same hemisphere which contains the candidate muon track.
With these selections we obtain 31 $`\nu _\mu `$ events for which the cosmic ray induced background is less than one event. The $`\nu _\mu `$ PCE rate predicted by our Monte Carlo with no oscillations is 40 events. The distribution of PCE data in $`\mathrm{log}\left(L/E_{vis}\right)`$ can be compared, as done previously for contained HiRes $`\nu _\mu `$ events, to representative oscillation scenarios having $`\mathrm{sin}^22\theta `$ = 1.0. In Fig. 7 we observe that, compared to the data (crosses) the oscillation prediction (dashed histogram) is relatively high for $`\mathrm{\Delta }m^2=10^4\text{eV}^2`$ (Fig. 7a). However the prediction for $`\mathrm{\Delta }m^2=10^1\text{eV}^2`$ (Fig. 7d) is too low relative to the data. Figs. 7b,c suggest that the scenario preferred by the data lies in the regime between $`\mathrm{\Delta }m^2=1\text{to}\mathrm{\hspace{0.25em}7}\times 10^3\text{eV}^2`$.
## 6 PLANS
In the near future, we will include the partially contained $`\nu _\mu `$ events into our $`\chi ^2`$ fit to $`L/E`$. Additionally, a sample of upward-stopping muon events initiated by neutrino reactions below the detector has been isolated and will be analyzed for oscillation effects. And of course we will continue to accumulate and analyze new data. Our goal is to keep the Soudan-2 detector running and tuned for the change-of-beams incident at our underground site. The change will be from atmospheric $`\nu `$’s to Fermilab $`\nu _\mu `$’s in Fall 2003, at which time the detector will serve the MINOS experiment .
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# SHEDDING LIGHT ON THE DARK SIDE OF GALAXY FORMATION: SUBMILLIMETRE SURVEYS THROUGH LENSING CLUSTERS
## 1 Historical perspective
Sub-mm/mm surveys have revolutionised our understanding of star formation in the early Universe $`^\mathrm{?}`$ through the discovery of a vast population of very luminous galaxies,$`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$ clarifying the relative importance of obscured and unobscured emission. Many are extremely red $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ (a factor $``$100 in flux between 1 and 2 $`\mu `$m) and most are optically invisible, $`BVRI>26`$, even to the Hubble Space Telescope$`^{\mathrm{?},\mathrm{?}}`$ (HST).
The impact of sub-mm/mm surveys has been due to the commissioning of revolutionary bolometer cameras such as SCUBA $`^\mathrm{?}`$ on the James Clerk Maxwell Telescope and MAMBO $`^\mathrm{?}`$ at the Institut de Radioastronomie Millimétrique and the sensitivity of those devices to heavily extinguished galaxies $`^\mathrm{?}`$ – to ‘the optically dark side of galaxy formation’. $`^\mathrm{?}`$ SCUBA, in particular, has made a huge impact in cosmology through its ability to measure the bolometric output of $`1<z<5`$ dust-enshrouded galaxies (albeit with a resolution of only 14<sup>′′</sup>) whose energy distributions peak in the sub-mm band.
The foundations of sub-mm/mm cosmology are already in place, only a few years after the commissioning of SCUBA, and the community is moving rapidly to build on them, developing new telescopes and instrumentation (e.g. the Atacama Large Millimeter Array in Chile, the Large Millimeter Telescope in Mexico, and the next-generation of ground-based bolometer cameras, SCUBA-2 and BOLOCAM).
## 2 Sub-mm/mm surveys and the nature of sub-mm galaxies
The first generation of sub-mm/mm surveys, completed and ongoing, are listed in Table 1.
It is apparent that conventional blank fields have soaked up most of the time spent on cosmology surveys. Areas and rms depths range from the UKSSC 8-mJy survey’s 200 arcmin<sup>2</sup>/2.7 mJy beam<sup>-1</sup> to the UKSSC HDF $`^\mathrm{?}`$ survey’s 5.6 arcmin<sup>2</sup>/0.5 mJy beam<sup>-1</sup>, and MAMBO has now completed its first deep 1250-$`\mu `$m survey $`^\mathrm{?}`$ (450 arcmin<sup>2</sup>/0.5 mJy beam<sup>-1</sup>, fwhm 10<sup>′′</sup>).
These blank-field surveys have been tremendously successful, determining the 850-$`\mu `$m source counts above 2 mJy and thereby resolving directly up to about half of the COBE background at 850 $`\mu `$m. The deepest map, of the HDF $`^\mathrm{?}`$, has also yielded a statistical detection of the sub-mm emission from Lyman-break galaxies $`^\mathrm{?}`$.
After initial uncertainty, there is now a growing consensus amongst the sub-mm/mm community that the sources uncovered by SCUBA (and now MAMBO) are massive, intensely star-forming galaxies at $`3`$ (possibly slightly closer $`^\mathrm{?}`$), resembling ultraluminous IRAS galaxies in some respects, usually with only a tiny fraction ($`<`$1%) of their luminosity released in the rest-frame UV $`^\mathrm{?}`$ (c.f. $`^{\mathrm{?},\mathrm{?}}`$) so that many qualify as ‘extremely red objects’ $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ (EROs, $`RK>6`$).
The road to this consensus has been paved by painstaking efforts to determine the nature of individual galaxies, largely through a process of trial and error, slowly determining the most efficient techniques for identifying near-IR or optical counterparts, investigating basic properties and, in pitifully few cases, measuring redshifts $`^{\mathrm{?},\mathrm{?}}`$.
To date, deep imaging in the radio and near-IR bands $`^{\mathrm{?},\mathrm{?}}`$ have been far and away the most effective techniques, pinpointing counterparts (see Figures 1 and 2 and their captions) and facilitating spectroscopic follow up. This has culminated in several CO detections that suggest molecular gas masses consistent with the formation of elliptical galaxies.$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$
Radio flux measurements or limits at 1.4 GHz have also provided a plausible redshift distribution $`^\mathrm{?}`$ based on new photometric techniques $`^{\mathrm{?},\mathrm{?}}`$. Other techniques – mm interferometry, for example $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ – have been less successful at elucidating what we know of the SCUBA galaxy population, but clearly hold promise for the future $`^\mathrm{?}`$, particularly for very bright sources ($`>`$8 mJy at 850 $`\mu `$m) found in the field, through cluster lenses or near luminous radio galaxies $`^\mathrm{?}`$. There are hopes that broad-band spectral devices may be able to determine spectroscopic redshifts using CO transitions, regardless of the availability of plausible optical/IR counterparts, though the technical challenges are immense.
The current samples of sub-mm/mm galaxies contain a small but significant fraction of active galactic nuclei (AGN), though deep, hard-X-ray imaging $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ has so far failed to uncover the large, heavily obscured AGN population that some had suspected from the earliest follow-up work $`^\mathrm{?}`$ and from theoretical arguments $`^\mathrm{?}`$.
## 3 The problem of confusion – lifting and separating with a lens
Had the galaxies discovered in sub-mm surveys been only fractionally fainter or less numerous, a second, more sensitive generation of bolometer cameras would have been required to discover them. Early surveys $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$ would collectively have uncovered only a couple of sources – the first, an obvious AGN $`^\mathrm{?}`$ (SMM J02399$``$0136) and the second,$`^\mathrm{?}`$ a puzzle with no optical or near-IR counterpart (HDF850.1). Who can say what conclusions might have been reached and how future surveys, e.g. with FIRST, may have suffered?
We have been fortunate, then, that first-generation bolometer arrays were sufficiently sensitive to enable rapid progress in sub-mm cosmology. We have been less fortunate regarding source confusion: few would have predicted that SCUBA would reach its effective confusion limit $`^\mathrm{?}`$ within a few months of being commissioned. The deepest direct counts $`^\mathrm{?}`$ are already at the confusion limit, suggesting that further progress in constraining the intensity of the sub-mm background and the nature of the faint sub-mm population requires an innovative approach.
To probe below the confusion limit using the existing sub-mm/mm cameras requires the use of the natural magnifying glasses that provide the raison d’$`\widehat{\mathrm{e}}`$tre for this conference: gravitational lenses. Massive clusters provide a magnified (although distorted) image of a small region of the background sky; thus both the effective resolution and sensitivity of the survey are increased, as measured on the background sky. This enables surveys to probe faint flux densities without suffering confusion, albeit at the price of a distorted view.
With an accurate cluster mass model, the distortion can be corrected. The first lens survey $`^{\mathrm{?},\mathrm{?}}`$ illustrated the advantages of this approach for the counts.$`^\mathrm{?}`$ About 100% of the COBE 850-$`\mu `$m background was resolved down to 0.5 mJy. Follow-up observations,$`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$ also benefitted from achromatic gravitational amplification: not only was the effective depth of the sub-mm maps increased, but the counterparts at all other wavelengths were similarly amplified. This allows useful follow-up observations to be obtained in several hours or tens of hours using the current generation of telescopes and instrumentation: it is no coincidence that of the $``$100 known sub-mm galaxies, only a handful have reliable spectroscopic redshifts and all of these were discovered through cluster lenses.
Another advantage of using clusters is that extraordinarily deep images – X-ray, optical, IR and radio – exist or are scheduled for these fields. The HDF is the only blank field that is equally blessed. Abell 851, 1835 and 2218 (and many other cluster fields) have superb HST images and near-IR data; Abell 370, 851 and 2125 have 1.4-GHz maps with $`<`$10-$`\mu `$Jy beam<sup>-1</sup> noise levels.
## 4 Future plans and concluding remarks
Following on from the success of the earliest sub-mm cluster survey $`^{\mathrm{?},\mathrm{?}}`$, groups in the UK, Holland and Hawaii are currently undertaking more surveys with SCUBA and MAMBO that exploit cluster lenses. The latest of these will combine a long integration (equal to that obtained on the HDF) with amplification by the cores of amongst the most massive, well-constrained cluster lenses known, A 370 and A 2218.
At modest amplifications ($`A`$ 2–5), it should be possible to detect the optically-identified arclet population; probing fainter, it is likely that a new, largely unexplored class of lensed feature may appear: multiply-imaged pairs, recognised in the first HST cluster images.
These appear in the optical/near-IR as symmetric images with typical separations of 5–10<sup>′′</sup> (i.e. within a single SCUBA or MAMBO beam) and can be simply and successfully modelled as highly magnified images ($`A`$ 10–100) of very faint, compact sources which lie close to a critical line. In a well-constrained lens such as A 2218, their location in the cluster, combined with the positions and separation of any radio/IR/optical counterparts, can give the source redshift and amplification to high precision. The area of the source plane in which pairs are formed can also be estimated from the lens models, allowing their rate of occurence to be converted into an estimate of the surface density of extremely faint (tens of $`\mu `$Jy) sub-mm/mm background sources, with the bonus of crude redshift information.
Using the superb recent HST imaging of A 2218, at least 4 highly magnified pairs have been identified (from a source population with a comparable surface density to that expected for the very faint sub-mm/mm population, $``$10 arcmin<sup>-2</sup>) suggesting that the cluster amplification cross-section is high and that the chance of finding such systems is good. Failure to detect any of these highly magnified sources using SCUBA and MAMBO would indicate convergence of the source counts and can be used to impose strong limits on the surface density of very faint sources and the total intensity in resolved sources in the sub-mm background.
## Acknowledgments
I acknowledge a PPARC Advanced Fellowship and support from the Training and Mobility of Researchers Programme.
## References
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# Managing Periodically Updated Data in Relational Databases: A Stochastic Modeling Approach
## 1 Introduction and motivation
Recent developments in information management involve the transcription of data onto secondary devices in a networked environment, *e.g.*, materialized views in data warehouses and search engines, and replicas in pervasive systems. Data transcription influences the way databases define and maintain consistency. In particular, the networked environment may require periodic (rather than continuous) synchronization between the database and secondary copies, either due to paucity of resources (*e.g.*, low bandwidth or limited night windows) or to the transient characteristics of the connection. Hence, the consistency of the information in secondary copies, with respect to the transcription origin, varies over time and depends on the rate of change of the base data and on the frequency of synchronization.
Systematic approaches to the proper scheduling of transcriptions necessarily involve optimizing a trade-off between the cost of transcribing fresh information versus the cost of using obsolescent data. To do so, one must quantify, at least in probabilistic terms, this latter cost, which we call *obsolescence cost* . This paper aims to provide a comprehensive stochastic framework for quantifying time-dependent data obsolescence in replicas. Suppose we are given a relation $`R`$, a start time $`s\mathrm{}`$, and some later time $`f>s`$. We denote the extension of a relation $`R`$ at time $`t\mathrm{}`$ by $`R(t)`$. Starting from a known extension $`R(s)`$, we are interested in making probabilistic predictions about the contents of the later extension $`R(f)`$. We also suggest a cost model schema to quantify the difference between $`R(s)`$ and $`R(f)`$. Such tools assist in optimizing the synchronization process, as demonstrated in this paper. Our approach is based on techniques from the field of stochastic processes, and provides several stochastic models for content evolution in a relational database, taking referential integrity constraints into account. In particular, we make use of compound nonhomogeneous Poisson models and Markov chains; see for example . We use Poisson processes to model the behavior of tuples entering and departing relations, allowing (nonhomogeneous) time-varying behavior — *e.g.*, more intensive activity during work hours, and less intensive activity after hours and on weekends — as well as compound (bulk) insertions, that is, the simultaneous arrival of several tuples. We use Markov chains in a general modeling approach for attribute modifications, allowing the assignment of a new value to an attribute in a tuple to depend on its current value. The approach is general enough to accommodate most of the common scenarios in databases, including batch insertions and memoryless, as well as time dependent, life spans.
As motivation, consider the following two examples:
###### Example 1 (Query optimization)
Query optimization relies heavily on estimating the cardinality and value distribution of relations in a database. If these statistics are outdated and inaccurate, even the best query optimizer may formulate poor execution plans. Typically, statistics are updated only periodically, usually at the discretion of the database administrator, using utilities such as DB2’s RUNSTATS. Although some research has been devoted to speeding up statistics collection through sampling and wavelet approximations , periodic updates are unavoidable in very large databases such as IBM’s Net.Commerce , an e-business software package with roughly one hundred relations, or an SAP application, which has more than 8,000 relations and 9,000 indices. Collection of statistics becomes an even more acute problem in database federations , where the federation members do not always “volunteer” their statistics (or their cost models for that matter ), and are unwilling to burden their resources with frequent statistics collection.
In current practice, cardinality or histogram data recorded at time $`s`$ are used unchanged until the next full analysis of the database at some later time $`s^{}>s`$. If a query optimization must be performed at some time $`f(s,s^{})`$, the optimizer simply uses the statistics gathered at time $`s`$, since the time spent recomputing them may overwhelm any benefits of the query optimization. As an alternative, we suggest using a probabilistic estimate of the necessary statistics at time $`f`$. Use of these techniques might make it possible to increase the interval between statistics-gathering scans, as will be discussed in Example 7.$`\mathrm{}`$
###### Example 2 (Replication management in distributed databases)
We now consider replication management in a distributed database. Since fully synchronous replication management, in which a user is guaranteed access to the most current data, comes at a significant computational cost, most commercial distributed database providers have adopted asynchronous replication management. That is, updates to relation replicas are performed after the original transaction has committed, in accordance with the workload of the machine on which the secondary copy is stored. Asynchronous replicas are also very common in Web applications such as search engines, where Web crawlers sample Web sites periodically, and in *pervasive systems* (*e.g.*, Microsoft’s Mobile Information Server<sup>1</sup><sup>1</sup>1http://www.microsoft.com/servers/miserver/ and Café Central<sup>2</sup><sup>2</sup>2http://www.comalex.com/central.htm). In a pervasive system, a server serves many different users, each with her own unpredictable connectivity schedule and dynamically changing device capabilities. Our modeling techniques would allow client devices to reduce the rate at which they poll the server, saving both server resources and network bandwidth. We demonstrate the usefulness of stochastic modeling in this setting in Sections 3.1 and 5.6.$`\mathrm{}`$
The novelty of this paper is in developing a formal framework for modeling content evolution in relational databases. The problem of content evolution with respect to materialized views (which may be regarded as a complex form of data transcription) in databases has already been recognized. For example, in , the incompleteness of data in views was noted as being a “dynamic notion since data may be constantly added/removed from the view.” Yet, we believe that there has been no prior formal modeling of the evolution process.<sup>3</sup><sup>3</sup>3Other research efforts involve probabilistic database systems (*e.g.*, ), but this work is concerned with uncertainty in the stored data, rather than data evolution. Related research involves the containment property of a materialized view with respect to its base data: a few of the many references in this area include . However, the temporal aspects of content evolution have not been systematically addressed in this work. In , for example, the containment relationships between a materialized view $`I`$ and the “true” query result $`𝒱(D)`$, taken from a database $`D`$, can be either $`I=𝒱(D)`$ or $`I𝒱(D)`$. The latter relationship represents a situation where the materialized view stores only a partial subset of the query result. However, taking content evolution into account, it is also possible that $`I𝒱(D)`$, if tuples may be deleted from $`𝒱(D)`$ and $`I`$ is periodically updated. Moreover, modifications to the base data may result in both $`I𝒱(D)`$ and $`I𝒱(D)`$.
Refresh policies for materialized views have been previously discussed in the literature (*e.g.*, and ). Typically, materialized views are refreshed immediately upon updates to the base data, at query time (as in ), or using snapshot databases (as in ). The latter approach can produce obsolescent materialized views. A combination of all three approaches appears in . Our methodology differs in that we do not assume an *a priori* association of a materialized view with a refresh policy, but instead design policies based on their transcription and obsolescence costs.
A preliminary attempt to describe the time dependency of updates in the context of Web management was given in , which suggests a simple homogeneous Poisson process to model the updating of Web pages. We suggest instead a nonhomogeneous compound Poisson model, which is far more flexible, and yet still tractable. In addition, the work in supposes that transcriptions are performed at uniform time intervals, mainly because “crawlers cannot guess the best time to visit each site.” We show in this paper that our model of content evolution gives rise to other, better transcription policies.
In , a trade-off mechanism was suggested to decide between the use of a cache or recomputation from base data by using range data, computed at the source. In this framework, an update is “pushed” to a replication site whenever updated data falls outside a predetermined interval, or whenever a query requires current data. The former requires the client and the server to be in touch continuously, in case the server needs to track down the client, which is not always realistic (either because the server does not provide such services, or because the overhead for such services undermines the cost-effectiveness of the client). The latter requirement puts the burden of deciding whether to refresh the data on the client, without providing it with any model for the evolution of the base data. We attempt to fill this gap by providing a stochastic model for content evolution, which allows a client to make judicious requests for current data. Other work in related areas (*e.g.*, ) has considered various alternatives for pushing updated data from a server to a cache on the client side. Lazy replica-update policies using replication graphs have also been discussed in, for example, . This work, however, does not take the data obsolescence into account, and is primarily concerned with transaction throughput and timely updates, subject to network constraints.
As with models in general, our model is an idealized representation of a process. To be useful, we wish to make predictions based on tractable analytical calculations, rather than detailed, computationally intensive simulations. Therefore, we restrict our modeling to some of the more basic tools of applied probability theory, specifically those relating to Poisson processes and Markov chains. Texts such as contain the necessary reference material on Markov chains and Poisson processes, and specifically on nonhomogeneous Poisson processes. Poisson processes can model a world where data updates are independent from one another. In databases with widely distributed access, *e.g.*, Web interfacing databases, such an independence assumption seems plausible, as was verified in .
The rest of the paper is organized as follows: Section 1.1 introduces some basic notation. Section 2 provides a content evolution model for insertions and deletions, while Section 3 discusses data modifications. We shall introduce preliminary results of fitting the insertion model parameters to real data feeds in Section 4. A cost model and transcription policies that utilize it follow in Section 5, highlighting the practical impact of the model. Conclusions and topics for further research are provided in Section 6.
### 1.1 Notational preliminaries
In what follows, we denote the set of attributes and relations in the database by $``$ and $``$, respectively. Each $`R`$ consists of a set of attributes $`𝒜(R)`$, and also has a *primary key* $`𝒦(R)`$, which is a nonempty subset of $`𝒜(R)`$. Each attribute $`A`$ has a *domain* $`\mathrm{dom}A`$, which we assume to be a finite set, and for any subset of attributes $`𝒜=\{A_1,A_2,\mathrm{},A_k\}`$, we let $`\mathrm{dom}𝒜=\mathrm{dom}A_1\times \mathrm{dom}A_2\times \mathrm{}\times \mathrm{dom}A_k`$ denote the compound domain of $`𝒜`$. We denote by $`r.A(t)`$ the value of attribute $`A`$ in tuple $`r`$ at time $`t`$, and similarly use $`r.𝒜(t)`$ for the value of a compound attribute. For a given time $`t`$, subset of attributes $`𝒜𝒜(R)`$, and value $`v=v_1,v_2,\mathrm{},v_k\mathrm{dom}𝒜`$, we define $`R_{𝒜,v}(t)=\{rR(t)|(r.A_1(t)=v_1)(r.A_2(t)=v_2)\mathrm{}(r.A_k(t)=v_k)\}`$. We also define $`\widehat{R}_𝒜(t)`$ to be the *histogram* of values of $`𝒜`$ at time $`t`$, that is, for each value $`v\mathrm{dom}𝒜`$, $`\widehat{R}_𝒜(t)`$ associates a nonnegative integer $`\widehat{R}_{𝒜,v}(t)`$, which is the cardinality of $`R_{𝒜,v}(t)`$.<sup>4</sup><sup>4</sup>4This vector can be computed exactly and efficiently using indices. Alternatively, in the absence of an index for a given attribute, statistical methods (such as “probabilistic” counting , sampling-based estimators , and wavelets ) can be applied. This notation, and well as other symbols used throughout the paper, are also summarized in Table 1.
## 2 Modeling insertions and deletions
This section introduces the stochastic models for insertions and deletions. Section 2.1 discusses insertions, while deletions are discussed in section 2.2. Section 2.3 combines the effect of insertions and deletions on a relation’s cardinality. We conclude with a discussion of non-exponential life spans in Section 2.4. We defer discussing model validation until Section 4.
### 2.1 Insertion
We use a nonhomogeneous Poisson process with instantaneous arrival rate $`\lambda _R:\mathrm{}[0,\mathrm{})`$ to model the occurrence of *insertion events* into $`R`$. That is, the number of insertion events occurring in any interval $`(s,f]`$ is a Poisson random variable with expected value $`\mathrm{\Lambda }_R(s,f)=_s^f\lambda _R(t)𝑑t.`$ A homogeneous Poisson process may be considered as the special case where $`\lambda _R(t)`$ is equal to a constant $`\lambda _R>0`$ for all $`t`$, yielding $`\mathrm{\Lambda }_R(s,f)=_s^f\lambda _R(t)𝑑t=_s^f\lambda _R𝑑t=\lambda _R(fs)`$.
We now consider the interarrival time distribution of the nonhomogeneous Poisson process. We first define the nonhomogeneous exponential distribution, as follows:
###### Definition 1 (Nonhomogeneous exponential distribution)
Let $`\varphi :\mathrm{}[0,\mathrm{})`$ be a integrable function. Given some $`s\mathrm{}`$, a random variable $`V`$ is said to have a *nonhomogeneous exponential* distribution (denoted by $`V\mathrm{Exp}_s(\varphi ())`$) if $`V`$’s density function is
$$p(\tau )=\{\begin{array}{cc}\varphi (s+\tau )\mathrm{exp}\left(_0^\tau \varphi (s+u)𝑑u\right),\hfill & \tau 0\hfill \\ 0,\hfill & \tau <0.\hfill \end{array}$$
It is worth noting that if $`\varphi (t)`$ is constant, $`p(\tau )`$ is just a standard exponential distribution. We shall now show that, as with homogeneous Poisson processes, the interarrival time of insertion events is distributed like an exponential random variable, $`L_{R,s}`$, but with a time-varying density function.
###### Lemma 1
At any time $`s`$, the amount of time $`L_{R,s}`$ to the next insertion event is distributed like $`\mathrm{Exp}_s(\lambda _R())`$. The probability of an insertion event occurring during $`(s,f]`$ is $`\mathrm{P}\{L_{R,s}<fs\}=1e^{\mathrm{\Lambda }_R(s,f)}`$.
Proof. Let $`\{N(t),t0\}`$ be a nonhomogeneous Poisson process with intensity function $`\lambda _R(t)`$, which implies $`\mathrm{P}\left\{N(f)N(s)=0\right\}=e^{\mathrm{\Lambda }_R(s,f)}`$. Now, the chance that no new tuple was inserted during $`(s,f]`$ is the same as the chance that the process $`N()`$ has no arrivals during $`(s,f]`$, that is, $`e^{\mathrm{\Lambda }_R(s,f)}`$. The chance that a new tuple was inserted during $`(s,f]`$ is just the complement of the chance of no arrivals, namely,
$$\mathrm{P}\left\{L_{R,s}<fs\right\}=\{N(f)N(s)1\}=1P\{N(f)N(s)=0\}=1e^{\mathrm{\Lambda }_R(s,f)}.$$
Taking the derivative of this expression with respect to $`f`$ and making a change of variables, the probability density of the time until the next insertion from time $`s`$ is $`p(\tau )=\lambda _R(s+\tau )e^{\mathrm{\Lambda }_R(s,s+\tau )}`$. Thus, $`L_{R,s}\mathrm{Exp}_s(\lambda _R())`$.
At insertion event $`i`$, a random number of tuples $`\mathrm{\Delta }_{R,i}^+`$ are inserted, allowing us to model bulk insertions. A *bulk insertion* is the simultaneous arrival of multiple tuples, and may occur because the tuples are related, or because of limitations in the implementation of the server. For example, e-mail servers may process an input stream periodically, resulting in bulk updates of a mailbox. Assuming that the $`\{\mathrm{\Delta }_{R,i}^+\}`$ are independent and identically distributed (IID), then the stochastic process $`\{B_R(t),t0\}`$ representing the cumulative number of insertions through time $`t`$ is a *compound Poisson* process (*e.g.*, , pp. 87-88). We let $`B_R(s,f)`$ denote the number of insertions falling into the interval $`(s,f]`$. The expected number of inserted tuples during $`(s,f]`$ may be computed via $`\mathrm{E}\left[B_R(s,f)\right]=_s^f\lambda _R(t)\mathrm{E}\left[\mathrm{\Delta }_R^+\right]𝑑t=\mathrm{E}\left[\mathrm{\Delta }_R^+\right]_s^f\lambda _R(t)𝑑t=\mathrm{E}\left[\mathrm{\Delta }_R^+\right]\mathrm{\Lambda }_R(s,f).`$ Here, $`\mathrm{\Delta }_R^+`$ represents a generic random variable distributed like the $`\{\mathrm{\Delta }_{R,i}^+\}`$.
We now consider three simple cases of this model:
##### General nonhomogeneous Poisson process:
Assume that $`\mathrm{E}\left[\mathrm{\Delta }_R^+\right]=1`$. The expected number of insertions simplifies to $`\mathrm{E}\left[B_R(s,f)\right]=\mathrm{E}\left[\mathrm{\Delta }_R^+\right]\mathrm{\Lambda }_R(s,f)=1\mathrm{\Lambda }_R(s,f)=\mathrm{\Lambda }_R(s,f)`$.
##### Homogeneous Poisson process:
Assume once more that $`\mathrm{E}\left[\mathrm{\Delta }_R^+\right]=1`$. Assume further that $`\lambda _R(t)`$ is a constant function, that is, $`\lambda _R(t)=\lambda _R`$ for all times $`t`$. In this case, as shown above, $`\mathrm{\Lambda }_R(s,f)`$ takes on the simple form of $`\lambda _R(fs)`$. Thus, $`\mathrm{E}\left[B_R(s,f)\right]=\mathrm{\Lambda }_R(s,f)=\lambda _R(fs)`$. The interarrival times are distributed as $`\mathrm{Exp}(\lambda _R)`$, the exponential distribution with parameter $`\lambda _R`$.
##### Recurrent piecewise-constant Poisson process:
A simple kind of nonhomogeneous Poisson process can be built out of homogeneous Poisson processes that repeat in a cyclic pattern. Given some length of time $`T`$, such as one day or one week, suppose that the arrival rate function $`\lambda _R(t)`$ of the recurrent Poisson process repeats every $`T`$ time units, that is, $`\lambda _R(t)=\lambda _R\left(tTt/T\right)`$ for all $`t`$. Furthermore, the interval $`[0,T)`$ is partitioned into a finite number of subsets $`J_1,\mathrm{},J_K`$, with $`\lambda _R(t)`$ constant throughout each $`J_k`$, $`k=1,\mathrm{},K`$. Finally, each $`J_k`$ is in turn composed of a finite number of half-open intervals of the form $`[s,f)`$. For instance, $`T`$ might be one day, with $`K=24`$ and $`J_1=[0\text{:}00,1\text{:}00),J_2=[1\text{:}00,2\text{:}00),\mathrm{},J_{24}=[23\text{:}00,0\text{:}00)`$. As another simple example, $`T`$ might be one week, and $`K=2`$. The subset $`J_1`$ would consist of a firm’s normal hours of operation, say $`[9`$:$`00,18`$:$`00)`$ for each weekday, and $`J_2=[0,T)\backslash J_1`$ would denote all “off-hour” times. Formalisms like those of could also be used to describe such processes in a more structured way. We term this class of Poisson processes to be *recurrent piecewise-constant* — abbreviated *RPC*.
It is worth noting that, in client-server environments, the insertion model should typically be formed from the client’s point of view. Therefore, if the server keeps a database from which many clients transcribe data, the modeling of insertions for a given client should only include the part of the database the client actually transcribes. Therefore, if a “road warrior” is interested only in new orders for the 08904 zip code area, the insertion model for that client should concentrate on that zip code, ignoring the arrival orders from other areas.
#### 2.1.1 The complexity of computing $`\mathrm{\Lambda }_R(s,f)`$
$`\mathrm{\Lambda }_R(s,f)`$, the Poisson expected value, is computed by integrating the model parameter $`\lambda _R(t)`$ over the interval $`[s,f]`$. Standard numerical methods allow rapid approximation of this definite integral even if no closed formula is known for the indefinite integral. However, the complexity of this calculation depends on the information-theoretic properties of $`\lambda _R(t)`$ \[36, Section 1\].
For our purposes, however, simple models of $`\lambda _R(t)`$ are likely to suffice. For example, if $`\lambda _R(t)`$ is a polynomial of degree $`d0`$, the integration can be performed in $`\mathrm{O}(d+1)`$ time. Consider next a piecewise-polynomial Poisson process: the time line is divided into intervals such that, in each time interval, $`\lambda _R(t)`$ can be written as a polynomial. The complexity of calculating $`\mathrm{\Lambda }_R(s,f)`$ in this case is $`\mathrm{O}(n(d+1))`$, where $`n`$ is the number of segments in the time interval $`(s,f]`$, and $`d`$ is the highest degree of the $`n`$ polynomials.
Further suppose that the piecewise-polynomial process is recurrent in a similar manner to the RPC process, that is, given some fixed time interval $`T`$, $`\lambda _R(t)=\lambda _R\left(tTt/T\right)`$ for all $`t`$. Note that the RPC Poisson process is the special case of this model in which $`d=0`$. If there are $`c`$ segments in the interval $`[0,T]`$, then the complexity of calculating $`\mathrm{\Lambda }_R(s,f)`$ becomes $`\mathrm{O}(c(d+1))`$, regardless of the length of the interval $`[s,f]`$. This reduction occurs because, for all intervals of the form $`[kT,(k+1)T][s,f]`$ for which $`k`$ is an integer, the integral $`_{kT}^{(k+1)T}\lambda _R(t)𝑑t`$ is equal to $`_0^T\lambda _R(t)𝑑t`$, which only needs to be calculated once.
In Section 4, we demonstrate the usefulness of the RPC model for one specific application. We hypothesize that a recurrent piecewise-polynomial process of modest degree (for example, $`d`$=3) will be sufficient to model most systems we are likely to encounter, and so the complexity of computing $`\mathrm{\Lambda }_R(s,f)`$ should be very manageable.
### 2.2 Deletion
We allow for two distinct deletion mechanisms. First, we assume individual tuples have their own intrinsic stochastic life spans. Second, we assume that tuples are deleted to satisfy referential integrity constraints when tuples in other relations are deleted. These two mechanisms are combined in a tuple’s overall probability of being deleted. Let $`R`$ and $`S`$ be two relations such that $`𝒦(S)`$ is a foreign key of $`S`$ in $`R`$. We refer to $`S`$ as a *primary relation* of $`R`$. Consider the directed multigraph $`G`$ whose vertices consist of all relations $`R`$ in the database, and whose edges are of the form $`R,S`$, where $`S`$ is a primary relation of $`R`$. The number of edges $`R,S`$ is the number of foreign keys of $`S`$ in $`R`$ for which integrity constraints are enforced. We assume that $`G(R)`$ has no directed cycles. Let $`G(R)`$ denote the subgraph of $`G`$ consisting of $`R`$ and all directed paths starting at $`R`$. We denote the vertices of this subgraph by $`S(R)`$.
###### Example 3 (Referential integrity constraints in Net.Commerce)
IBM’s Net.Commerce is supported by a DB2 database with about a hundred relations interrelated through foreign keys. For demonstration purposes, consider a sample of seven relations in the Net.Commerce database. Figure 1 is a pictorial representation of the multigraph $`G`$ of these seven relations. The MERCHANT relation provides data about merchant profiles, the SCALE and DISCCALC relations are for computing price discounts, the CATEGORY and CGRYREL relations assist in categorizing products, and the ORDERS and SHIPTO relations contain information about orders. The six relations, SCALE, DISCCALC, CATEGORY, CGRYREL, ORDERS, and SHIPTO have a foreign key to the MERCHANT relation, through MERCHANT’s primary key (MERFNBR). Integrity constraints are enforced between the SCALE relation and the MERCHANT relation, as long as SCALE.SCLMENBR (the foreign key to MERCHANT.MERFNBR) does not have the value NULL. That is, unless a NULL value is assigned to the MERCHANT.MERFNBR attribute, a deletion of a tuple in MERCHANT results in a deletion of all tuples in SCALE such that SCALE.SCLMENBR$`=`$MERCHANT.MERFNBR. DISCCALC has a foreign key to the SCALE relation, through SCALE’s primary key (SCLRFNBR). There are two attributes of CGRYREL that serve as foreign keys to the CATEGORY relation, through CATEGORY’s primary key (CGRFNBR). Finally, SHIPTO contains shipment information of each product in an order, and therefore it has a foreign key to ORDERS through its primary key (ORFNBR). $`\mathrm{}`$
With regard to intrinsic deletions within a relation, we assume that each tuple $`rR(s)`$ has a stochastic remaining life span $`L_{R,s}^\mathrm{I}`$. This random variable is identically distributed for each $`rR(s)`$, and is independent of the remaining life span of any other tuple and of $`r`$’s age at time $`s`$ (see Section 2.4 for a discussion of tuples with a non-memoryless life span). Specifically, we will assume that the chance of $`rR(t)`$ being deleted in the time interval $`[t,t+\mathrm{\Delta }t]`$ approaches $`\mu _R(t)\mathrm{\Delta }t`$ as $`\mathrm{\Delta }t0`$, for some function $`\mu _R:\mathrm{}[0,\mathrm{})`$. We define $`M_R(s,f)=_s^f\mu _R(t)𝑑t`$.
###### Lemma 2
$`L_{R,s}^\mathrm{I}\mathrm{Exp}_s(\mu _R())`$. The probability that a tuple $`rR(s)`$ is deleted by time $`f`$, given that no corresponding tuple in $`S(R)\backslash \{R\}`$ is deleted, is $`\mathrm{P}\{L_{R,s}^\mathrm{I}<fs\}=1e^{M_R(s,f)}`$.
Proof. Let $`rR(s)`$ be a randomly chosen tuple, and assume that no corresponding tuple to $`r`$ in $`S(R)\backslash \{R\}`$ is deleted. The proof is identical to that of Lemma 1, replacing $`\lambda _R(t)`$ with $`\mu _R(t)`$ and $`\mathrm{\Lambda }_R(s,f)`$ by $`M_R(s,f)`$.
#### 2.2.1 Deletion and referential integrity
For any $`rR(s)`$ and any relation $`SS(R)`$, we define $`w(r,S)`$ to be the number of tuples in $`S`$ whose deletion would force deletion of $`r`$ in order to maintain referential integrity. This value can be between $`0`$ and the number of paths from $`R`$ to $`S`$ in $`G(R)`$. For example, if $`rR=\mathrm{𝙲𝙶𝚁𝚈𝚁𝙴𝙻}`$ of Figure 1, then $`0w(r,\mathrm{𝙲𝙰𝚃𝙴𝙶𝙾𝚁𝚈})2`$ and $`0w(r,\mathrm{𝙼𝙴𝚁𝙲𝙷𝙰𝙽𝚃})3`$. For completeness, we define $`w(r,R)=1`$. Each tuple in $`S`$ has an independent remaining lifetime distributed as $`\mathrm{Exp}_s(\mu _S())`$, and if any of the $`w(r,S)`$ tuples corresponding to $`r`$ is deleted, then $`r`$ must be immediately deleted, to maintain referential integrity constraints. We use $`p_R(s,f)`$ to denote the probability that a randomly chosen tuple in $`R(s)`$ survives until time $`f`$.
###### Lemma 3
$`p_R(s,f)=\mathrm{E}_{rR(s)}\left[\mathrm{exp}\left(_{SS(R)}w(r,S)M_S(s,f)\right)\right]`$, where $`\mathrm{E}_{rR(s)}[]`$ denotes expectation over random selection of tuples in $`R(s)`$.
Proof. Considering all $`SS(R)`$, and using the well-known fact that if $`L_i\mathrm{Exp}_s(\mu _i())`$ for $`i=1,\mathrm{},k`$ are independent, then
$$\mathrm{min}\{L_1,\mathrm{},L_k\}\mathrm{Exp}_s\left(\underset{i=0}{\overset{k}{}}\mu _i()\right),$$
(1)
we conclude that the remaining lifetime of $`r`$ (denoted $`L_{R,s}`$) has a nonhomogeneous exponential distribution with intensity function $`_{SS(R)}w(r,S)\mu _S()`$. The probability of a given tuple $`rR(s)`$ surviving through time $`f`$ is thus
$$\mathrm{exp}\left(_s^f\left(\underset{SS(R)}{}w(r,S)\mu _S(t)\right)𝑑t\right)=\mathrm{exp}\left(\underset{SS(R)}{}w(r,S)M_S(s,f)\right),$$
and the probability that a randomly chosen tuple in $`R(s)`$ survives until time $`f`$ is therefore
$$p_R(s,f)=\mathrm{E}_{rR(s)}\left[\mathrm{exp}\left(\underset{SS(R)}{}w(r,S)M_S(s,f)\right)\right].$$
(2)
The complexity analysis of integrating $`\mu _S(t)`$ over time to obtain $`M_S(s,f)`$ is similar to that of Section 2.1.1. However, the computation required by Lemma 3 may be prohibitive, in the most general case, because it requires knowing the empirical distribution of the $`w(r,S)`$ over all $`rR(s)`$ for all $`SS(R)`$. This empirical distribution can be computed accurately by computing for each tuple, upon insertion, the number of tuples in any $`SS(R)`$ with a comparable foreign key, using either histograms or by directly querying the database. Maintaining this information requires $`\mathrm{O}(\left|R(s)\right|\left|S(R)\right|)`$ space. This complexity can be reduced using a manageably-sized sample from $`R(s)`$. Our initial analysis of real-world applications, however, indicates that in many cases, $`w(r,S)`$ takes on a much simpler form, in which $`w(r,S)`$ is identical for all $`rR(s)`$. We term such a typical relationship between $`R`$ and $`SS(R)`$ a *fixed multiplicity*, as defined next:
###### Definition 2
The pair $`R,S`$, where $`SS(R)`$, has *fixed multiplicity* if $`w(r,S)`$ is identical for all tuples in $`R`$. In this case, we denote its common value by $`w(R,S)`$. $`\mathrm{}`$
###### Example 4 (Fixed multiplicies in Net.Commerce)
Consider the example multigraph of Figure 1. Both DISCALC and SCALE reference MERCHANT. It is clear that the discount calculation of a product (as stored in DISCALC) cannot reference a different merchant than SCALE. The only exception is when the foreign key in SCALE is assigned with a null value. If this is the case, however, there is only a single tuple in MERCHANT whose deletion requires the deletion of a tuple in DISCALC. Thus, for any tuple $`r`$ DISCALC, $`w(r,\text{SCALE})=w(r,\text{MERCHANT})=1`$ and therefore $`\text{DISCALC},\text{SCALE}`$ and $`\text{DISCALC},\text{MERCHANT}`$ both have fixed multiplicity of $`1`$. Now consider CGRYREL. Since each tuple in CGRYREL describes the relationship between a category and a subcategory, it is clear that its two foreign keys to CATEGORY must always have distinct values. Thus, $`\text{CGRYREL},\text{CATEGORY}`$ has a fixed multiplicity, and $`w(\text{CGRYREL},\text{CATEGORY})=2`$.$`\mathrm{}`$
As the following lemma shows, fixed multiplicities permit great simplification in computing $`p_R(s,f)`$.
###### Lemma 4
If $`R,S`$ has fixed multiplicity for all $`SS(R)`$, $`p_R(s,f)=\mathrm{exp}(\stackrel{~}{M}_R(s,f))`$, where $`\stackrel{~}{M}_R(s,f)=_s^f\stackrel{~}{\mu }_R(t)𝑑t`$ and $`\stackrel{~}{\mu }_R(t)=_{SS(R)}w(R,S)\mu _S(t)`$.
Proof.
$`p_R(s,f)`$ $`=\mathrm{E}_{rR(s)}\left[\mathrm{exp}\left({\displaystyle \underset{SS(R)}{}}w(r,S)M_S(s,f)\right)\right]`$
$`=\mathrm{E}_{rR(s)}\left[\mathrm{exp}\left({\displaystyle _s^f}\left({\displaystyle \underset{SS(R)}{}}w(r,S)\mu _S(t)\right)𝑑t\right)\right]`$
$`=\mathrm{E}_{rR(s)}\left[\mathrm{exp}\left({\displaystyle _s^f}\left({\displaystyle \underset{SS(R)}{}}w(R,S)\mu _S(t)\right)𝑑t\right)\right]`$
$`=\mathrm{exp}\left({\displaystyle \underset{SS(R)}{}}\left(w(R,S){\displaystyle _s^f}\mu _S(t)𝑑t\right)\right)`$
$`=\mathrm{exp}\left({\displaystyle \underset{SS(R)}{}}w(R,S)M_S(s,f)\right).`$
Since $`w(R,S)`$ is fixed and constant over time, no additional statistics need to be collected for it. As a final note, it is worth noting that in certain situations, another alternative may also be available. Let $`\{N_R(t),t0\}`$ be a nonhomogeneous Poisson process with intensity function $`\widehat{\mu }_R(t)`$, modeling the occurrence of *deletion events* in $`R`$. At deletion event $`i`$, a random number $`\mathrm{\Delta }_i^{}`$ tuples are deleted from $`R`$. Generally speaking, this kind of model cannot be accurate, since it ignores that each deletion causes a reduction in the number of remaining tuples, and thus presumably a change in the spacing of subsequent deletion events. However, it may be reasonably accurate for large databases with either a stable or steadily growing number of tuples, or whenever the time interval $`(s,f]`$ is sufficiently small. Statistical analysis of the database log would be required to say whether the model is applicable. If the model is valid, then the stochastic process $`\{D_R(t),t0\}`$ representing the cumulative number of deletions through time $`t`$, can be taken to be a compound Poisson process. The expected number of deleted tuples during $`(s,f]`$ may be computed via
$$\mathrm{E}\left[D_R(t)\right]=_s^f\mu _R(t)\mathrm{E}\left[\mathrm{\Delta }^{}\right]𝑑t=M_R(s,f)\mathrm{E}\left[\mathrm{\Delta }^{}\right],$$
where $`\mathrm{\Delta }^{}`$ is a generic random variable distributed like the $`\{\mathrm{\Delta }_i^{}\}`$.
### 2.3 Tuple survival: the combined effect of insertions and deletions
Some tuples inserted during $`(s,f]`$ may be deleted by time $`f`$. Let the random variable $`X_R(s,f)`$ denote the number of tuples inserted during the interval $`(s,f]`$ that survive through time $`f`$. Consider any tuple inserted into $`R`$ at time $`t(s,f]`$, and denote its chance of surviving through time $`f`$ by $`\widehat{p}_R(t,f)`$. For any $`SS(R)`$ and $`t(s,f]`$, let $`W(R,S,t)`$ be a random variable denoting the value of $`w(r,S)`$, given that $`r`$ was inserted into $`R`$ at time $`t`$. Let $`W(R,t)`$ denote the random vector, of length $`\left|S(R)\right|`$, formed by concatenating the $`W(R,S,t)`$ for all $`SS(R)`$.
###### Lemma 5
$`\widehat{p}_R(t,f)=\mathrm{E}_{W(R,t)}\left[\mathrm{exp}\left(_{SS(R)}W(R,S,t)M_S(t,f)\right)\right]`$. When $`R,S`$ has fixed multiplicity for all $`SS(R)`$, then $`\widehat{p}_R(t,f)=p_R(t,f)=\mathrm{exp}(\stackrel{~}{M}_R(t,f))`$.
Proof. Let $`L_{R,t}`$ denote the lifetime of a tuple inserted into $`R`$ at time $`t`$. Similarly to the proof of Lemma 3, we know that $`L_{R,t}\mathrm{Exp}_t(_{SS(R)}W(R,S,t)\mu _S())`$. The probability such a tuple survives through time $`f`$ is the random quantity
$$\mathrm{exp}\left(_t^f\left(\underset{SS(R)}{}W(R,S,t)\mu _S(\tau )\right)𝑑\tau \right)=\mathrm{exp}\left(\underset{SS(R)}{}W(R,S,t)M_S(t,f)\right).$$
Considering all the possible elements of the vector $`W(R,t)`$, we then obtain
$$\widehat{p}_R(t,f)=\mathrm{E}_{W(R,t)}\left[\mathrm{exp}\left(\underset{SS(R)}{}W(R,S,t)M_S(t,f)\right)\right],$$
Assume now that $`R,S`$ has fixed multiplicity for all $`SS(R)`$. Consequently, we replace $`W(R,S,t)`$ with $`w(R,S)`$. Drawing on the proof of the previous lemma,
$`\widehat{p}_R(t,f)`$ $`=\mathrm{E}_{W(R,t)}\left[\mathrm{exp}\left({\displaystyle \underset{SS(R)}{}}w(R,S)M_S(t,f)\right)\right]`$
$`=\mathrm{exp}\left({\displaystyle _t^f}\left({\displaystyle \underset{SS(R)}{}}w(R,S)\mu _S(\tau )\right)𝑑\tau \right)`$
$`=p_R(t,f)\text{.}`$
The following proposition establishes the formula for the expected value of $`X_R(s,f)`$.
###### Proposition 1
$`\mathrm{E}\left[X_R(s,f)\right]=\stackrel{~}{\mathrm{\Lambda }}_R(s,f)\mathrm{E}\left[\mathrm{\Delta }_R^+\right]`$, where $`\stackrel{~}{\mathrm{\Lambda }}_R(s,f)=_s^f\lambda _R(t)\widehat{p}_R(t,f)𝑑t`$. In the simple case where each insertion involves exactly one tuple, $`X_R(s,f)\mathrm{Poisson}(\stackrel{~}{\mathrm{\Lambda }}_R(s,f))`$.
Proof. Let $`N`$ be the number of insertion events in $`(s,f]`$, and let their times be $`\{T_1,T_2,\mathrm{},T_N\}`$. Suppose that $`N=n`$ and that insertion event $`i`$ happens at time $`t_i(s,f]`$. Event $`i`$ inserts a random number of tuples $`\mathrm{\Delta }_{R,i}^+`$, each of which has probability $`\widehat{p}_R(t_i,f)`$ of surviving through time $`f`$. Therefore, the expected number of tuples surviving through $`f`$ from insertion event $`i`$ is $`\mathrm{E}\left[\mathrm{\Delta }_R^+\right]\widehat{p}_R(t_i,f)`$. Consequently,
$$\mathrm{E}[X_R(s,f)|N=n,T_1=t_1,T_2=t_2,\mathrm{},T_n=t_n]=\mathrm{E}\left[\mathrm{\Delta }_R^+\right]\underset{i=1}{\overset{n}{}}\widehat{p}_R(t_i,f).$$
Next, we recall, given that $`N=n`$, that the times $`T_i`$ of the insertion events are distributed like $`n`$ independent random variables with probability density function $`\lambda _R(t)/\mathrm{\Lambda }_R(s,f)`$ on the interval $`(s,f]`$. Thus,
$`\mathrm{E}\left[X_R(s,f)|N=n\right]`$ $`=\mathrm{E}_{T_1,\mathrm{},T_n}\left[\mathrm{E}\left[\mathrm{\Delta }_R^+\right]{\displaystyle \underset{i=1}{\overset{n}{}}}\widehat{p}_R(T_i,f)\right]`$
$`=\mathrm{E}\left[\mathrm{\Delta }_R^+\right]{\displaystyle \underset{i=1}{\overset{n}{}}}\left({\displaystyle _s^f}\widehat{p}_R(t,f){\displaystyle \frac{\lambda _R(t)}{\mathrm{\Lambda }_R(s,f)}}𝑑t\right)`$
$`=n\left({\displaystyle \frac{\mathrm{E}\left[\mathrm{\Delta }_R^+\right]\stackrel{~}{\mathrm{\Lambda }}_R(s,f)}{\mathrm{\Lambda }_R(s,f)}}\right).`$
Finally, removing the conditioning on $`N=n`$, we obtain
$`\mathrm{E}\left[X_R(s,f)\right]`$ $`=\mathrm{E}_N\left[N\left({\displaystyle \frac{\mathrm{E}\left[\mathrm{\Delta }_R^+\right]\stackrel{~}{\mathrm{\Lambda }}_R(s,f)}{\mathrm{\Lambda }_R(s,f)}}\right)\right]`$
$`=\mathrm{\Lambda }_R(s,f)\left({\displaystyle \frac{\mathrm{E}\left[\mathrm{\Delta }_R^+\right]\stackrel{~}{\mathrm{\Lambda }}_R(s,f)}{\mathrm{\Lambda }_R(s,f)}}\right)`$
$`=\mathrm{E}\left[\mathrm{\Delta }_R^+\right]\stackrel{~}{\mathrm{\Lambda }}_R(s,f).`$
In the case that $`\mathrm{\Delta }_R^+`$ is always $`1`$, we may use the notion of a *filtered* Poisson process: if we consider only tuples that manage to survive until time $`f`$, the chance of a single insertion in time interval $`[t,t+\mathrm{\Delta }t]`$ no longer has the limiting value $`\lambda _R(t)\mathrm{\Delta }t`$, but instead $`\lambda _R(t)\widehat{p}_R(t,f)\mathrm{\Delta }t`$. Therefore, the insertion of surviving tuples can be viewed as a nonhomogeneous Poisson process with intensity function $`\lambda _R(t)\widehat{p}_R(t,f)`$ over the time interval $`(s,f]`$, so $`X_R(s,f)\mathrm{Poisson}(\stackrel{~}{\mathrm{\Lambda }}_R(s,f))`$.
In the general case, the computation of $`\stackrel{~}{\mathrm{\Lambda }}_R(s,f)`$ will require approximation by numerical integration techniques; the complexity of this calculation will depend on the information-theoretic properties of $`\lambda _R()`$ and the $`\mu _S()`$, $`SS(R)`$, but is unlikely to be burdensome if these functions are reasonably smoothly-varying. In one important special case, however, the complexity of computing $`\stackrel{~}{\mathrm{\Lambda }}_R(s,f)`$ is essentially the same as that of calculating $`\mathrm{\Lambda }_R(s,f)`$: suppose that for some constants $`\alpha (R,S)`$, $`SS(R)`$, one has that $`\mu _S(t)=\alpha (R,S)\lambda _R(t)`$ for all $`t`$. That is, the general insertion and deletion activity level of the relations in $`S(R)`$ all vary proportionally to some common fluctuation pattern. In this case, we have $`\stackrel{~}{\mu }_R(t)=\alpha (R)\lambda _R(t)`$ and $`\stackrel{~}{M}_R(s,f)=\alpha (R)\mathrm{\Lambda }_R(s,f)`$ for all $`t,s,f`$, where $`\alpha (R)=_{SS(R)}\alpha (R,S)`$. Making a substitution $`u(t)=\mathrm{\Lambda }_R(t,f)`$, we have:
$`\stackrel{~}{\mathrm{\Lambda }}_R(s,f)`$ $`={\displaystyle _s^f}\lambda _R(t)\mathrm{exp}(\alpha (R)\mathrm{\Lambda }_R(t,f))𝑑t`$
$`={\displaystyle _s^f}\left({\displaystyle \frac{d\mathrm{\Lambda }_R(t,f)}{dt}}\right)\mathrm{exp}(\alpha (R)\mathrm{\Lambda }_R(t,f))𝑑t`$
$`={\displaystyle _s^f}\mathrm{exp}(\alpha (R)u(t))du(t)`$
$`={\displaystyle _{u(s)}^{u(f)}}e^{\alpha (R)u}𝑑u`$
$`={\displaystyle \frac{1}{\alpha (R)}}\left(1e^{\alpha (R)\mathrm{\Lambda }(s,f)}\right),`$
so $`\stackrel{~}{\mathrm{\Lambda }}_R(s,f)`$ can be calculated directly from $`\mathrm{\Lambda }(s,f)`$.
We define the random variable $`Y_R(s,f)`$ to be the number of tuples in $`R(s)`$ that survive through time $`f`$.
###### Proposition 2
$`\mathrm{E}\left[Y_R(s,f)\right]=p_R(s,f)\left|R(s)\right|`$ and $`\mathrm{E}\left[\left|R(f)\right|\right]=p_R(s,f)\left|R(s)\right|+\stackrel{~}{\mathrm{\Lambda }}_R(s,f)\mathrm{E}\left[\mathrm{\Delta }_R^+\right]`$.
Proof. Each tuple in $`R(s)`$ has a survival probability of $`p_R(s,f)`$, which yields that $`\mathrm{E}\left[Y_R(s,f)\right]=p_R(s,f)\left|R(s)\right|`$. By the definitions of $`Y_R(s,f)`$ and $`X_R(s,f)`$, one has that
$$\left|R(f)\right|=Y_R(s,f)+X_R(s,f),$$
so therefore
$$\mathrm{E}\left[\left|R(f)\right|\right]=\mathrm{E}\left[Y_R(s,f)\right]+\mathrm{E}\left[X_R(s,f)\right]=p_R(s,f)\left|R(s)\right|+\stackrel{~}{\mathrm{\Lambda }}_R(s,f)\mathrm{E}\left[\mathrm{\Delta }_{R}^{}{}_{}{}^{+}\right].$$
In cases where deletions may also be accurately modeled as a compound Poisson process, we have
$`\mathrm{E}\left[\left|R(f)\right|\right]`$ $`=\mathrm{E}\left[\left|R(s)\right|\right]+B_R(s,f)D_R(s,f)`$
$`=\left|R(s)\right|+\mathrm{\Lambda }_R(s,f)\mathrm{E}\left[\mathrm{\Delta }^+\right]M_R(s,f)\mathrm{E}\left[\mathrm{\Delta }^{}\right].`$
###### Example 5 (The homogeneous case)
Assume that $`\mathrm{E}\left[\mathrm{\Delta }_R^+\right]=1`$, that $`R,S`$ has fixed multiplicity for all $`SS(R)`$, and furthermore $`\lambda _R(t)`$ and $`\mu _S(t)`$, for all $`SS(R)`$, are constant functions, that is, $`\lambda _R(t)=\lambda _R`$ for all times $`t`$ and $`\mu _S(t)=\mu _S`$ for all $`SS(R)`$ and times $`t`$. Then $`\mathrm{\Lambda }_R(s,f)=\lambda _R(fs)`$ and $`M_R(s,f)=\mu _R(fs)`$. Thus, letting $`\stackrel{~}{\mu }_R=_{SS(R)}w(R,s)\mu _S`$,
$$\stackrel{~}{\mathrm{\Lambda }}_R(s,f)=_s^f\lambda _Re^{\stackrel{~}{\mu }_R(ts)}𝑑t=\frac{\lambda _R}{\stackrel{~}{\mu }_R}\left(1e^{\stackrel{~}{\mu }_R(fs)}\right),$$
and assuming $`\mathrm{\Delta }_{R,i}^+=1`$ for all $`i>0`$,
$$\mathrm{E}\left[\left|R(f)\right|\right]=\left|R(s)\right|e^{\stackrel{~}{\mu }_R(fs)}+\frac{\lambda _R}{\stackrel{~}{\mu }_R}\left(1e^{\stackrel{~}{\mu }_R(fs)}\right)=\frac{\lambda _R}{\stackrel{~}{\mu }_R}+e^{\stackrel{~}{\mu }_R(fs)}\left(\left|R(s)\right|\frac{\lambda _R}{\stackrel{~}{\mu }_R}\right).$$
$`\mathrm{}`$
### 2.4 Tuples with non-exponential life spans
We now consider the possibility that tuples in $`R`$ have a stochastic life span $`L_R^\mathrm{I}`$ that is not memoryless, but rather has some general cumulative distribution function $`G_R`$. For example, if tuples in $`R`$ correspond to pieces of work in process on a production floor, the likelihood of deletion might rise the longer the tuple has been in existence. Let us consider a single relation, and thus no referential integrity constraints. For any tuple $`r`$, let $`b(r)`$ denote the time it was created. We next establish the expected cardinality of $`R`$ at time $`f`$.
###### Proposition 3
In the case that tuples in $`R`$ have lifetimes with a general cumulative distribution function $`G_R`$,
$$\mathrm{E}\left[\left|R(f)\right|\right]=\left|R(s)\right|\mathrm{E}_{rR(s)}\left[\frac{1G_R(fb(r))}{1G_R(sb(r))}\right]+\mathrm{E}\left[\mathrm{\Delta }_R^+\right]_s^f\lambda _R(t)\left(1G_R(ft)\right)𝑑t.$$
(3)
Proof. Let $`L_R^\mathrm{I}`$ denote a generic random variable with cumulative distribution $`G_R`$. The probability of $`rR(s)`$ surviving throughout $`(s,f]`$ is then
$$\mathrm{P}\left\{L_R^\mathrm{I}fb(r)|L_R^\mathrm{I}sb(r)\right\}=\frac{1G_R(fb(r))}{1G_R(sb(r))},$$
and therefore the expected number of tuples in $`R(s)`$ that survive through time $`f`$ is
$$\mathrm{E}\left[Y_R(s,f)\right]=\underset{rR(s)}{}\left(\frac{1G_R(fb(r))}{1G_R(sb(r))}\right)=\left|R(s)\right|\mathrm{E}_{rR(s)}\left[\frac{1G_R(fb(r))}{1G_R(sb(r))}\right].$$
We now consider a tuple $`r`$ inserted at some time $`t(s,f]`$. The probability that such a tuple survives through time $`f`$ is simply $`\widehat{p}_R(t,f)=1G_R(ft)`$. By reasoning similar to the proof of Lemma 5,
$$\mathrm{E}\left[X_R(s,f]\right]=\mathrm{E}\left[\mathrm{\Delta }_{R}^{}{}_{}{}^{+}\right]\stackrel{~}{\mathrm{\Lambda }}_R(s,f)=\mathrm{E}\left[\mathrm{\Delta }_{R}^{}{}_{}{}^{+}\right]_s^f\lambda _R(t)\left(1G_R(ft)\right)𝑑t.$$
The conclusion then follows from $`\mathrm{E}\left[\left|R(f)\right|\right]=\mathrm{E}\left[Y_R(s,f)\right]+\mathrm{E}[X_R(s,f)]`$
It is worth noting that, as opposed to the memoryless case presented above, the calculation of $`\mathrm{E}\left[Y_R(s,f)\right]`$ requires remembering the commit times $`b(r)`$ of all tuples $`rR(s)`$, or equivalently the ages of all such tuples. Of course, for large relations $`R`$, a reasonable approximation could be obtained by using a manageably-sized sample to estimate
$$\mathrm{E}_{rR(s)}\left[\frac{1G_R(fb(r))}{1G_R(sb(r))}\right].$$
It is likely that the integral in (3) will require general numerical integration, depending on the exact form of $`G_R`$.
### 2.5 Summary
In this section, we have provided a model for the insertion and deletion of tuples in a relational database. The immediate benefit of this model is the computation of the expected relation cardinality ($`\mathrm{E}\left[\left|R(f)\right|\right]`$), given an initial cardinality and insertion and tuple life span parameters. Relation cardinality has proven to be an important property in many database tools, including query optimization and database tuning. Reasonable assumptions regarding constant multiplicity allow, once appropriate statistics have been gathered, a rapid computation of cardinalities in this framework. Section 4 elaborates on statistics gathering and model validation.
A note regarding tuples with non-exponential life spans is now warranted. For the case of a single relation, non-exponential life spans add only a moderate amount of complexity to our model, namely the requirement to store at least an approximation of the distribution of tuples ages in $`R(s)`$. For multiple relations with referential integrity constraints, however, the complexity of dealing with general tuple life spans is much greater. First, to estimate the cardinality of $`R(f)`$, we must keep (approximate) tuple age distributions for all relations in $`S(R)`$. Second, because the tuple life span distributions of some of the members of $`S(R)`$ are not memoryless, we cannot combine them with a simple relation like (1). Furthermore, in attempting to find the distribution of the remaining life span of a particular tuple $`rR(s)`$, it may become necessary to consider the issue of the correlation of ages of tuples in $`R(s)`$ with the ages of the corresponding tuples in other relations of $`S(R)`$. Because of these complications, we defer further consideration of non-exponential tuple life spans to future research.
## 3 Modeling data modification
This section describes various ways to model the modification of the contents of tuples. We start with a general approach, using Markov chains, followed by several special cases where the amount of computation can be greatly reduced.
### 3.1 Content-dependent updates
In this section, we model the modification of the contents of tuples as a finite-state continuous-time Markov chain, thus assuming dependence on tuples’ previous contents. For each relation $`R`$, we allow for some (possibly empty) subset $`𝒞(R)𝒜(R)`$ of its attributes to be subject to change over the lifetime of a tuple. We do not permit primary key fields to be modified, that is, $`𝒞(R)𝒦(R)=\mathrm{}`$.
Attribute values may change at time instants called *transition events*, which are the transition times of the Markov chain. We assume that the spacing of transition events is memoryless with respect to the age of a tuple (although it may depend on the time and the current value of the attribute, as demonstrated below). For any attribute $`A`$, tuple $`r`$, time $`s`$, and value $`v\mathrm{dom}A`$ with $`r.A(s)=v`$, the time remaining until the next transition event for $`r.A`$ is a random variable $`\tau _{v,s}^{R,A}`$ with the distribution $`\mathrm{Exp}_s(\mathrm{}_v^{R,A}\gamma _{R,A}())`$, where $`\gamma _{R,A}:\mathrm{}[0,\mathrm{})`$ is a function giving the general instantaneous rate of change for the attribute, and $`\mathrm{}_v^{R,A}`$ is a nonnegative scalar which we call the *relative exit rate* of $`v`$. We define $`\mathrm{\Gamma }_{R,A}(s,f)=_s^f\gamma _{R,A}(t)𝑑t`$. When a transition event occurs from state $`u\mathrm{dom}A`$, attribute $`A`$ changes to $`v\mathrm{dom}A`$ with probability $`P_{u,v}^{R,A}`$.
Suppose $`𝒜=\{A_1,A_2,\mathrm{},A_k\}`$ is an independently varying set of attributes, and $`v=v_1,v_2,\mathrm{},v_k\mathrm{dom}𝒜`$ is a compound value. Then the time until the next transition event for $`r.𝒜`$ is $`\tau _{v,s}^{R,𝒜}=\mathrm{min}\{\tau _{v_1,s}^{R,A_1},\mathrm{},\tau _{v_k,s}^{R,A_k}\}`$. As a rule, we will assume that the modification processes for the attributes of a relation are independent, so $`\tau _{v,s}^{R,𝒜}\mathrm{Exp}_s(_{i=1}^k\mathrm{}_{v_i}^{R,A_i}\gamma _{A_i,R}())`$. When the functions $`\gamma _{R,A_i}`$ are identical for $`i=1,\mathrm{},k`$, we define $`\gamma _{R,𝒜}=\gamma _{R,A_i}`$ and $`\mathrm{}_v^{R,𝒜}=_{i=1}^k\mathrm{}_{v_i}^{R,A_i}`$, so $`\tau _{v,s}^{R,𝒜}\mathrm{Exp}_s(\mathrm{}_v^{R,𝒜}\gamma _{R,𝒜}())`$. To justify the assumption of independence, we note that coordinated modifications among attributes can be modeled by replacing the coordinated attributes with a single compound attribute (this technique requires that the attributes have identical $`\gamma _{R,𝒜}()`$ functions, which is reasonable if they change in a coordinated way).
Under these assumptions, let $`\overline{𝒞}(R)`$ denote a partition of $`𝒞(R)`$ into subsets $`𝒜`$ such that any two attributes $`A_1,A_2𝒞(R)`$ vary dependently iff they are in the same $`𝒜\overline{𝒞}(R)`$.
###### Example 6 (First alteration time)
For a relation $`R`$ and time $`s`$, we define $`\mathrm{{\rm Y}}_{R,s}`$ to be the amount of time until the next change in $`R`$, be it a tuple insertion, a tuple deletion, or an attribute modification. Also, for any $`SS(R)`$, let $`D(R,S,s)`$ denote the number of tuples in $`S(s)`$ whose deletion would force the deletion of some tuple in $`R(s)`$ $`(`$and therefore $`D(R,R,s)=\left|R(s)\right|)`$. The following proposition establishes the distribution of $`\mathrm{{\rm Y}}_{R,s}`$.
###### Proposition 4
$`\mathrm{{\rm Y}}_{R,s}\mathrm{Exp}_s(\zeta _R())`$, where
$$\zeta _R(t)=\lambda _R(t)+\underset{SS(R)}{}D(R,S,s)\mu _S(t)+\underset{𝒜\overline{𝒞}(R)}{}h(R,A,s)\gamma _{R,𝒜}(t)$$
and $`h(R,A,s)=_{v\mathrm{dom}𝒜}\widehat{R}_{𝒜,v}(s)\mathrm{}_v^{R,A}`$. The probability of any alteration to $`R`$ in the time interval $`(s,f]`$ is $`1e^{Z_R(s,f)}`$, where
$$Z_R(s,f)=\mathrm{\Lambda }_R(s,f)+\underset{SS(R)}{}D(R,S,s)M_S(t)+\underset{𝒜\overline{𝒞}(R)}{}h(R,A,s)\mathrm{\Gamma }_{R,𝒜}(s,f)$$
Proof. Let $`\mathrm{{\rm Y}}_{R,s}^\mathrm{I}`$, $`\mathrm{{\rm Y}}_{R,s}^\mathrm{M}`$, and $`\mathrm{{\rm Y}}_{R,s}^\mathrm{D}`$ be the times until the next insertion, modification, and deletion in $`R`$, respectively. From Section 2, we have that $`\mathrm{{\rm Y}}_{R,s}^\mathrm{I}\mathrm{Exp}_s(\lambda _R())`$. Now, for each $`SS(R)`$, there are $`D(R,S,s)`$ tuples whose deletion would cause a deletion in $`R`$. The time until deletion of any such $`rSS(R)`$ is distributed like $`\mathrm{Exp}_s(\mu _S())`$. The deletion processes for all these tuples are independent across all of $`S(R)`$, so we can use (1) to conclude that
$$\mathrm{{\rm Y}}_{R,s}^\mathrm{D}\mathrm{Exp}_s\left(\underset{SS(R)}{}D(R,S,s)\mu _S()\right).$$
From the preceding discussion, we have
$$\mathrm{{\rm Y}}_{R,s}^\mathrm{M}=\tau _{v,s}^{𝒞(R),R}\mathrm{Exp}_s\left(\underset{𝒜\overline{𝒞}(R)}{}\mathrm{}_{r.𝒜(s)}^{R,𝒜}\gamma _{R,𝒜}()\right).$$
Since $`\mathrm{{\rm Y}}_{R,s}=\mathrm{min}\{\mathrm{{\rm Y}}_{R,s}^\mathrm{I},\mathrm{{\rm Y}}_{R,s}^\mathrm{M},\mathrm{{\rm Y}}_{R,s}^\mathrm{D}\}`$, we therefore have, again using independence and (1), that $`\mathrm{{\rm Y}}_{R,s}\mathrm{Exp}_s(\zeta _R())`$, where
$`\zeta _R(t)`$ $`=\lambda _R(t)+{\displaystyle \underset{SS(R)}{}}D(R,S,s)\mu _S(t)+{\displaystyle \underset{rR(s)}{}}\left[{\displaystyle \underset{𝒜\overline{𝒞}(R)}{}}\mathrm{}_{r.𝒜(s)}^{R,𝒜}\gamma _{R,𝒜}(t)\right]`$
$`=\lambda _R(t)+{\displaystyle \underset{SS(R)}{}}D(R,S,s)\mu _S(t)+{\displaystyle \underset{𝒜\overline{𝒞}(R)}{}}h(R,A,s)\gamma _{R,𝒜}(t).`$
Integrating over $`(s,f]`$ results in
$`Z_R(s,f)`$ $`={\displaystyle _s^f}\zeta _R(t)𝑑t`$
$`=\mathrm{\Lambda }_R(s,f)+{\displaystyle \underset{SS(R)}{}}D(R,S,s)M_S(t)+{\displaystyle \underset{𝒜\overline{𝒞}(R)}{}}h(R,A,s)\gamma _{R,𝒜}(s,f).`$
Therefore, the probability of any alteration to $`R`$ in the time interval $`(s,f]`$ is
$$\mathrm{P}\left\{\mathrm{{\rm Y}}_{R,s}<fs\right\}=1e^{Z_R(s,f)}$$
Example 6 continued (First alteration transcription policy) Suppose the user wishes to refresh her replica of relation $`R`$ whenever the probability that it contains any inaccuracy exceeds some threshold $`\pi `$, a tactic we call the *first alteration policy*. Then, a refresh is required at time $`f`$ if $`1e^{Z_R(s,f)}>\pi `$.$`\mathrm{}`$
Given $`𝒜`$, we describe the transition process for the value $`r.𝒜`$ over time by probabilities
$$P_{u,v}^{R,𝒜}(s,f)=\mathrm{P}\{r.𝒜(f)=v|r.𝒜(s)=u\},$$
for any two values $`u,v\mathrm{dom}𝒜`$ and times $`s<f`$. Under the assumption of independence,
$$P_{u,v}^{R,𝒜}(s,f)=\underset{i=1}{\overset{k}{}}P_{u_i,v_i}^{R,A_i}(s,f).$$
(4)
Given any simple attribute $`A`$, we define $`q_{u,v}^{R,A}`$, the *relative transition rate* from $`u`$ to $`v`$, by
$$q_{u,v}^{R,A}=\mathrm{}_u^{R,A}P_{u,v}^{R,A}.$$
Given a set of attributes $`𝒜`$ with identical $`\gamma _{R,A}()`$ functions, the compound transition rate $`q_{u,v}^{R,A}`$ may be computed via
$$q_{u,v}^{R,𝒜}=\mathrm{}_u^{R,𝒜}P_{u,v}^{R,𝒜}=\left(\underset{i=1}{\overset{k}{}}\mathrm{}_{u_i}^{R,A_i}\right)\left(\underset{i=1}{\overset{k}{}}P_{u_i,v_i}^{R,A_i}\right).$$
(5)
Let $`Q^{R,A}`$ be the matrix of $`q_{u,v}^{R,A}`$, where $`q_{u,u}^{R,A}=\mathrm{}_u^{R,A}`$.
###### Proposition 5
The matrix $`P^{R,A}(s,f)`$ of elements $`P_{u,v}^{R,A}(s,f)`$ is given by the matrix exponential formula
$$P^{R,A}(s,f)=\mathrm{exp}\left(\mathrm{\Gamma }_{R,A}(s,f)Q^{R,A}\right)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }_{R,A}(s,f)^n}{n!}\left(Q^{R,A}\right)^n.$$
(6)
Proof. Consider a continuous-time Markov chain on the same state space $`\mathrm{dom}A`$, and with the same instantaneous transition probabilities $`P_{u,v}^{R,A}`$, where $`u,v\mathrm{dom}A`$. However, in the new chain, the holding time in each state $`v`$ is simply a homogeneous exponential random variable with arrival rate $`\mathrm{}_v^{R,A}`$. We call this system the *linear-time* chain, to distinguish it from the original chain. Define $`\overline{P}_{u,v}^{A,R}(t)`$ to be the chance that the linear-time chain is in state $`v`$ at time $`t`$, given that it is in state $`u`$ at time $`0`$. Standard results for finite-state continuous time Markov chains imply that
$$\overline{P}_{u,v}^{R,A}(t)=\mathrm{exp}\left(tQ^{R,A}\right)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{t^n}{n!}\left(Q^{R,A}\right)^n.$$
By a transformation of the time variable, we then assert that
$$P_{u,v}^{R,A}(s,f)=\overline{P}_{u,v}^{R,A}(\mathrm{\Gamma }_{R,A}(s,f)),$$
from which the result follows.
###### Example 7 (Query optimization, revisited)
As the following proposition shows, our model can be used to estimate the histogram of a relation $`R`$ at time $`f`$. A query optimizer running at time $`f`$ could use expected histograms, calculated in this manner, instead of the old histograms $`\widehat{R}_A(s)`$.
###### Proposition 6
Assume that $`w(r,S)`$, for all $`SS(R)`$, is independent of the attribute values $`r.A(s)`$ for all $`A𝒞(R)`$. Let $`\widehat{\omega }_u^{R,A}(t)`$ denote the probability that $`r.A(t)=u`$, given that $`r`$ is inserted into $`R`$ at time $`t`$. Then, for all $`v\mathrm{dom}A`$,
$`\mathrm{E}\left[\widehat{R}_{A,v}(f)\right]`$ $`=p_R(s,f){\displaystyle \underset{u\mathrm{dom}A}{}}\widehat{R}_{A,u}(s)P_{u,v}^{R,A}(s,f)`$
$`+\mathrm{E}[\mathrm{\Delta }_R^+]{\displaystyle \underset{u\mathrm{dom}A}{}}\left({\displaystyle _s^f}\widehat{\omega }_u^{R,A}(t)\widehat{p}_R(s,f)\lambda _R(t)P_{u,v}^{R,A}(t,f)𝑑t\right).`$ (7)
Proof. We first compute the expected number of surviving tuples $`r`$ whose values $`r.A`$ migrate to $`v`$. Given a value $`u\mathrm{dom}A`$, there are $`\widehat{R}_{A,u}(s)`$ tuples at time $`s`$ such that $`r.A(s)=u`$. Using the previous results, the expected number of these tuples surviving through time $`f`$ is $`\widehat{R}_{A,u}(s)p_R(s,f)`$, and the probability of each surviving tuple $`r`$ having $`r.A(f)=v`$ is $`P_{u,v}^{R,A}(s,f)`$. Using the independence assumption and summing over all $`u\mathrm{dom}A`$, one has that the expected numbers of tuples in $`R(s)`$ that survive through $`f`$ and have $`r.A(f)=v`$ is
$$\underset{u\mathrm{dom}A}{}\widehat{R}_{A,u}(s)p_R(s,f)P_{u,v}^{R,A}(s,f)=p_R(s,f)\underset{u\mathrm{dom}A}{}\widehat{R}_{A,u}(s)P_{u,v}^{R,A}(s,f).$$
We next consider newly inserted tuples. Recall that $`\widehat{\omega }_u^{R,A}(t)`$ denotes the probability that $`r.A(t)=u`$, given that $`r`$ is inserted into $`R`$ at time $`t`$. Suppose that an insertion occurs at time $`t(s,f]`$. The expected number of tuples $`r`$ created at this insertion that both survive until $`f`$ and have $`r.A(f)=v`$ is
$$\widehat{p}_R(t,f)\underset{u\mathrm{dom}A}{}\omega _u^{R,A}(t)P_{u,v}^{R,A}(t,f).$$
By logic similar to Proposition 5, one may then conclude that the expected number of newly-inserted tuples that survive through time $`f`$ and have $`r.A(f)=v`$ is
$$\mathrm{E}[\mathrm{\Delta }_R^+]\underset{u\mathrm{dom}A}{}\left(_s^f\widehat{\omega }_u^{R,A}(t)\widehat{p}_R(s,f)\lambda _R(t)P_{u,v}^{R,A}(t,f)𝑑t\right).$$
The result follows by adding the last two expressions.
Example 7 continued *We next consider whether the complexity of calculating (7) is preferable to recomputing the histogram vector $`\widehat{R}_A(f)`$. This topic is quite involved and depends heavily on the specific structure of the database (*e.g.*, the availability of indices) and the specific application (*e.g.*, the concentration of values in a small subset of an attribute’s domain). In what follows, we lay out some qualitative considerations in deciding whether calculating (7) would be more efficient than recalculating $`\widehat{R}_A(f)`$ “from scratch.” Experimentation with real-world application is left for further research.*
*Generally speaking, direct computation of the histogram of an attribute $`A`$ (in the absence of an index for $`A`$) can be done by either scanning all tuples (although sampling may also be used) or scanning a modification log to capture changes to the prior histogram vector $`\widehat{R}_A(s)`$ during $`(s,f]`$. Therefore, the recomputation can be performed in $`\mathrm{O}(\mathrm{min}\{\left|R(f)\right|,T(s,f)\})`$ time, where $`T(s,f)`$ denotes the total number of updates during $`(s,f]`$. Whenever $`\left|R(f)\right|`$ and $`T(s,f)`$ are both large — *i.e.*, the database is large and the transaction load is high — the straightforward techniques will be relatively unattractive. As for the estimation technique, it will probably work best when $`\left|domA\right|`$ is small (for example, for a binary attribute) or whenever the subset of actually utilized values in the domain is small. In addition, commercial databases recompute the entire histogram as a single, atomic task. Formula (7), on the other hand, can be performed on a subset of the attribute values. For example, in the case of exact matching (say, a condition of the form $`\mathrm{𝚆𝙷𝙴𝚁𝙴}A=v`$), it is sufficient to compute $`\widehat{R}_{A,v}(f)`$, rather than the full $`\widehat{R}_A(f)`$ vector. Finally, it is worth noting that the computing the expected value of $`\widehat{R}_{A,v}(f)`$ via (7) does not require locking $`R`$, while a full histogram recomputation may involve extended periods of locking.*$`\mathrm{}`$
We next consider the number of tuples in $`R(s)`$ that have survived through time $`f`$ without being modified, which we denote $`Y_R^{}(s,f)`$. The expectation of this random variable is
$$\mathrm{E}\left[Y_R^{}(s,f)\right]=p_R(s,f)\underset{v\mathrm{dom}𝒜}{}\widehat{R}_{𝒜,v}(s)P_{v,v}^{R,𝒜}(s,f)$$
We let $`Y_R^+(s,f)=Y_R(s,f)Y_R^{}(s,f)`$ denote the number of tuples in $`R(s)`$ that have survived through time $`f`$ and were modified; it follows from the linearity of the $`\mathrm{E}[]`$ operator that
$$\mathrm{E}\left[Y_R^+(s,f)\right]=\mathrm{E}\left[Y_R(s,f)\right]\mathrm{E}\left[Y_R^{}(s,f)\right].$$
#### 3.1.1 Complexity analysis of content-dependent updates
In practice, as with computing a scalar exponential, only a limited number of terms will be needed to compute the sum (6) to machine precision. It is worth noting that efficient means of calculating (6) are a major topic in the field of computational probability.
In the case of a compound attribute $`𝒜=\{A_1,A_2,\mathrm{},A_k\}`$ with independently varying components, it will be computationally more efficient to first calculate the individual transition probability matrices $`P^{A_i,R}(s,f)`$ via (6), and then calculate the joint probability matrix $`P_{u,v}^{R,𝒜}(s,f)`$ using (4), rather than first finding the joint exit rate matrix $`Q^{R,𝒜}`$ via (5) and then applying (6). The former approach would involve repeated multiplications of square matrices of size $`\left|\mathrm{dom}A_i\right|`$, for $`i=1,\mathrm{},k`$, resulting in a computational complexity of $`\mathrm{O}(_{i=1}^kn_i\left|\mathrm{dom}A_i\right|^\nu )`$, where $`n_i`$ is the number of iterations needed to compute the sum (6) to machine precision, and the complexity of multiplying two $`n\times n`$ matrices is $`\mathrm{O}(n^\nu )`$.<sup>5</sup><sup>5</sup>5$`\nu =3`$ for the standard method and $`\nu =\mathrm{log}_27`$ for Strassen’s and related methods. The latter would involve multiplying square matrices of size $`_{i=1}^k\left|\mathrm{dom}A_i\right|`$, resulting in the considerably worse complexity of $`\mathrm{O}(n(_{i=1}^k\left|\mathrm{dom}A_i\right|)^\nu )`$, where $`n`$ is the number of iterations needed to obtain the desired precision.
### 3.2 Simplified modification models
We next introduce several possible simplifications of the general Markov chain case. To do so, we start by differentiating numeric domains from non-numeric domains. Certain database attributes $`A𝒜`$, such as prices and order quantities, represent numbers, and numeric operations such as addition are meaningful for these attributes. For such attributes, one can easily define a distance function between two attribute values, as we shall see below. We call the domains $`\mathrm{dom}A`$ of such attributes *numeric domains*, and denote the set of all attributes with numeric domains by $`𝒩`$. All other attributes and domains are considered *non-numeric*.<sup>6</sup><sup>6</sup>6Distance metrics can also be defined for complex data types such as images. We leave the handling of such cases to further research. It is worth noting that not all numeric data necessarily constitute a numeric domain. Consider, for example, a customer relation $`R`$ whose primary key is a customer number. Although the customer number consists of numeric symbols, it is essentially an arbitrary identification string for which arithmetic operations like addition and subtraction are not intrinsically meaningful for the database application. We consider such attributes to be non-numeric.
#### 3.2.1 Domain lumping
To make our data modification model more computationally tractable, it may be appropriate, in many cases, to simplify the Markov chain state space for an attribute $`A`$ so that it is much smaller than $`\mathrm{dom}A`$. Suppose, for example, that $`A`$ is a 64-character string representing a street address. Restricting to 96 printable characters, $`A`$ may assume on the order of $`96^{64}10^{126}`$ possible values. It is obviously unnecessary, inappropriate, and intractable to work with a Markov chain with such an astronomical number of states.
One possible remedy for such situations is referred to as *lumping* in the Markov chain literature . In our terminology, suppose we can partition $`\mathrm{dom}A`$ into a collection of sets $`\{V\}_{V𝒱}`$ with the property that $`\left|𝒱\right|\left|\mathrm{dom}A\right|`$ and
$$U,V𝒱,u,u^{}U\underset{vV}{}q_{u,v}^{R,A}=\underset{vV}{}q_{u^{},v}^{R,A}.$$
Then, one can model the transitions between the “lumps” $`V𝒱`$ as a much smaller Markov chain whose set of states is $`𝒱`$, with the transition rate from $`U𝒱`$ to $`V𝒱`$ being given by the common value of $`_{vV}q_{u,v}^{R,A}`$, $`uU`$. If we are interested only in which lump the attribute is in, rather than its precise value, this smaller chain will suffice. Using lumping, the complexity of the computation is directly dependent on the number of lumps. We now give a few simple examples:
###### Example 8 (Lumping into a binary domain)
Consider the street address example just discussed. Fortunately, if an address has changed since time $`s`$, the database user is unlikely to be concerned with how different it is from the address at time $`s`$, but simply whether it is different. Thus, instead of modeling the full domain $`\mathrm{dom}A`$, we can represent the domain via the simple binary set $`\{0,1\}`$, where $`0`$ indicates that the address has not changed since time $`s`$, and $`1`$ indicates that it has. We assume that the exit rates $`q_{v,r.A(s)}^{R,A}`$ from all other addresses $`v\mathrm{dom}A`$ back to the original value $`r.A(s)`$ all have the identical value $`\theta ^{}`$. In this case, one has $`P_{0,1}^{R,A}=P_{1,0}^{R,A}=1`$, and the behavior of the attribute is fully captured by the exit rates $`\mathrm{}_0^{R,A}=q_{0,1}^{R,A}`$ and $`\mathrm{}_1^{R,A}=q_{1,0}^{R,A}`$. We will abbreviate these quantities by $`\theta `$ and $`\theta ^{}`$, respectively.
Using standard results for a two-state continuous-time Markov chain \[35, Section VI.3.3\], we conclude that
$`P_{0,0}^{R,A}(s,f)`$ $`={\displaystyle \frac{\theta ^{}+\theta e^{(\theta +\theta ^{})\mathrm{\Gamma }_{R,A}(s,f)}}{\theta +\theta ^{}}}`$ (8)
$`P_{0,1}^{R,A}(s,f)`$ $`={\displaystyle \frac{\theta \theta e^{(\theta +\theta ^{})\mathrm{\Gamma }_{R,A}(s,f)}}{\theta +\theta ^{}}}.`$ (9)
$`\mathrm{}`$
###### Example 9 (Web crawling)
As an even simpler special case, consider a Web crawler (*e.g.*, ). Such a crawler needs to visit Web pages upon change to re-process their content, possibly for the use of a search engine. Recalling Example 8, one may define a boolean attribute Modified in a relation that collects information on Web pages. Modified is set to True once the page has changed, and back to False once the Web crawler has visit the page. Therefore, once a page has been modified to True, it cannot be modified back to False before the next visit of the Web crawler. In the analysis of Example 8, one can set $`\theta ^{}=0`$, resulting in $`P_{0,0}^{A,R}(s,f)=e^{\theta \mathrm{\Gamma }_{R,A}(s,f)}`$ and $`P_{0,1}^{A,R}(s,f)=1e^{\theta \mathrm{\Gamma }_{R,A}(s,f)}`$. $`\mathrm{}`$
#### 3.2.2 Random walks
Like large non-numeric domains, many numeric domains may also be cumbersome to model directly via Markov chain techniques. For example, a 32-bit integer attribute can, in theory, take $`2^{32}4\times 10^9`$ distinct values, and it would be virtually impossible to directly form, much less exponentiate, a full transition rate matrix for a Markov chain of this size.
Fortunately, it is likely that such attributes will have “structured” value transition patterns that can be modeled, or at least closely approximated, in a tractable way. As an example, we consider here a random walk model for numeric attributes.
In this case, we still suppose that the attribute $`A`$ is modified only at transition event times that are distributed as described above. Letting $`t_i`$ denote the time of transition event $`i`$, with $`t_0=s`$, we suppose that at transition event $`i`$, the value of attribute $`A`$ is modified according to
$$r.A(t_i)=r.A(t_{i1})+\mathrm{\Delta }A_i,$$
where $`\mathrm{\Delta }A_i`$ is a random variable. We suppose that the random variables $`\left\{\mathrm{\Delta }A_i\right\}`$ are IID, that is, they are independent and share a common distribution with mean $`\delta `$ and variance $`\sigma ^2`$. Defining
$$\mathrm{\Delta }A(s,f)=\underset{i:t_i(s,f]}{}\mathrm{\Delta }A_i,$$
we obtain that $`\{\mathrm{\Delta }A(s,f),fs\}`$ is a nonhomogeneous compound Poisson process, and $`r.A(f)=r.A(s)+\mathrm{\Delta }A(s,f)`$. From standard results for compound Poisson processes, we then obtain for each tuple $`rR(s)`$ that $`\mathrm{E}[r.A(f)]=r.A(s)+\mathrm{\Gamma }_{R,A}(s,f)\delta `$.
It should be stressed that such a model must ultimately be only an approximation, since a random walk model of this kind would, strictly speaking, require an infinite number of possible states, while $`\mathrm{dom}A`$ is necessarily finite for any real database. However, we still expect it to be accurate and useful in many situations, such as when $`r.A(s)`$ and $`\mathrm{E}[r.A(f)]`$ are both far from largest and smallest possible values in $`\mathrm{dom}A`$.
#### 3.2.3 Content-independent overwrites
Consider the simple case in which $`P_{u,v}^{R,A}(s,f)`$ is independent of $`u`$ once a transition event has occurred. Let $`𝒜𝒞(R)`$ be a set of attributes $`A`$ with identical $`\gamma _{R,A}`$ functions, and let $`\mathrm{\Gamma }_{R,𝒜}(s,t)=\mathrm{\Gamma }_{R,A}(s,t)`$ for any $`A𝒜`$. We define a probability distribution $`\omega _{R,𝒜}`$ over $`\mathrm{dom}𝒜`$, and assume that at each transition event, a new value for $`𝒜`$ is selected at random from this distribution, without regard to the prior value of $`r.𝒜`$. It is thus possible that a transition event will leave $`r.𝒜`$ unchanged, since the value selected may be the same one already stored in $`r`$. For any tuple $`rR(s)R(f)`$ and $`u\mathrm{dom}𝒜`$, we thus compute the probability $`P_{u,u}^{R,𝒜}(s,f)`$ that the value of $`r.𝒜`$ remains unchanged at $`u`$ at time $`f`$ to be
$`P_{u,u}^{R,𝒜}(s,f)`$ $`=\mathrm{P}\left\{\tau _u^{R,𝒜}>fs\right\}+\mathrm{P}\left\{\tau _u^{R,𝒜}fs\right\}\omega _{R,𝒜}(u)`$
$`=e^{\mathrm{}_u^{R,𝒜}\mathrm{\Gamma }_{R,𝒜}(s,f)}+\left(1e^{\mathrm{}_u^{R,𝒜}\mathrm{\Gamma }_{R,𝒜}(s,f)}\right)\omega _{R,𝒜}(u)`$
$`=e^{\mathrm{}_u^{R,𝒜}\mathrm{\Gamma }_{R,𝒜}(s,f)}\left(1\omega _{R,𝒜}(u)\right)+\omega _{R,𝒜}(u)`$
For $`u,v\mathrm{dom}𝒜`$ such that $`uv`$, we also compute the probability $`P_{u,v}^{R,A}(s,f)`$ that $`r.𝒜`$ changes from $`u`$ to $`v`$ in $`[s,f)`$ to be
$`P_{u,v}^{R,𝒜}(s,f)`$ $`=\mathrm{P}\left\{\tau _u^{R,𝒜}fs\right\}\omega _{R,𝒜}(v)`$
$`=\left(1e^{\mathrm{}_u^{R,𝒜}\mathrm{\Gamma }_{R,𝒜}(s,f)}\right)\omega _{R,𝒜}(v).`$
Content-independent overwrites are a special case of the Markov chain model discussed above. To apply the general model formulae when content-independent updates are present, each $`\mathrm{}_v^{R,A}`$ is multiplied by $`1\omega _{R,A}(v)`$ and $`P_{u,v}^{R,A}=\omega _{R,A}(v)/(1\omega _{R,A}(u))`$ for all $`u,v\mathrm{dom}A`$, $`uv`$.
### 3.3 Summary
In this section we have introduced a general Markov-chain model for data modification, and discussed three simplified models that allows tractable computation. Using these models, one can compute, in probabilistic terms, value histograms at time $`f`$, given a known initial set of value histograms at time $`s<f`$. Such a model could be useful in query optimization, whenever the continual gathering of statistics becomes impossible due to either heavy system loads or structural constraints (*e.g.*, federations of databases with autonomous DBMSs).
Generally speaking, computing the transition matrix for an attribute $`A`$ involves repeated multiplications of square matrices of size $`\left|\mathrm{dom}A\right|`$, resulting in a computational complexity of $`\mathrm{O}(n\left|\mathrm{dom}A\right|^\nu )`$, where $`n`$ is the number of iterations needed to compute the sum (6) to machine precision. While $`n`$ is usually small, $`\left|\mathrm{dom}A\right|`$ may be very large, as demonstrated in Section 3.2.1 and Section 3.2.2. Methods such as domain lumping would require $`\mathrm{O}(nX^\nu )`$ time, where $`X\left|\mathrm{dom}A\right|`$.<sup>7</sup><sup>7</sup>7Here, $`n`$ may also be affected by the change of domain. As for random walks and independent updates, both methods no longer require repeated matrix multiplications, but rather the computation of $`\mathrm{\Gamma }_{R,A}(s,f)`$. The complexity of calculating $`\mathrm{\Gamma }_{R,A}(s,f)`$ is similar to that for $`\mathrm{\Lambda }_R(s,f)`$ in Section 2.1.1.
## 4 Insertion model verification
It is well-known that Poisson processes model a world where data updates are independent from one another. While in databases with widely distributed access, *e.g.*, incoming e-mails, postings to newsgroups, or posting of orders from independent customers, such an independence assumption seems plausible, we still need to validate the model against real data. In this section we shall present some initial experiments as a “proof of concept.” These experiments deal only with the insertion component of the model. Further experiments, including modification and deletion operations, will be reported in future work.
Our data set is taken from postings to the DBWORLD electronic bulletin board. The data were collected over more than seven months and consists of about 750 insertions, from November 9$`^{\text{th}}`$, 2000 through May 14$`^{\text{th}}`$, 2001. Figure 2 illustrates a data set with 580 insertions during the interval \[2000/11/9:00:00:00, 2001/3/31:00:00:00). We used the Figure 2 data as a *training set*, *i.e.*, it serves as our basis for parameter estimation. Later, in order to test the model, we applied these parameters to a separate *testing set* covering the period \[2001/3/31:00:00:00, 2001/5/15:00:00:00). In the experiments described below, we tried fitting the training data with two insertion-only models, namely a homogeneous Poisson process and an RPC Poisson process (see Section 2.1). For each of these two models, we have applied two variations, either as a compound or as a non-compound model. In the experiments described below, we have used the Kolmogorov-Smirnov goodness of fit test (see for example \[16, Section 7.7\]). For completeness, we first overview the principles of this statistical test.
The Kolmogorov-Smirnov test evaluates the likelihood of a *null hypothesis* that a given sample may have been drawn from some hypothesized distribution. If the null hypothesis is true, and sample set has indeed been drawn from the hypothesized distribution, then the empirical cumulative distribution of the sample should be close to its theoretical counterpart. If the sample cumulative distribution is too far from the hypothesized distribution at any point, that suggests that the sample comes from a different distribution. Formally, suppose that the theoretical distribution is $`F(x)`$, and we have $`n`$ sample values $`x_1,\mathrm{},x_n`$ in nondecreasing order. We define an empirical cumulative distribution $`F_n(x)`$ via
$$F_n(x)=\{\begin{array}{cc}0,& \text{if }x<x_1\hfill \\ \frac{k}{n},& \text{if }x_kx<x_{k+1}\hfill \\ 1,& \text{if }x>x_n,\hfill \end{array}$$
and then compute $`D_n=sup_{k=1,\mathrm{},n}\{|F_n(x_k)F(x_k)|\}`$. For large $`n`$, given a significance level $`\alpha `$, the test measures $`D_n`$ against $`X(\alpha )/\sqrt{n}`$, where $`X(\alpha )`$ is a factor depending on the *significance level* $`\alpha `$ at which we reject the null hypothesis. For example, $`X(0.05)=1.36`$ and $`X(0.1)=1.22`$. The value of $`\alpha `$ is the probability of a “false negative,” that is, the chance that the null hypothesis might be rejected when it is actually true. Larger values of $`\alpha `$ make the test harder to pass.
### 4.1 Fitting the homogeneous Poisson process
Based on the training set, we computed the parameter for a homogeneous Poisson process by averaging the 580 interarrival times, an unbiased estimator of the Poisson process parameter. The average interarrival time was computed to be 5:15:19, and thus $`\lambda =4.57`$ per day. Figure 3(a) provides a pictorial comparison of the cumulative distribution functions of the interarrival times with their theoretical counterpart. We applied the Kolmogorov-Smirnov test to the distribution of interarrival times, comparing it with an exponential distribution with a parameter of $`\lambda =4.57`$. The outcome of the test is $`D_n=0.106`$, which means we can reject the null hypothesis at any reasonable level of confidence $`\alpha 0.005`$ (for $`\alpha =0.005`$, the rejection threshold is $`0.0718`$ for $`n=580`$). In all likelihood, then, the data are not derived from a homogeneous Poisson process.
Next, we have applied a compound homogeneous Poisson model. Our rationale in this case is that DBWORLD is a moderated list, and the moderators sometimes work on postings in batches. These batches are sometimes posted to the group in tightly-spaced clusters. For all practical purposes, we treat each such cluster as a single batch insertion event. To construct the model, any two insertions occurring within less than one minute from one another were considered to be a single event occurring at the insertion time of the first arrival. For example, on November 14, 2000, we had three arrivals, one at 13:43:19, and two more at 13:43:23. All three arrivals are considered to occur at the same insertion arrival event, with an insertion time of 13:43:19. Using the compound variation, the data set now has 557 insertion events. The revised average interarrival time is now 5:28:20, and thus $`\lambda =4.39`$ per day. Figure 3(b) provides a pictorial comparison of the cumulative distribution functions of the interarrival times, assuming a compound model, with their theoretical counterpart. We have applied the Kolmogorov-Smirnov test to the distribution of interarrival times, comparing it with an exponential distribution with a parameter of $`\lambda =4.39`$. The outcome was somewhat better than before. $`D_n=0.094`$, which means we can still reject the null hypothesis at any level of confidence $`\alpha 0.005`$ (for $`\alpha =0.005`$, the rejection threshold is $`0.0733`$ for $`n=557`$). Although the compound variant of the model fits the data better, it is still not statistically plausible.
### 4.2 Fitting the RPC Poisson process
Next, we tried fitting the data to an RPC model. Examining the data, we chose a cycle of one week. Within each week, we used the same pattern for each weekday, with one interval for work hours (9:00-18:00), plus five additional three-hour intervals for “off hours”. We treated Saturday and Sunday each as one long interval. Table 2 shows the arrival rate parameters for each segment of the RPC Poisson model, calculated in much the same manner as the for the homogeneous Poisson model.
The specific methodology for structuring the RPC Poisson model is beyond the scope of this paper and can range from *ad hoc* “look and feel” crafting (as practiced here) to more established formal processes for statistically segmenting, filtering, and aggregating intervals . It is worth noting, however, that from experimenting with different methods, we have found that the model is not sensitive to slight changes in the interval definitions. Also, the model we selected has only $`8`$ segments, and thus only $`8`$ parameters, so there is little danger of “overfitting” the training data set, which has over $`500`$ observations.
Next, we attempted to statistically validate the RPC model. To this end, we use the following lemma:
###### Lemma 6
Given a nonhomogeneous Poisson process with arrival intensity $`\lambda (t)`$, the random variable $`U_s=_s^{s+L_{R,s}}\lambda (t)𝑑t`$ is of the distribution $`\mathrm{Exp}(1)`$.
Proof. Let $`f_s(t)=\mathrm{\Lambda }(s,s+t)`$, which is a monotonically nondecreasing function. From Lemma 1, $`\mathrm{P}\{L_{R,s}<t\}=1e^{f_s(t)}`$ for all $`t0`$. We have $`U_s=f_s(L_{R,s})`$. By applying the monotonic function $`f_s`$ to both sides of the inequality $`L_{R,s}<t`$, one has that $`\mathrm{P}\{f_s(L_{R,s})<f_s(t)\}=\mathrm{P}\{L_{R,s}<t\}=1e^{f_s(t)}`$ for all $`t0`$. Substituting in the definitions of $`U_s`$ and $`u=f_s(t)`$, one then obtains $`\mathrm{P}\{U{}_{s}{}^{}<u\}=1e^u`$ for all $`u0`$, and therefore $`U_s\mathrm{Exp}(1)`$.
Thus, given an instantaneous arrival rate $`\lambda (t)`$, and a sequence of observed arrival events $`\{t_n\}_{n=0}^N`$, we compute the set of values $`u_n=_{t_{n1}}^{t_n}\lambda (t)𝑑t`$, $`n=1,\mathrm{},N,`$ and perform a Kolmogorov-Smirnov test of them versus the unit exponential distribution.
Figure 4(a) provides a comparison of the theoretical and empirical cumulative distribution of the random variable $`U`$. We applied the Kolmogorov-Smirnov test to $`U`$, comparing it with an exponential distribution with $`\lambda =1`$, based on Lemma 6. The outcome of the test is $`D_n=0.080`$, which is better than either homogeneous model, but is still rejected at any reasonable level of significance (recall that for $`\alpha =0.005`$, the rejection threshold is again $`0.0718`$ for $`n=580`$).
Finally, we evaluated a compound version of the RPC model, combining successive postings separated by less than one minute. We kept the same segmentation as in Table 2, but recalculated the arrival intensities in each segment, as shown in Table 3.
Next we recalculated the sample of the random variable $`U`$ for the compound RPC Poisson model, and applied the Kolmogorov-Smirnov test. In this case, we have $`D_n=0.050`$, which cannot be rejected at any reasonable confidence level through $`\alpha =0.10`$ (for $`\alpha =0.10`$, the rejection threshold is $`0.0517`$ for $`n=557`$). Figure 4(b) shows the theoretical and empirical distributions of $`U`$ in this case.
As a final confirmation of the applicability of the compound RPC Poisson model, we attempted to validate the assumption that the number of postings in successive insertion events are independent and identically distributed (IID). In the sample, 536 insertion events were of size 1, 19 were of size 2, and 2 were of size 3. Thus, we approximate the random variable $`\mathrm{\Delta }_R^+`$ as having a $`536/557.962`$ probability of being 1, a $`19/557.034`$ probability of being 2, and a $`2/557.004`$ probability of being 3. Validating that the observed insertion batch sizes $`\mathrm{\Delta }_{R,i}^+`$ appear to be independently drawn from this distribution is somewhat delicate, since they nearly always take the value 1. To compensate, we performed our test on the *runs* in the sample, that is, the number of consecutive insertion events of size 1 between insertions of size 2 or 3. Our sample contains 21 runs, ranging from 0 to 112. If the insertion batch sizes $`\{\mathrm{\Delta }_{R,i}^+\}`$ are independent with the distribution $`\mathrm{\Delta }_R^+`$, then the length of a run should be a geometric random variable with parameter $`536/557.962`$. We tested this hypothesis via a Kolmogorov-Smirnov test, as shown in Figure 5. The $`D_n`$ statistic is $`0.207`$, which is within the $`\alpha =0.1`$ acceptance level for a sample of size $`n=21`$ (although the divergence of the theoretical and empirical curves in Figure 5 is more visually pronounced than in the prior figures, it should be remembered that the sample is far smaller). Thus, the assumption that the insertion batch sizes $`\{\mathrm{\Delta }_{R,i}^+\}`$ are IID is plausible.
Table 4 compares the goodness-of-fit of the four models to the test data. For each of the models, we have specified the KS test result ($`D_n`$) and the level at which one can reject the null hypothesis. The higher the level of confidence is, the better the fit is. The RPC compound Poisson model models best the data set, accepting the null hypothesis at any level up to $`0.1`$ (which practically means that the model can fit to the data well). The main conclusion from these experiments is that the simple model of homogeneous Poisson process is limited to the modeling of a restricted class of applications (one of which was suggested in ). Therefore, there is a need for a more elaborate model, as suggested in this paper, to capture a broader range of update behaviors. A nonhomogeneous model consisting of just 8 segments per week, as we have constructed, seems to model the arrivals significantly better than the homogeneous approach.
## 5 Content evolution cost model
We now develop a cost model suitable for transcription-scheduling applications such as those described in Example 2. The question is how often to generate a remote replica of a relation $`R`$. We have suggested one such policy in Example 6. In this section, we shall introduce two more policies and show an empirical comparison based on the data introduced in Section 4.
A transcription policy aims to minimize the combined cost of *transcription cost* and *obsolescence cost* . The former includes the cost of connecting to a network and the cost of transcribing the data, and may depend on the time at which the transcription is performed (*e.g.*, as a function of network congestion), and the length of connection needed to perform the transcription. The obsolescence cost captures the cost of using obsolescent data, and is basically a function of the amount of time that has passed since the last transcription.
In what follows, let the set $`\{b_i,e_i\}_{i=1}^{\mathrm{}}`$ represents an infinite sequence of connectivity periods between a client and a server. During session $`i`$, the client data is synchronized with the state of the server at time $`b_i`$, the information becoming available at the client at time $`e_i`$. At the next session, beginning at time $`b_{i+1}`$, the client is updated with all the information arriving at the server during the interval $`(b_i,b_{i+1}]`$, which becomes usable at time $`e_{i+1}`$, and so forth. We define $`b_0=e_0=0`$, and require that $`0<b_1e_1<b_2e_2<\mathrm{}`$.
Let $`C_{R,\text{u}}(s,f)`$ denote the cost of performing a transcription of $`R`$ starting at time $`f`$, given that the last update was started at time $`s`$. Let $`C_{R,\text{o}}(s,f)`$, to be described in more detail later, denote the obsolescence cost through time $`f`$ attributable to tuples inserted into $`R`$ at the server during the time interval $`(s,f]`$. Then the total cost $`C_R(t)`$ through time $`t`$ is
$$C_R(t)=\underset{i:b_it}{}\left(\alpha C_{R,\text{u}}(b_{i1},b_i)+(1\alpha )C_{R,\text{o}}(b_{i1},b_i)\right)+(1\alpha )C_{R,\text{o}}(b_{i^{}(t)},t),$$
(10)
where $`i^{}(t)=\mathrm{max}\left\{i|b_it\right\}`$ and $`\alpha `$ serves as the ratio of importance a user puts on the transcription cost versus the obsolescence cost. Traditionally, $`\alpha =0`$, and therefore $`C_R(t)`$ is minimized for $`C_{R,\text{o}}(b_{i1},b_i)=0`$, $`b_i<t`$, allowing the use of current data only. In this section we shall look into another, more realistic approach, where data currency is sacrificed (up to a level defined by the user through $`\alpha `$) for the sake of reducing the transcription cost. Ideally, one would want to choose the sequence $`\{b_i,e_i\}_{i=1}^{\mathrm{}}`$ of connectivity periods, subject to any constraints on their durations $`e_ib_i`$, to minimize $`C_R(t)`$ over some time horizon $`t`$. One may also consider the asymptotic problem of minimizing the average cost over time, $`lim_t\mathrm{}C_R(t)/t`$. We note that the presence of $`\alpha `$ is not strictly required, as its effects could be subsumed into the definitions of the $`C_{R,\text{u}}`$ and $`C_{R,\text{o}}`$ functions, especially if both are expressed in natural monetary units. However, we retain $`\alpha `$ in order to demonstrate some of the parametric properties of our model.
In general, modeling transcription and obsolescence costs may be difficult and application-dependent. They may be difficult to quantify and difficult to convert to a common set of units, such as dollars or seconds. Some subjective estimation may be needed, especially for the obsolescence costs. However, we maintain that, rather than avoiding the subject altogether, it is best to try construct these cost models and then use them, perhaps parametrically, to evaluate transcription policies. Any transcription policy implicitly makes some trade-off between consuming network resources and incurring obsolescence, so it is best to try quantify the trade-off and see if a better policy exists. In particular, one should try to avoid policies that are clearly *dominated*, meaning that there is another policy with the same or lower transcription cost, and strictly lower obsolescence, or *vice versa*. Below, for purposes of illustration, we will give one simple, plausible way in which the cost functions may be constructed; alternatives are left to future research.
### 5.1 Transcription costing example
In determining the transcription cost, one may use existing research into costs of distributed query execution strategies. Typically, (*e.g.*, ) the transcription time can be computed as some function of the CPU and I/O time for writing the new tuples onto the client and the cost of transmitting the tuples over a network. There is also some fixed setup time to establish the connection, which can be substantial. For purposes of example, suppose that
$`C_{R,\text{u}}(s,f)`$ $`=c+\beta \left(X_R(s,f)+Y_R^+(s,f)+\left|R(s)\right|Y_R(s,f)\right)`$
$`=c+\beta \left(X_R(s,f)+\left|R(s)\right|Y_R^{}(s,f)\right)`$
Here, $`c0`$ denotes the fixed setup cost, $`\beta 0`$, $`X_R(s,f)`$ denotes the number of tuples inserted during the interval $`(s,f]`$ that survive through time $`f`$, $`Y_R^+(s,f)`$ is the number of tuples that survive but are modified, by time $`f`$, and $`\left|R(s)\right|Y_R(s,f)`$ is the number of deleted tuples. For the latter, it may suffice to transmit only the primary key of each deleted tuple, incurring a unit cost of less than $`\beta `$. For sake of simplicity, however, we use the same cost factor $`\beta `$ for deletion, insertion, and modification. We note that, under this assumption,
$$\underset{i:b_it}{}C_{R,\text{u}}(b_{i1},b_i)=n(t)c+\beta \left|R(s)\right|+\beta \underset{i:b_it}{}\left(X_R(b_{i1},b_i)Y_R^{}(b_{i1},b_i)\right),$$
where $`n(t)`$ is the number of transcriptions in the interval $`[0,t]`$. For the special case that there are no deletions or modifications, $`\beta \left|R(s)\right|+\beta \left(X_R(s,f)Y_R^{}(s,f)\right)=\beta B(s,f)`$ and
$$\underset{i:b_it}{}C_{R,\text{u}}(b_{i1},b_i)=n(T)c+\beta B(0,b_{i^{}(T)}).$$
For large $`t`$, one would expect the $`\beta B(0,b_{i^{}(t)})`$ term to be roughly comparable across most reasonable polices, whereas the $`n(t)c`$ term may vary widely for any value of $`t`$. It is worth noting that $`c`$ and $`\beta `$ could be generalized to vary with time or other factors. For example, due to network congestion, certain times of day may have higher unit transcription costs than others. Also, transcribing via airline-seat telephone costs substantially more than connecting via a cellular phone. For simplicity, we have refrained from discussing such variations in the transcription cost.
### 5.2 Obsolescence costing example
We next turn our attention to the obsolescence cost, which is clearly a function of the update time of tuples and the time they were transcribed to the client. Intuitively, the shorter the time between the update of a tuple and its transcription to the client, the better off the client would be. As a basis for the obsolescence cost, we suggest a criterion that takes into account user preferences, as well as the content evolution parameters. For any relation $`R`$, times $`s<f`$, and tuple $`rR(s)R(f)`$, let $`b(r)`$ and $`d(r)`$ denote the time $`r`$ was inserted into and deleted from $`R`$, respectively. We let $`\iota _r(s,f)`$ be some function denoting the contribution of tuple $`r`$ to the obsolescence cost over $`(s,f]`$; we will give some more specific example forms of this function later. We then make the following definition:
###### Definition 3
The total *obsolescence cost* of a relation $`R`$ over the time interval $`(s,f]`$ (annotated *$`C_{R,\text{o}}(s,f)`$*) is defined to be *$`C_{R,\text{o}}(s,f)_{rR(s)R(f)}\iota _r(s,f).`$*$`\mathrm{}`$
Our principal concern is with the *expected* obsolescence cost, that is, the expected value of $`C_{R,\text{o}}(s,f)`$,
$$\mathrm{E}\left[C_{R,\text{o}}(s,f)\right]=\mathrm{E}\left[\underset{rR(s)R(f)}{}\iota _r(s,f)\right].$$
To compute $`\mathrm{E}\left[C_{R,\text{o}}(s,f)\right]`$, we note that
$$\mathrm{E}\left[C_{R,\text{o}}(s,f)\right]=\mathrm{E}\left[\underset{rR(s)R(f)}{}\iota _r(s,f)\right]+\mathrm{E}\left[\underset{rR(s)\backslash R(f)}{}\iota _r(s,f)\right]+\mathrm{E}\left[\underset{rR(f)\backslash R(s)}{}\iota _r(s,f)\right].$$
The three terms in the last expression represent potentially modified tuples, deleted tuples, and inserted tuples, respectively. We denote these three terms by $`\widehat{\iota }_R^\mathrm{M}(s,f)`$, $`\widehat{\iota }_R^\mathrm{D}(s,f)`$, and $`\widehat{\iota }_R^\mathrm{I}(s,f)`$, respectively, whence
$$\mathrm{E}\left[C_{R,\text{o}}(s,f)\right]=\widehat{\iota }_R^\mathrm{M}(s,f)+\widehat{\iota }_R^\mathrm{D}(s,f)+\widehat{\iota }_R^\mathrm{I}(s,f).$$
### 5.3 Obsolescence for insertions
We will now consider a specific metric for computing the obsolescence stemming from insertions in $`(s,f]`$, as follows:
$$\iota _r^\mathrm{I}(s,f)=\{\begin{array}{cc}g^\mathrm{I}(s,f,b(r))\hfill & s<b(r)f<d(r)\hfill \\ 0\hfill & \text{otherwise},\hfill \end{array}$$
(11)
where $`g^\mathrm{I}(s,f,t)`$ is some application-dependent function representing the level of importance a user assigns, over the interval $`(s,f]`$, to a tuple arriving at a time $`t`$. For example, in an e-mail transcription application, a user may attach greater importance to messages arriving during official work hours, and a lesser measure of importance to non-work hours (since no one expects her to be available at those times). Thus, one might define
$$g^\mathrm{I}(s,f,t)=_t^fa(\tau )𝑑\tau ,\text{where }a(\tau )=\{\begin{array}{cc}a_1,\hfill & \text{if }\tau \text{ is during work hours}\hfill \\ a_2,\hfill & \text{if }\tau \text{ is after hours,}\hfill \end{array}$$
(12)
and $`a_1a_2`$. For $`a_1=a_2=1`$, $`g^\mathrm{I}(s,f,t)`$ takes a form resembling the age of a local element in . More complex forms of $`g^\mathrm{I}(s,f,t)`$ are certainly possible. In this simple case, we refer to $`a_1/a_2`$ as the *preference ratio*.
Using the properties of nonhomogeneous Poisson processes, we calculate
$`\widehat{\iota }_R^\mathrm{I}(s,f)`$ $`=\mathrm{E}\left[{\displaystyle \underset{rR(f)\backslash R(s)}{}}\iota _r(s,f)\right]`$
$`=\mathrm{E}\left[X_R(s,f)\right]\mathrm{E}\left[f(s,f,b(r))|s<b(r)f<d(r)\right]`$
$`=\stackrel{~}{\mathrm{\Lambda }}_R(s,f)\mathrm{E}\left[\mathrm{\Delta }_R^+\right]{\displaystyle _s^f}{\displaystyle \frac{\lambda _R(t^{})}{\stackrel{~}{\mathrm{\Lambda }}_R(s,f)}}g^\mathrm{I}(s,f,t^{})𝑑t^{}`$
$`=\mathrm{E}\left[\mathrm{\Delta }_R^+\right]{\displaystyle _s^f}\lambda _R(t^{})g^\mathrm{I}(s,f,t^{})𝑑t^{}.`$
###### Example 10 (Transcription policies using the expected obsolescence cost)
Consider the insertion-only data set of Section 4. Figure 6 compares two transcription policies for the week $`[2001/4/2:0:00,2001/4/8:0:00)`$. The transcription policy in Figure 6(a) (referred to below as the *uniform synchronization point* — USP — policy) was suggested in . According to this policy, the intervals $`(s,f]`$ are always of the same size. The decision regarding the interval size $`fs`$ may be either arbitrary (*e.g.*, once a day) or may depend on $`\lambda `$, the Poisson model parameter (in which case a homogeneous Poisson process is implicitly assumed). The policy may be expressed as $`f=s+M/\lambda `$ for some multiplier $`M>0`$. According to this policy with $`M=1`$ (as suggested in ), and $`\lambda =4.57`$ per day as computed from the training data. Therefore, one would refresh the database every 5:15:19. Figure 6(a) shows the transcription times resulting from the USP policy.
Consider now another transcription policy, dubbed the *threshold* policy. With this policy, given that the last connection started at time $`s`$, we transcribe at time $`f`$ if the expected obsolescence cost from insertions ($`\widehat{\iota }_R^\mathrm{I}(s,f)`$) exceeds $`\mathrm{\Pi }`$, where $`\mathrm{\Pi }`$ is a threshold that measures the user’s tolerance to obsolescent data. In comparing the two policies, one can compute $`\mathrm{\Pi }`$, given $`M`$, as follows. Consider the homogeneous case where $`\mathrm{E}\left[\mathrm{\Delta }_R^+\right]=1`$ and $`\lambda _R(t)=\lambda _R`$ for all $`t`$. Assume further that $`a_1=a_2=1`$ for all $`t`$. In this case,
$$\widehat{\iota }_R^\mathrm{I}(s,f)=_s^f\lambda _R(ft)𝑑t=\lambda _R_s^f(ft)𝑑t=\frac{1}{2}\lambda _R(fs)^2$$
Setting $`f=s+M/\lambda _R`$ and $`\mathrm{\Pi }=\widehat{\iota }_R^\mathrm{I}(s,f)`$, one has that
$$\mathrm{\Pi }=(1/2)\lambda _R(fs)^2=(1/2)\lambda _R\left(M/\lambda _R\right)^2=M^2/2\lambda _R.$$
Figure 6(b) shows the transcription times using the RPC arrival model (see Section 4.1) and the threshold policy with $`\mathrm{\Pi }=0.109`$ (obtained by setting $`M=1`$ and $`\lambda =4.57`$ per day, and letting $`\mathrm{\Pi }=M^2/2\lambda `$). It is worth noting that transcriptions are more frequent when the $`\lambda `$ intensity is higher and less frequent whenever the arrival rate is expected to be more sluggish.
We have performed experiments comparing the performance of the threshold policy for the homogeneous Poisson model (equivalent to the USP policy) and the RPC Poisson model. Figure 7 shows representative results, with costs computed over the testing set. Figure 7(a) displays the obsolescence cost and the number of transcriptions for various $`M`$ values, with a preference ratio $`a_2/a_1=4`$. For all $`M`$ values, there is no dominant model. For example, for $`M=1`$, the RPC model has a slightly higher obsolescence cost (43.02 versus 42.35, a 1.6% increase) and a significantly lower number of transcriptions (137 versus 204, a 32.8% decrease).
Figure 7(b) provides a comparison of combined normalized obsolescence and transcription costs for both insertion models and $`\alpha \{0.6,0.7,0.8\}`$ (still assuming a 4:1 preference ratio). Solid lines represent results related with the homogeneous Poisson model, while dotted lines represent results related with the RPC Poisson model. Generally speaking, the RPC model performs better for small $`M`$ values ($`M7`$), while the homogeneous model performs better for the largest $`M`$ values ($`M8`$). $`\mathrm{}`$
###### Example 11 (Comparison of USP, threshold and FA policies)
Once again with the data from Section 4, we consider one more transcription policy, the first alteration (FA) policy derived from the analysis of Example 6. Since there are no deletions or modifications, $`Z_R(s,f)`$ simplifies to $`\mathrm{\Lambda }_R(s,f)`$. We choose $`\pi `$ in the FA policy to be a function of $`M`$ such that the transcription intervals agree with the USP policy in the case of the homogeneous model. Figure 9 compares the performance of all three transcription policies: USP, threshold, and FA, for a 4:1 preference ratio and $`\alpha =0.8`$, using the testing data set to compute the costs. For $`M=1`$, the threshold policy and the FA policy perform similarly, where the FA policy performs slightly better than the Threshold policy. Both policies outperform the USP policy. The threshold policy is best for $`M\{2\mathrm{}8\}`$. For all $`M>8`$, the USP policy is best. The best policy for this choice of $`a_1/a_2`$ and $`\alpha `$ is threshold with $`M=6`$, followed closely by FA with $`M\{5,6\}`$. We have conducted our experiments with various $`\alpha `$ values and our conclusion is that the Threshold model is preferred over the USP model for larger $`\alpha `$, that is, the more the user is willing to sacrifice currency for the sake of reducing transcription cost.$`\mathrm{}`$
### 5.4 Obsolescence for deletions
In a similar manner to Section 5.3, we will consider the following metric for computing the obsolescence stemming from deletions in $`(s,f]`$. We compute $`\iota _r^\mathrm{D}(s,f)`$ via
$$\iota _r^\mathrm{D}(s,f)=\{\begin{array}{cc}g^\mathrm{D}(s,f,d(r))\hfill & b(r)s<d(r)f\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
(13)
where $`g^\mathrm{D}(s,f,t)`$ is some application-dependent function, possibly similar to $`g^\mathrm{I}(s,f,t)`$ above.
Using the properties of nonhomogeneous Poisson processes, we calculate
$`\widehat{\iota }_R^\mathrm{D}(s,f)`$ $`=\mathrm{E}\left[{\displaystyle \underset{rR(s)\backslash R(f)}{}}\iota _r(s,f)\right]`$
$`=\left(\left|R(s)\right|\mathrm{E}\left[Y_R(s,f)\right]\right)\mathrm{E}\left[g^\mathrm{D}(s,f,d(r))|b(r)s<d(r)f\right]`$
$`=\left(\left|R(s)\right|p_R(s,f)\left|R(s)\right|\right)\mathrm{E}\left[g^\mathrm{D}(s,f,d(r))|b(r)s<d(r)f\right]`$
$`=\left|R(s)\right|\left(1p_R(s,f)\right)\mathrm{E}\left[g^\mathrm{D}(s,f,d(r))|b(r)s<d(r)f\right]`$
In the case $`R,S`$ has fixed multiplicity for all $`SS(R)`$, $`p_R(s,f)=\mathrm{exp}(\stackrel{~}{M}_R(s,f))`$, where $`\stackrel{~}{M}_R(s,f)=_s^f\stackrel{~}{\mu }_R(t)𝑑t`$ and $`\stackrel{~}{\mu }_R(t)=_{SS(R)}w(R,S)\mu _S(t)`$. Therefore,
$`\widehat{\iota }_{R,A}^\mathrm{D}(s,f)`$ $`=\left|R(s)\right|\left(1\mathrm{exp}(\stackrel{~}{M}_R(s,f))\right){\displaystyle _s^f}{\displaystyle \frac{\stackrel{~}{\mu }_R(t^{})}{\stackrel{~}{M}_R(s,f)}}g^\mathrm{D}(s,f,t^{})𝑑t^{}`$
$`=\left|R(s)\right|\left({\displaystyle \frac{1\mathrm{exp}(\stackrel{~}{M}_R(s,f))}{\stackrel{~}{M}_R(s,f)}}\right){\displaystyle _s^f}\stackrel{~}{\mu }_R(t^{})g^\mathrm{D}(s,f,t^{})𝑑t^{}`$
### 5.5 Obsolescence for modification
We now consider obsolescence costs relating to modifications. While, in some applications, a user may be primarily concerned with how many tuples were modified during $`[s,f)`$, we believe that a more general, attribute-based framework is warranted here, taking into account exactly how each tuple was changed. Therefore, we define $`\iota _{r,A}(s,f)`$ to be some function denoting the contribution of attribute $`A𝒜(R)`$ in tuple $`r`$ to the obsolescence cost over $`(s,f]`$ and assume that
$$\iota _r(s,f)=\underset{A𝒜(R)}{}\iota _{r,A}(s,f)$$
Therefore,
$`\widehat{\iota }_R^\mathrm{M}(s,f)`$ $`=\mathrm{E}\left[{\displaystyle \underset{rR(s)R(f)}{}}\iota _r(s,f)\right]`$
$`=\mathrm{E}\left[{\displaystyle \underset{A𝒜(R)}{}}{\displaystyle \underset{rR(s)R(f)}{}}\iota _{r,A}(s,f)\right]`$
$`={\displaystyle \underset{A𝒜(R)}{}}\widehat{\iota }_{R,A}^\mathrm{M}(s,f)`$
where $`\widehat{\iota }_{R,A}^\mathrm{M}(s,f)`$ is the expected obsolescence cost due to modifications to $`A`$ during $`(s,f]`$. Assuming that attributes not in $`𝒞(R)`$ incur zero modification cost, the last sum may be taken over $`𝒞(R)`$ instead of $`𝒜(R)`$.
We start the section by introducing the notion of distance metric and provide two models of $`\iota _{r,A}(s,f)`$, for numeric and non-numeric domains. We then provide an explicit description of $`\widehat{\iota }_{R,A}^\mathrm{M}`$, based on distance metrics.
#### 5.5.1 General distance metrics
Let $`c_{u,v}^{R,A}`$, where $`u,v\mathrm{dom}A`$ denote the elements of a matrix of costs for an attribute $`A`$. We declare that if $`r.A(s)=u`$ and $`r.A(f)=v`$, then $`\iota _{r,A}(s,f)=c_{u,v}^{R,A}`$, or equivalently,
$$\iota _{r,A}(s,f)=c_{r.A(s),r.A(f)}^{R,A}.$$
Consequently, we require that $`c_{u,u}^{R,A}=0`$ for all $`u\mathrm{dom}A`$, so that an unchanged attribute field yields a cost of zero.
##### A squared-error metric for numeric domains:
For numeric domains, that is, $`A𝒩`$, we propose a squared-error metric, as is standard in statistical regression models. In this case, we let
$$\iota _{r,A}(s,f)=c_{r.A(s),r.A(f)}^{R,A}=k_{R,A}(s)(r.A(f)r.A(s))^2,$$
where $`k_{R,A}(s)`$ is a user-specified scaling factor. A typical choice for the scaling factor would be the reciprocal $`1/\left(\mathrm{Var}_{rR(s)}[r.A(s)]\right)`$ of the variance of attribute $`A`$ in $`R`$ at time $`s`$,
$`\mathrm{Var}_{rR(s)}[r.A(s)]`$ $`=\mathrm{E}_{rR(s)}\left[(r.A(s)\mathrm{E}_{rR(s)}[r.A(s)])^2\right]`$
$`=\mathrm{E}_{rR(s)}[r.A(s)^2]\mathrm{E}_{rR(s)}[r.A(s)]^2`$
$`={\displaystyle \frac{1}{\left|R(s)\right|}}\left({\displaystyle \underset{v\mathrm{dom}A}{}}v^2\widehat{R}_{A,v}(s)\right)\left({\displaystyle \frac{1}{\left|R(s)\right|}}{\displaystyle \underset{v\mathrm{dom}A}{}}v\widehat{R}_{A,v}(s)\right)^2.`$
Other choices for the scaling factor $`k_{R,A}(s)`$ are also possible. In any case, we may calculate the expected alteration cost for attribute $`A`$ in tuple $`r`$ via
$`\mathrm{E}\left[\iota _{r,A}(s,f)\right]`$ $`=\mathrm{E}\left[k_{R,A}(s)(r.A(f)r.A(s))^2\right]`$
$`=k_{R,A}(s)\mathrm{E}[r.A(f)^22r.A(f)r.A(s)+r.A(s)^2]`$
$`=k_{R,A}(s)(\mathrm{E}[r.A(f)^2]2r.A(s)\mathrm{E}[r.A(f)]+r.A(s)^2).`$ (14)
##### A general metric for non-numeric domains:
For non-numeric domains, it may not be possible or meaningful to compute the difference of $`r.A(s)`$ and $`r.A(f)`$. In such cases, we shall use a general cost matrix $`[`$$`c_{u,v}^{R,A}`$$`]_{u,v\mathrm{dom}A}`$ and compute
$`\mathrm{E}\left[\iota _{r,A}(s,f)\right]`$ $`={\displaystyle \underset{v\mathrm{dom}A}{}}\left(P_{r.A(s),v}^{R,A}(s,f)\right)\left(c_{r.A(s),v}^{R,A}\right)`$
$`={\displaystyle \underset{{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{v\mathrm{dom}A}{vr.A(s)}}}{}}\left(P_{r.A(s),v}^{R,A}(s,f)\right)\left(c_{r.A(s),v}^{R,A}\right).`$
For domains that have no particular structure, a typical choice might be $`c_{u,v}^{R,A}=1`$ whenever $`uv`$. In this case, the expected cost calculation simplifies to
$`\mathrm{E}\left[\iota _{r,A}(s,f)\right]`$ $`=\mathrm{P}\{r.A(f)r.A(s)\}`$
$`=1P_{r.A(s),r.A(s)}^{A,R}(s,f).`$
We are now ready to consider the calculation of $`\widehat{\iota }_{R,A}^\mathrm{M}(s,f)`$.
#### 5.5.2 The expected modification cost
We next consider computing the expected modification cost $`\widehat{\iota }_{R,A}^\mathrm{M}(s,f)`$. To do so, we partition the tuples $`r`$ in $`R(s)R(f)`$ according to their initial value $`r.A(s)`$ of the attribute $`A`$. Consider the subset $`R_{A,u}(s)R(f)`$ of all $`rR(s)R(f)`$ that have $`r.A(s)=u`$. Since all such tuples are indistinguishable from the point of view of the modification process for $`(R,A)`$, their $`\iota _{r,A}(s,f)`$ random variables will be identically distributed. The number of tuples $`rR(s)`$ with $`r.A(s)=u`$ is, by definition, $`\widehat{R}_{A,u}(s)`$. The number $`\left|R_{A,u}(s)R(f)\right|`$ that are also in $`r.A(f)`$ is a random variable whose expectation, by the independence of the deletion and modification processes, must be $`p_R(s,f)\widehat{R}_{A,u}(s)`$. Using standard results for sums of random numbers of IID random variables, we conclude that
$`\widehat{\iota }_{R,A}^\mathrm{M}(s,f)`$ $`=\mathrm{E}\left[{\displaystyle \underset{rR(s)R(f)}{}}\iota _{r,A}(s,f)\right]`$
$`={\displaystyle \underset{u\mathrm{dom}A}{}}\left(p_R(s,f)\widehat{R}_{A,u}(s)\right)\mathrm{E}\left[\iota _{r,A}(s,f)\right|r.A(s)=u]`$
$`=p_R(s,f){\displaystyle \underset{{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{u\mathrm{dom}A}{\widehat{R}_{A,u}(s)>0}}}{}}\widehat{R}_{A,u}(s)\mathrm{E}\left[\iota _{r,A}(s,f)\right|r.A(s)=u]`$
$`=p_R(s,f){\displaystyle \underset{{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{u\mathrm{dom}A}{\widehat{R}_{A,u}(s)>0}}}{}}\widehat{R}_{A,u}(s)\widehat{\iota }_{R,A,u}^\mathrm{M}(s,f),`$
where we define $`\widehat{\iota }_{R,A,u}^\mathrm{M}(s,f)=\mathrm{E}\left[\iota _{r,A}(s,f)\right|r.A(s)=u]`$. We now address the calculation of the $`\widehat{\iota }_{R,A,u}^\mathrm{M}(s,f)`$.
For a non-numeric domain, we have from Section 5.5.1 that
$$\widehat{\iota }_{R,A,u}^\mathrm{M}(s,f)=\underset{v\mathrm{dom}A}{}\left(P_{u,v}^{R,A}(s,f)\right)\left(c_{u,v}^{R,A}\right),$$
and in the simple case of $`c_{u,v}^{R,A}=1`$ whenever $`uv`$,
$$\widehat{\iota }_{R,A,u}^\mathrm{M}(s,f)=1P_{u,u}^{R,A}(s,f).$$
In any case, $`P_{u,v}^{R,A}(s,f)`$ and $`P_{u,u}^{R,A}(s,f)`$ may be computed using the results of Section 3.
For a numeric domain, we have from (14) that
$`\widehat{\iota }_{R,A,u}^\mathrm{M}(s,f)`$ $`=k_{R,A}(s)(\mathrm{E}\left[(r.A(f))^2\right|r.A(s)=u]2u\mathrm{E}[r.A(f)|r.A(s)=u]+u^2)`$
$`=k_{R,A}(s)\left(\left({\displaystyle \underset{v\mathrm{dom}A}{}}(v^22uv)P_{u,v}^{R,A}(s,f)\right)+u^2\right).`$
In cases where a random walk approximation applies, however, the situation simplifies considerably, as demonstrated in the following proposition.
###### Proposition 7
When a random walk model with mean $`\delta `$ and variance $`\sigma ^2`$ accurately describes modifications to a numeric attribute $`A`$, $`\widehat{\iota }_{R,A,u}^\mathrm{M}(s,f)k_{R,A}(s)\mathrm{\Gamma }_{R,A}(s,f)\left(\sigma ^2+2\mathrm{\Gamma }_{R,A}(s,f)\delta ^2\right).`$
Proof. In this case, we note that the random variable $`r.A(f)r.A(s)`$ is identical to $`\mathrm{\Delta }A(s,f)`$ (using the notation of section 3.2.2), and is independent of $`r.A(s)`$. The number $`N`$ of modification events in $`(s,f]`$ has a Poisson distribution with mean $`\mathrm{\Gamma }_{R,A}(s,f)`$, and hence variance $`\mathrm{\Gamma }_{R,A}(s,f)`$<sup>2</sup>. Therefore we have, for any $`u\mathrm{dom}A`$,
$`\widehat{\iota }_{R,A,u}^\mathrm{M}(s,f)`$ $`k_{R,A}(s)\mathrm{E}\left[\left(\mathrm{\Delta }A(s,f)\right)^2\right]`$
$`=k_{R,A}(s)(\mathrm{Var}\left[\mathrm{\Delta }A(s,f)\right]+\mathrm{E}\left[\mathrm{\Delta }A(s,f)\right]^2)`$
$`=k_{R,A}(s)(\mathrm{E}\left[N\right]\sigma ^2+\delta ^2\mathrm{Var}\left[N\right]+\mathrm{E}\left[N\right]^2\delta ^2)`$
$`=k_{R,A}(s)\mathrm{\Gamma }_{R,A}(s,f)\left(\sigma ^2+2\mathrm{\Gamma }_{R,A}(s,f)\delta ^2\right).`$
### 5.6 Example: the use of the cost model in Web crawling
The following example concludes the introduction of the cost function. We show how, by using the cost model, one can generate an optimal transcription policy for Web crawling.
###### Example 12 (Web Monitoring)
WebSQL is a Web monitoring tool which uses a virtual database schema to query the structural properties of Web documents. The database schema consists of two relations, Document with six attributes, namely url, title, text, type, length, and modif, and Anchor with four attributes, namely base, label, href, and context. Each tuple in Anchor indicates that document base contains a link to document href. Consider the following query (taken from http://www.cs.toronto.edu/~websql/), which identifies locally reachable documents that contain some hyperlink to a compressed Postscript File:
SELECT d.url, d.modif
FROM Document d SUCH THAT ‘‘http://www.OtherDoc.html’’ -$`>`$-$`>`$* d,
Anchor a SUCH THAT base = d
WHERE filename(a.href) CONTAINS ‘‘.ps.Z’’;
(We refrain from dwelling on the language specification;he interested reader is referred to the cited Web site.) Assume that the cost of performing the query at time $`t`$ is $`_{dD(t)}\psi _d`$, where $`D(t)`$ represents the set of scanned documents and $`\psi _d`$ is a random variable representing the size of document $`d`$ in bytes. Assuming the $`\{\psi _d\}`$ are IID, the expected cost of performing the query at time $`t`$ is thus
$$\mathrm{E}\left[\underset{dD(t)}{}\psi _d\right]=\mathrm{E}[\left|D(t)\right|]\mathrm{E}[\psi ],$$
where $`\psi `$ is a generic random variable distributed like the $`\{\psi _d\}`$.
A modification to a document is identified using changes to the modif attribute of the Document relation. For brevity in what follows, we let $`R=\text{Document}`$ and $`A=\text{modif}`$. We assign the following costs to changes in $`A`$:
* $`g^\mathrm{D}(s,f,t)=0`$ for all $`s<t<f`$, that is, the user has no interest in being notified of deleted documents.
* For all $`s<t<f`$ and $`u,v\mathrm{dom}A`$, $`uv`$, $`c_{u,v}^{R,A}=g^\mathrm{I}(s,f,t)=\mathrm{E}[\psi ]`$, where $`c_{R,A}^\mathrm{M}`$ is the cost for a modified document. For all other attribute $`A^{}A`$, $`c_{u,v}^{R,A^{}}=0`$ for all $`u,v\mathrm{dom}A^{}`$.
Suppose that a query was performed at time $`s`$, scanning the set of documents $`D(s)`$, and returning the set of documents $`B(s)`$, where $`\left|B(s)\right|\left|D(s)\right|`$. A user is interested in refreshing the query result without overloading system resources, thus balancing the cost of refreshing the query results against the cost of using partial or obsolescent data. This trade-off can be captured by the following policy: refresh the query at time $`f`$, after performing it at time $`s`$ iff
$$\mathrm{E}\left[\underset{dD(f)}{}\psi _d\right]<\mathrm{E}\left[C_{R,\mathrm{o}}(s,f)\right]$$
Thus, an equivalent conditions is
$`\mathrm{E}[\left|D(f)\right|]\mathrm{E}[\psi ]`$ $`<{\displaystyle \underset{A^{}𝒜(R)}{}}\widehat{\iota }_{R,A^{}}^\mathrm{M}(s,f)+\widehat{\iota }_R^\mathrm{D}(s,f)+\widehat{\iota }_R^\mathrm{I}(s,f)`$
$`=\widehat{\iota }_{R,A}^\mathrm{M}(s,f)+\widehat{\iota }_R^\mathrm{I}(s,f),`$
or
$`\left(p_R(s,f)\left|D(s)\right|+\stackrel{~}{\mathrm{\Lambda }}_R(s,f)\mathrm{E}\left[\mathrm{\Delta }_R^+\right]\right)\mathrm{E}[\psi ]`$
$`<\left(p_R(s,f){\displaystyle \underset{u\mathrm{dom}A}{}}\widehat{R}_{A,u}(s)(1P_{u,u}^{A,R}(s,f))+\stackrel{~}{\mathrm{\Lambda }}_R(s,f)\mathrm{E}\left[\mathrm{\Delta }_R^+\right]p^\mathrm{I}\right)\mathrm{E}[\psi ],`$
where $`p^\mathrm{I}`$ is the probability of a newly-inserted document being relevant to the query. Cancelling the factor of $`\mathrm{E}[\psi ]`$, another equivalent condition is
$$p_R(s,f)\left|D(s)\right|+\stackrel{~}{\mathrm{\Lambda }}_R(s,f)\mathrm{E}\left[\mathrm{\Delta }_R^+\right]<p_R(s,f)\underset{u\mathrm{dom}A}{}\widehat{R}_{A,u}(s)(1P_{u,u}^{A,R}(s,f))+\stackrel{~}{\mathrm{\Lambda }}_R(s,f)\mathrm{E}\left[\mathrm{\Delta }_R^+\right]p^\mathrm{I},$$
which is independent of the expected document size. Further assume that $`P_{u,u}^{RA}(s,f)=P_,^{R,A}(s,f)`$ is independent of $`u`$. Then the refresh condition can be expressed as
$$p_R(s,f)D(s)+\stackrel{~}{\mathrm{\Lambda }}_R(s,f)\mathrm{E}\left[\mathrm{\Delta }_R^+\right]<p_R(s,f)|B(s)|(1P_,^{A,R}(s,f))+\stackrel{~}{\mathrm{\Lambda }}_R(s,f)\mathrm{E}\left[\mathrm{\Delta }_R^+\right]p^\mathrm{I}$$
$`\mathrm{}`$
## 6 Conclusion and topics for future research
This paper represents a first step in a new research area, the stochastic estimation of the consistency of transcribed data over time. We have also suggested one possible technique for assigning a cost to the differences between two relation extensions, including a means of computing the expected value of this cost under our stochastic model. We have discussed a number of potential applications relating managing replicas, query management, and Web crawling. We have also examined several strategies for refreshing replicas, although other strategies are certainly possible.
As an illustration of the low client-side computational demands of the insertion-only transcription application of our model, a Java-based demo, based on the transcription policies described in and in this paper, can be accessed at http://rbs.rutgers.edu:6677/. The demo compares the performance of various policies using data that exist at a backend mSQL database.
We hope to extend our work to the case where the materialized views are not simple replications, but are produced by SQL queries that involve selections, projections, natural joins, and certain types of aggregations. This work will involve a *propagation algebra* for tracing the base data changes through a series of relational operators.
This development should make it possible to apply the theory to the management of more complex queries than presented here. In particular, it will facilitate a possible approach to managing general materialized view obsolescence on a query-by-query basis, taking into account current user preferences for query accuracy and speed. The refresh rate of materialized views in a periodically-updated data source (such as a data warehouse) can be defined in terms of data obsolescence, which in turn can be stochastically estimated using our model for content evolution. In this case, we advocate a three-way cost model for query optimization , in which the query optimizer evaluates various query plans using three complementary factors, namely generation cost, transmission cost, and obsolescence cost. The first two factors take on a conventional interpretation and the obsolescence cost of a query represents a penalty for basing the query result on possibly obsolescent materialized views. A query plan using only selection from a local materialized view, for example, might have lower generation and transmission costs, but a higher obsolescence cost, than a plan fetching complete base relations from an extranet and then processing them through a series of join operations. Our model, when combined with additional techniques to propagate updates through relational operators, can be used as a basis for estimating the obsolescence cost. However, developing the propagation algebra may require some enrichment of our basic model, in particular the introduction of dependency between the deletion and modification processes.
We foresee several additional future research directions. One direction involves the design of efficient algorithms for the numerical computations required by our model. As it stands so far, the most demanding computations required are general numerical integration and the matrix exponentiation formula (6). With regard to integration, we note that, in practice, the nonhomogeneous Poisson arrival rate functions $`\lambda _R()`$, $`\mu _R()`$, and $`\gamma _{R,A}()`$ will most likely be chosen to be periodic piecewise low-order polynomials, as suggested in Section 4. In such cases, many of the integrals needed by the model could be performed in closed form within each time period.
Further calibration and verification of the models in real situations is also needed. So far, we have demonstrated that the insertion model has plausible applications, but this work needs to be extended to the deletion and modification models. Furthermore, the insertion model may need to be generalized to handle situations where there is “burstiness” or autocorrelation in the interarrival times that may require more involved techniques than simply combining very closely spaced arrivals.
Another future research direction involves applying the model to real-life settings such as managing a data warehouse. While the model is quite flexible, a methodology is still needed for structuring Markov chains and estimating the stochastic model’s parameters. Finally, in order to calibrate the cost model, the issue of measuring user tolerance for data obsolescence should be considered.
## Acknowledgments
We would like to thank Benny Avi-Itzhak, Adi Ben-Israel, David Shanno, Andrzej Ruszczynski, Ben Melamed, Zachary Stoumbos, and Bob Vanderbei for their help. Also, we thank Kumaresan Chinnusamy and Shah Mitul for their comparative research on statistics gathering methods and Connie Lu and Gunjan Modha for their assistance in designing and implementing the demo.
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# Anomalous radial expansion in central heavy-ion reactions
## I Introduction
In heavy-ion collisions, the stage of compression and heating is followed by the expansion of nuclear matter. Expansion dynamics as a collective motion of excited matter is characterized by certain space-momentum correlations and has been well ascertained by experimental data such as the collective flow. In particular, in central heavy-ion collisions with the beam energy ranging from the Fermi energy to almost 200 GeV/A, the radial flow is clearly manifested through the flattening of the transverse spectra with the particle mass and this effect, as expected, is stronger for heavier systems .
The collective expansion scenario is important also in other issues. In studying the quantum statistical correlations which describe the space-time characteristics of expanding systems, one can infer the information about the freeze-out configuration. The recent finding here is that the size parameters of an effective source are determined not only by the geometrical length scale which measures the region of homogeneity but also by the thermal length scale which is related to the region in the coordinate space from which identical particles with similar momenta may emerge . Alongside with the temperature and the freeze-out time, the thermal size of a source is determined by the velocity gradients in this source. The thermal length dominates the correlation function if the geometrical length scale is sufficiently large . Similarly, the question about the particle velocity profile at the freeze-out configuration arises when one considers the formation of light nuclei in the framework of a coalescence model .
In almost all papers devoted to this subject, the velocity profile is assumed to be linear :
$`\stackrel{}{v}={\displaystyle \frac{\dot{R}(t)}{R(t)}}\stackrel{}{r}`$ (1)
where $`\stackrel{}{v}`$ and $`\stackrel{}{r}`$ are three- or two-dimensional vectors for the spherically or cylindrically symmetric expansion, respectively. The relation (1) is prompted by an analytical scaling solution of the equations of non-relativistic hydrodynamics with the ideal gas equation of state for a slowly expanding fireball . In fact, Eq. (1) is a consequence of a regular motion governed by the continuity equation when a fluid does not influence the expansion rate. This relation is well known in cosmology where the Hubble constant : $`\dot{R}(t)/R(t)65`$ km/s/Mpc, at the right-hand side of Eq. (1), characterizes the expansion of homogeneous and isotropic galaxies .
It is important to note that Eq. (1) need not be valid in general. As was shown by Dumitru , while the longitudinal expansion of the fireball three-volume is independent of the energy density of the fluid, the transverse collective motion in the case of (3+1)-dimensional expansion may couple the expansion rate to the properties of the fluid, i.e., to the equation of state. In particular, the hydrodynamic solution for a fireball expanding in the longitudinal and transverse direction with a possible first order hadronization phase transition affects the three-volume expansion rate on the hadronization hypersurface .
Analyzing the experimental data on nuclear multifragmentation within an extended statistical microcanonical model which takes into account an interplay of the radial expansion with a non-spherical shape of the fragmenting nucleus , it turned out to be necessary to postulate a non-Hubblean velocity profile $`\stackrel{}{v}(\stackrel{}{r})`$ for fragments in the freeze-out volume :
$`\stackrel{}{v}(\stackrel{}{r})=v_0\left({\displaystyle \frac{r}{R_0}}\right)^\alpha {\displaystyle \frac{\stackrel{}{r}}{r}}`$ (2)
with $`\alpha `$ in between 1.5 and 2 , to explain the experimental charge-number dependence of the mean kinetic energy of fragments in central $`Xe+Sn`$ collisions at 50 MeV/A. The answer to the question : why the exponent $`\alpha `$ of the radial expansion differs from 1 ($`\alpha =1`$ yields the Hubble expansion), cannot be given within the statistical multifragmentation model. An application of the hydrodynamics for this kind of problem is also questionable because the nucleon density at the freeze-out configuration is low and, moreover, the dynamical processes at short time scales are not correctly described. For that reason, in this work we study both the particle velocity profile and the particle density profile in central HI collisions using the framework of a nonlocal quantum kinetic theory . To gain an insight into the dynamics of collective expansion of small fermionic systems such as the atomic nuclei, the kinetic approach which uses the quasiparticle interaction as input and takes into account consistently the two-particle correlations is probably more reliable, even though the dynamical formation of clusters is absent in this approach.
The paper is organized as follows. In Sect. II.A, the main ingredients of the nonlocal quantum kinetic approach are presented. The time evolution of central $`Ta+Au`$ collisions at 33 MeV/A and 60 MeV/A is studied in Sect. II.B by looking at the transversal and longitudinal profiles of the nucleon velocity, the nucleon density and the proton to neutron ratio. The expansion velocity profile is discussed in more details in Sect. II.C , separately for bulk and surface particles. The qualitative evolution of the radial expansion profile with the collision energy is compared in Sect. II.D with the dynamical trajectories of excited system in the temperature - particle density plane. The possible consequences of the long-range tail in the particle density on the small momenta behavior of the Bose-Einstein correlations in discussed in Sect. II.E. Finally, Sect. III summarizes main results of the paper.
## II The kinetic approach
### A The nonlocal quantum kinetic equation
The observables of interest are : the particle density $`n(r,t)`$, the current density $`J(r,t)`$ and the kinetic energy density $`E(r,t)`$, which can be expressed by the one - particle phase - space distribution function $`f(p,r,t)`$ as follows :
$`n(r,t)`$ $`=`$ $`{\displaystyle \frac{dp}{(2\pi )^3}f(p,r,t)}`$ (3)
$`J(r,t)`$ $`=`$ $`{\displaystyle \frac{dp}{(2\pi )^3}pf(p,r,t)}`$ (4)
$`E(r,t)`$ $`=`$ $`{\displaystyle \frac{dp}{(2\pi )^3}\frac{p^2}{2m}f(p,r,t)}.`$ (5)
The one - particle distribution function obeys a nonlocal Boltzmann - Uehling - Uhlenbeck (BUU) kinetic equation :
$`{\displaystyle \frac{f_1}{t}}+{\displaystyle \frac{\epsilon _1}{k}}{\displaystyle \frac{f_1}{r}}{\displaystyle \frac{\epsilon _1}{r}}{\displaystyle \frac{f_1}{k}}={\displaystyle \underset{b}{}}{\displaystyle \frac{dpdq}{(2\pi )^6}𝒫}`$ (6)
$`\times `$ $`\left[f_3f_4\left(1f_1\right)\left(1f_2\right)\left(1f_3\right)\left(1f_4\right)f_1f_2\right],`$ (7)
with the Enskog-type shifts of the arguments :
$`f_1`$ $``$ $`f(k,r,t)`$ (8)
$`f_2`$ $``$ $`f(p,r\mathrm{\Delta }_2,t)`$ (9)
$`f_3`$ $``$ $`f(kq\mathrm{\Delta }_K,r\mathrm{\Delta }_3,t\mathrm{\Delta }_t)`$ (10)
$`f_4`$ $``$ $`f(p+q\mathrm{\Delta }_K,r\mathrm{\Delta }_4,t\mathrm{\Delta }_t).`$ (11)
The arguments of the effective scattering measure $`𝒫`$ are centered in all $`\mathrm{\Delta }`$ \- shifts. The quasiparticle energy $`\epsilon `$ contains the mean field as well as the correlated self energy. The shifts or displacements are a compact form of gradient corrections and ensure that the conservation laws contain both the mean-field and the two-particle correlations. In particular, the momentum and the energy gain arises from the finite duration of collisions . All shifts in (8) are proportional to derivatives of the scattering phase shift and have been calculated for realistic nuclear potentials . When neglecting these shifts, one recovers the usual BUU scenario.
It should be noted that using the nonlocal BUU kinetic equations for the description of the proton spectra in central $`Xe+Sn`$ collisions at 50 MeV/A leads to a significant enhancement of the high energy tail and a better agreement with the experimental data than obtained using the standard BUU equations.
### B The evolution plots
The result of the nonlocal BUU scenario for the reaction $`Ta+Au`$ at 33 MeV/A can be seen in Figs. 1 and 2. Let us concentrate first on the corresponding velocity and density profiles (the first and the second column in Fig. 1) and the arrows characterizing the mass momentum (the first column in Fig. 2). One sees that at around 40 fm/c the nuclei start to squeeze out the matter side-wards (the first column in Fig. 2), what is characterized by the momentum focusing at both sides perpendicular to the beam direction , predominantly in peripheral regions. The surface matter is stopped and bounced back in longitudinal direction during the times $`2080`$ fm/c (see the first column in Fig. 1). The inner (bulk) matter exhibiting a quite clear spatial boundary (see the second column in Fig. 2) is still moving inwards. This leads at short time scales ($`40`$ fm/c) to an enhancement of matter density. The strong velocity gradient, which is seen at $`60`$ fm/c, disappears at about $`100`$ fm/c and the inner matter comes to a rest. The recoil of the splashing matter at the surface and the attractive mean field force start effectively to reaccelerate an inner matter towards the center of mass. Since the surface particles are still accelerating in the outward direction, therefore there is a zone of matter in between which comes to a rest. This evolution leads to the dumb-bell structure in the transversal density profile at $`120240`$ fm/c.
The different behavior of the surface matter and the bulk matter leads to the development of a nonlinear velocity profile in the surface region, which can be seen in the log - log representation of the angular averaged velocity in the second column of Fig. 2. For $`t>80`$ fm/c, the velocity - radius scaling with an asymptotically stable coefficient $`\alpha _{surf}1.75\pm 0.05`$ (see the second column in Fig. 2) appears definitely in the surface region. In this region, the particle density drops nearly as a power law : $`n(r)r^\beta `$, with the asymptotically stable (for $`t>120`$ fm/c) coefficient : $`\beta 3\pm 0.2`$. As can be seen from the $`N/Z`$ \- ratios in the third column in Fig. 1, the Coulomb interaction expels protons from the surface in the early stage of the evolution. In particular, the proton rms radius is larger than the neutron one. This effect becomes weaker at later times ($`t200`$ fm/c) but, nevertheless, it survives indicating that proton and neutron distributions are different in the surface region. This may lead to different source temperatures and temperature gradients for protons and neutrons and, hence, to different quantum statistical corrections for protons and neutrons in the interferometry experiments. This possibility has been suggested in the phenomenological analysis of the asymmetric reaction $`Ar+Au`$ at 30 MeV/A .
These two unusual effects : the nonlinear velocity profile with $`\alpha _{surf}[1.5,2.0]`$, and an approximately power law fall-off of the particle density with $`\beta 3`$, characterize the transitional region in the central collisions of symmetric HI collisions in the Fermi energy domain.
For later times ($`t>200`$ fm/c), one sees the formation of an oblate configuration which is connected to the inversion of the velocity of the inner matter and to the accumulation of the density. In agreement with earlier observations , we see that the formed hot and nearly fused matter is not spherically symmetric. We would like to remark that this deformation is specific for energies around the Fermi energy and, moreover, is impact-parameter dependent. At nonzero impact parameters, the shape of the matter distribution becomes prolate due to the spectator matter keeping its initial direction of motion and also due to the angular momentum effects.
The evolution picture shown in Figs. 1 and 2 is essentially changed at higher bombarding energies. At $`E_{lab}/A=60`$ MeV (see Figs. 3 and 4), the time interval where the Coulomb force counter-balance the nuclear forces is becoming very short (e.g., see the plots for $`t=40`$ fm/c) and the system enters very fast in the phase of a smooth radial expansion. In this case, the velocity - radius scaling can be well approximated asymptotically by a single exponent $`\alpha _{bulk}=\alpha _{surf}1`$. The particle density is expanding almost uniformly and no characteristic power law dependence is seen in the surface region. The proton excess in the surface region is less pronounced than at $`E_{lab}/A=33`$ MeV/A (see Fig. 1) and $`Z/N1`$ at later times. One expects that with increasing bombarding energy in symmetric HI reactions the effective source parameters for protons and neutrons become close to each other.
### C The expansion velocity
Let us now discuss the dependence of the expansion velocity on the radius in more details. As was noted above, two different slopes can be distinguished at lower collision energies. We call the ’inner’ and the ’outer’ parts of the density profile the ’bulk’ and the ’surface’ regions, respectively. They are separated here at around $`R=10`$ fm. Therefore, we plot in Fig. 5 the time dependence of the exponent $`\alpha `$ in these two regions for different bombarding energies.
Let us start with the reaction at $`E_{lab}/A=33`$ MeV. The very first stage of the collision, where the nuclei are slightly overlapping, is characterized by similar values of the exponents $`\alpha `$ in the bulk and the surface regions. This feature continues until the surface particles begin to be evaporated. At around this time, the surface develops a much steeper velocity gradient with $`\alpha _{surf}`$ as large as $`2`$. At the same time, the bulk matter velocity profile is quite smooth and even $`\alpha _{bulk}`$ changes the sign . After this overshooting of the surface exponent during a time interval $`100200`$ fm/c, the bulk matter develops the radial expansion with the coefficient $`\alpha _{bulk}`$ which is approaching $`\alpha _{surf}`$. Note, that even at this stage the mass current is characterized by the nonlinear scaling. Hence, using the Hubble ansatz for the radial flow in the analysis of experimental data in the Fermi energy domain is not justified.
At lower energies ($`E_{lab}/A=15`$ MeV in Fig. 5), one observes that the difference between radial flow patterns in the surface and bulk regions becomes even stronger than at $`E_{lab}/A=33`$ MeV. Consequently, higher values of the exponent $`\alpha _{surf}`$ are reached asymptotically. Moreover, we see that a giant resonance with a period of $`T=60`$ fm/c, or an energy of $`2\pi /T=20.6`$ MeV, is excited. This can be considered as a complete fusion event.
If one proceeds to energies higher than the Fermi energy, one sees that the general trend of evolution is conserved but the deviation between surface and bulk matter expansion patterns decreases, i.e., the coefficients $`\alpha _{bulk}`$ and $`\alpha _{surf}`$ become close one to another. One can see also that the maximum value of the exponents $`\alpha `$ is reached faster than at lower energies, and the exponents $`\alpha `$ are becoming close to 1 at late times . The nuclear system at this high excitation energy is expanding continuously. Due to the higher initial velocity, there is practically no inversion of the velocity profile and no time-periodic structures like giant resonances are excited. The Hubblean expansion pattern is reached faster and no anomalous behavior ($`\alpha _{surf}>1`$) is observed for energies higher than $`90`$ MeV/A.
### D Dynamical trajectories
To elucidate the connection of the anomalous velocity profile with the multifragmentation, let us characterize an instantaneous state of the system in terms of the average nucleon density $`n`$ and the temperature $`T`$. In order to define a global time dependent temperature $`T(t)`$ we adopt the Fermi liquid relation :
$`E(t)={\displaystyle \frac{3}{5}}E_F(t)+E_{\mathrm{coll}}(t)+{\displaystyle \frac{\pi ^2}{4E_F(t)}}T(t)^2,`$ (12)
where the global kinetic energy $`E(t)`$, the collective energy $`E_{\mathrm{coll}}`$ and the Fermi energy $`E_F(n)=(3\pi ^2n/2)^{2/3}/2m`$ are given by a spatial integration of the local quantities (3) :
$`E(t)`$ $`=`$ $`{\displaystyle \frac{{\displaystyle 𝑑rE(r,t)}}{{\displaystyle 𝑑rn(r,t)}}}`$ (13)
$`E_F(t)`$ $`=`$ $`{\displaystyle \frac{{\displaystyle 𝑑rE_F\left(n(r,t)\right)n(r,t)}}{{\displaystyle 𝑑rn(r,t)}}}`$ (14)
$`E_{\mathrm{coll}}(t)`$ $`=`$ $`{\displaystyle \frac{{\displaystyle 𝑑r\frac{J(r,t)^2}{mn(r,t)}}}{{\displaystyle 𝑑rn(r,t)}}}.`$ (15)
It is more problematical to define a density. We will present here the two possibilities. The first one is to consider the density of matter inside the evolving mean square radius. The other possibility is to consider the density of matter contained in the static separation of bulk matter by the radius $`R`$ found earlier in the velocity profile.
The dynamical trajectories in the $`(Tn)`$ \- plane for different collision energies are shown in Fig. 6 for the bulk matter ($`R<10`$ fm) . Let us first look at the static density definition inside the bulk region. For $`E_{lab}/A=`$ 15 and 33 MeV, the system evolves inside the spinodal region and, therefore, is mechanically unstable. In the time interval between 150 fm/c and 200 fm/c, the system is in a configuration with $`nn_0/3`$ and $`T78`$ MeV, which are the typical values for the nuclear multifragmentation. At higher energies, the freeze-out density of an evolving system is shifted towards lower densities and finally it ends in a gaseous phase. If one chooses to look at the interior part of the system within the radial size of the mean squared radius (dotted lines in Fig. 6), the main difference is seen in the initial stage of the evolution where larger nucleon densities are reached but, at the same time, they are passed through much faster. The freeze-out configurations are practically the same as when the above static definition of bulk matter is considered.
### E The long tail of the density distribution
Let us now discuss the particle density profile in more details. As noted in Fig. 2, the decrease of the particle density in the surface region is algebraic (power law) rather than exponential. In Fig. 7, we plot the time dependence of the exponent $`\beta `$ which is extracted from the power law fit : $`n(r)r^\beta `$, for $`R>10`$ fm.
There are here two distinct behaviors. For $`E_{lab}/A=15`$ and 33 MeV, after an initial build-up of the surface region as characterized by the decrease of $`\beta `$ in time, the value of exponent $`\beta `$ stabilizes asymptotically. The limiting value of $`\beta `$, called $`\beta _{lim}`$, decreases with bombarding energy from $`\beta _{lim}3.5`$ at $`E_{lab}/A=15`$ MeV to $`\beta _{lim}3.1`$ at $`E_{lab}/A=33`$ MeV. Actually, the most interesting long-range region of the nucleon density ends at about $`E_{lab}/A50`$ MeV and the smallest value reached is about $`\beta _{lim}=2`$. In the energy interval $`15`$ MeV $`\text{ }<E_{lab}/A\text{ }<50`$ MeV, the anomalous power law tail of the nucleon density accompanies the anomalous profile of the expansion velocity in the surface region, and both effects are leading to the non-Gaussian shape of the emitting source and the anomalous short-range correlations.
For energies $`E_{lab}/A60`$ MeV in the Hubblean expansion regime, the exponent $`\beta `$ decreases monotonously in time and no asymptotically stable surface region with the power law dependence is seen. At later times , one sees however the appearance of a new expansion regime corresponding to the negative values of the exponent $`\beta `$. This indicates the formation of a shell-like structure in the system. The formation time of the shell-like structure decreases rapidly with increasing bombarding energy ($`t_{sh}280`$ fm/c at $`E_{lab}/A=60`$ MeV and $`t_{sh}170`$ fm/c at $`E_{lab}/A=90`$ MeV). The $`\beta `$ values found are not strongly negative, indicating that the expanding shell is a very diffused object due to the long mean-free path of nucleons in the kinetic approach. Similar unusual solution of spherically expanding scaling hydrodynamic has been used in the analysis of the Bose-Einstein correlations .
The specific shape fluctuations in the long-range region can be a source of the power-law Bose-Einstein correlations. This problem has been discussed by Białas in the context of strong interaction physics at relativistic energies. The long tail of the particle density for the systems produced in the symmetric heavy-ion collisions in the Fermi energy domain, may result in the unusual quantum statistical correlations. The $`N`$-particle interferometry reduced densities can be written as :
$`D_N(k_1,\mathrm{},k_N)=`$ (16)
$`=`$ $`{\displaystyle }dx_1\mathrm{}{\displaystyle }dx_N\left({\displaystyle \underset{per}{}}\mathrm{exp}[i(x_1p_{a_1}+\mathrm{}+x_Np_{a_N})]\right)\times `$ (17)
$`\times `$ $`n_N(x_1,\mathrm{},x_N)\left[N!\left({\displaystyle 𝑑xn_1(x)}\right)^N\right]^1,`$ (18)
where $`n_N`$ is the $`N`$-point density of emitting sources and the sum runs over all permutations of the indices $`a_i`$. The above formula supposes the incoherent emission from the source and neglects the final state interaction. In the case of uncorrelated emission in space-time :
$$n_N(x_1,\mathrm{},x_N)=n_1(x_1)\mathrm{}n_1(x_N),$$
the $`N`$-particle cumulant can be written using the Fourier transform of the source density $`n_1(x)`$ :
$`n_N(k_1,\mathrm{},k_N)={\displaystyle n_1(k_1k_{a_1})\mathrm{}n_1(k_Nk_{a_N})},`$ (19)
where the sum runs over all permutations of the indices $`a_i`$ with $`a_ii`$. If the source density has a power law tail :
$`n_1(x)x^\beta x^{\gamma D},`$ (20)
its Fourier transform shows also a power law in some range of small momenta :
$`n_1(k)|q|^\gamma .`$ (21)
Thus there is the relation between the power law tail in the source density distribution and the power law in the two-particle Bose-Einstein correlations which are given in terms of the Fourier transform of the source density :
$`C_2(q)|q|^{2\gamma }.`$ (22)
Similarly, the higher order cumulants $`C_i(q)`$ ($`i=3,\mathrm{}`$) are expected also to have a power law dependence on the rescaling of momenta with an index $`i\gamma `$.
A particularly interesting case of the Bose-Einstein correlations, increasing as the power of $`|q|`$, corresponds to $`\gamma >0`$ ($`\beta <D`$) in Eq. (20). In the case of central heavy-ion collisions, our analysis suggests that this effect could be seen in symmetric systems for energies : $`35`$ MeV $`\text{ }<E_{lab}/A\text{ }<50`$ MeV.
## III The summary and outlook
The dynamical behavior of heavy-ion reactions in the Fermi energy domain during the first $`200`$ fm/c is clearly associated with a nonlinear $`(\alpha >1)`$ radial velocity profile (2). The existence of such a non-Hubblean radial expansion at the freeze-out configuration was postulated by Le Févre et al. in the analysis of experimental kinetic energies of fragments in the $`Xe+Sn`$ reaction at $`E_{lab}/A=50`$ MeV in the framework of the statistical microcanonical model. The present studies using the nonlocal kinetic theory show that indeed such an unusual radial flow velocity profile is a plausible freeze-out configuration for symmetric heavy-ion collisions. One can roughly determine the instant of time when the compression turns into the expansion by looking at the crossing point of the bulk $`\alpha _{bulk}`$ and surface $`\alpha _{surf}`$ scaling exponents. While at the compression stage these two exponents are almost equal, the expansion stage exhibits larger exponent for the surface matter than for the bulk matter and $`\alpha _{surf}`$ takes values significantly larger than 1. This dynamical behavior of the surface matter disappears for energies significantly higher than the Fermi energy.
For central heavy-ion collisions in the Fermi energy domain, we see that the expansion stage is characterized by very small values of the exponent $`\alpha `$ in the bulk ($`\alpha _{bulk}[0,1]`$), indicating its slow evolution and possible mechanical instabilities. At these energies, the system spends a long time in the spinodal region, what may result in its multifragmentation decay. With increasing collision energy, the system passes quickly through this unstable region or, at even higher energies, its thermodynamic trajectories go above the spinodal region. In this case, the multifragmentation due to the spinodal instabilities is hardly possible . For details see
The nonlinear radial velocity profile in the surface region is accompanied by the long tail of the particle density. Both these effects may have important consequences on the quantum statistical correlations and their evolution with the bombarding energy. The effect of the radial expansion on the two-particle correlations have been studied assuming the linear scaling solution ($`\alpha =1`$) of the scaling hydrodynamics . It was found that the expansion makes the effective radius of the two-particle correlation functions smaller than the geometrical size of the source. We expect this effect to be even stronger in the presence of the nonlinear radial expansion flow, leading to even stronger discrepancy between the effective radius and the geometrical radius. Moreover, the commonly used Gaussian approximation for the source shape is certainly hazardous in this non-Hubblean expansion regime. The existence of the power law tail in the particle density can in turn lead to the power law two-particle Bose-Einstein correlations at small relative momenta of the particles in the range of collisions energies $`35`$ MeV $`\text{ }<E_{lab}/A\text{ }<50`$ MeV. It should be stressed that the ratio $`Z/N`$ in the surface region is strongly different from 1, in particular at early collision times and at low bombarding energies ($`E_{lab}/A\text{ }<60`$ MeV). This specific effect in the Fermi energy domain, which leads to an effective increase of the proton rms radius and to the change in the temperature gradients in proton/neutron source, may have measurable consequences in the quantum statistical correlations for protons and neutrons. This effect has been studied by Helgesson et al. assuming a linear scaling solution ($`\alpha =1`$) of hydrodynamics for asymmetric HI reactions at $`30`$ MeV/A. Our results suggest that the simultaneous description of $`n`$ and $`p`$ spectra, as well as $`nn`$ and $`pp`$ correlation functions may require different source parametrizations for neutrons and protons, though the detailed dynamics for HI reactions at $`E_{lab}/A30`$ MeV is very different from those assumed by Helgesson et al. . It is interesting to notice , that the nonlocal kinetic theory predicts the appearance of the solution somewhat similar to the linear scaling solution ($`\alpha =1`$) of hydrodynamics with $`\beta <0`$ for higher energies bombarding energies ($`E_{lab}/A\text{ }>60`$ MeV.
The different behavior of the surface and bulk matter before equilibration can also be of importance for the description of super nova where a surface - like ring of matter (crust) is expanding with enormous velocities and is clearly separated from the remaining bulk matter collapsing back into neutron stars. The clearest experimental observation of the expanding shell-like structure noted above comes from stellar astronomy. The envelope material ejected by the stars forms an expanding shell of gas that is known as a planetary nebula. The space-time evolution of these objects is in many aspects similar to the considered evolutions : $`\alpha \text{ }<2`$, $`\beta >0`$ and $`\alpha 1`$, $`\beta <0`$. The latter solution can be successfully simulated by a scaling solution of the non-relativistic hydrodynamics .
We suggest that the anomalies found in the kinetic expansion reflect the nature of effective interactions among the elementary constituents of the system. Their manifestation is twofold. Firstly, the interplay between the repulsive Coulomb interaction and the attractive mean field results in the formation of a rather sharp surface of the system. Secondly, the evolution of the bulk matter is not a simple uniform expansion of a homogeneous ideal fluid. In the unstable spinodal region, the local interaction of quasiparticles in the bulk phase affects noticeably the subsequent evolution of the system. In this respect, one should stress again an attractive possibility of implementing both the nonlinear velocity profile and the algebraic long-range density tail into the analysis of the Bose-Einstein correlations. These effects are not only important for the nuclear multifragmentation process but also for the hadron interferometry at ultrarelativistic collisions where the deconfinement phase transition can have a strong influence on the expansion stage. This is of a particular interest for the statistical mixed phase equation of state which gives the crossover type of the deconfinement phase transition and allows for a small admixture of unbound quarks at the freeze-out point .
An interesting example which illustrates an important difference between the linear and quadratic scalings of the expansion velocity with the radius can be found in cosmology . If there is an attractive force decreasing as an inverse power of the radius, the equation of motion for a radially symmetric matter follows from the total energy $`h`$ which reads :
$`{\displaystyle \frac{m}{2}}\dot{R}(t)^2{\displaystyle \frac{G}{R(t)^\delta }}=h.`$ (23)
Assuming a homogeneous matter density $`n`$, the mass is : $`m=4\pi nR^3/3`$. In the case of the escaping matter $`(h=0)`$, one gets :
$`\dot{R}(t)=\sqrt{{\displaystyle \frac{6G}{4\pi n}}}R(t)^{{\displaystyle \frac{3\delta }{2}}},`$ (24)
what corresponds to taking :
$`\alpha ={\displaystyle \frac{3\delta }{2}}`$ (25)
in Eq. (2). We see that $`\delta =1`$ for the Coulomb or gravitational forces and, hence, the Hubble expansion ($`\alpha =1`$) follows in these cases. However if $`\alpha =2`$, as found in the surface region of nuclei formed in central HI collisions at around the Fermi energy (see Fig. 5), the above relations lead to a string-like force with $`\delta =1`$. Since this force is used to describe the confinement in the effective theories motivated by the Quantum Chromodynamics, its occurrence as a consequence of the dynamical behavior is worth of attention. As was shown above in the solution of nonlocal kinetic equations, the interplay between Coulomb and mean field leads to such a string-like behavior for surface particles. This is here a clear non-equilibrium effect. The above discussed oscillation in the time-dependence of the $`\alpha `$ exponent may be considered as a possible manifestation of this effective dynamical string-like force in the surface region.
Perhaps the best way to demonstrate the existence of both the non-Hubblean radial expansion and the algebraic long-range tail of the particle density, would correspond to finding a non-Gaussian deformation of the source in the Bose-Einstein interference experiments for central, symmetric heavy-ion reactions in the narrow range of collisions energies ($`35`$ MeV $`\text{ }<E_{lab}/A\text{ }<50`$ MeV). In the same narrow range unusual scale - dependence of the many-particle correlations at small momenta should be induced via the same Bose-Einstein quantum interference effect if the source has an algebraic (power law) density tail (Sect. II.E). On the other hand, as discussed in , the kinematical observables related mainly to the intermediate mass fragments, provide an independent and sensitive measure of the velocity profile in the deformed, expanding source. To put together all these different pieces of evidence into the circumstantial proof for the non-Hubblean collective expansion of nuclear aggregates, remains a difficult and exciting challenge for the future experimental and theoretical studies of heavy-ion collision in the Fermi energy domain.
## IV Acknowledgements
This work was supported by the IN2P3-JINR agreement No 0049. K. M. likes to thank the LPC for hospitality and friendly atmosphere.
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# Dynamical masses, time-scales, and evolution of star clusters
## 1. Introduction
The Milky Way and probably all large galaxies contain old globular cluster populations (see the review by Harris 1991). These old star clusters have an approximately log-normal luminosity function, and the mean cluster luminosity is somewhat brighter than M$`{}_{V}{}^{}=7`$ with little dependence on the host galaxy luminosity. In the Milky Way their typical mass-to-light ratios are $`M/L_V2`$, and typical total masses are $`2\times 10^5\mathrm{M}_{}`$ (Pryor & Meylan 1993). It is widely assumed that the globular clusters we see today must be the part of an initially larger population that survived the internal and external dynamical processes leading to cluster destruction (e.g., Ostriker 1988).
One of the exciting results from HST has been the discovery of young star clusters in starburst and interacting galaxies. Whitmore & Schweizer (1995) found many hundreds of young clusters in the Antennae galaxies. Young cluster systems have now been discovered in other interacting and merging galaxies, in barred and starburst galaxies, and even dwarf starburst galaxies (e.g., ESO 338-IG04, Oestlin, Bergvall & Roennback 1998). The luminosity functions of the young clusters are not log-normal, but seem to be better described by power-laws, about $`L^2`$. Carlson et al. (1999) use population synthesis models to determine the ages of the blue clusters in the young merger remnant NGC 3597. Based on these models they argue that the difference in the observed luminosity function when compared to the Galactic globular clusters cannot simply be an age effect, even if the young clusters formed with an intrinsic age spread. Are these young cluster systems then a good model for what the Milky Way’s globular cluster population could have looked like at birth?
This review gives a brief discussion of dynamical methods to determine masses of distant and nearby star clusters (Section 2). It then goes on to describe a number of dynamical processes and their time-scales which will lead to evolution and potentially destruction of star clusters over long time-scales. Finally, the results of some evolution calculations for globular cluster systems are briefly summarized (Section 3).
## 2. Dynamical Mass Determination for Star Clusters
In this Section we discuss methods for estimating star cluster masses from structural and kinematic measurements. Mass estimates based on stellar population properties are discussed in U. Fritze von Alvensleben’s article in these proceedings.
### 2.1. Virial Masses
A simple global mass estimate for a star cluster can be obtained from the virial theorem. This says that, in equilibrium, the radius of a stellar system is proportional to $`GM/V^2`$, where $`M`$ is the total mass and $`V`$ the rms three-dimensional velocity of the stars. The constant of proportionality generally depends on the stellar density profile, but Spitzer (1969) showed that if the relation is expressed in terms of the half–mass radius $`r_h`$, this dependence is weak and the constant is approximately 0.4 for realistic cluster profiles. If we furthermore assume that the cluster is spherical, $`V^2=3\sigma _{}^2`$, where $`\sigma _{}`$ is the one-dimensional rms velocity dispersion along the line-of-sight, and write $`\sigma _{10}=\sigma _{}/10\mathrm{km}\mathrm{s}^1`$ and $`r_5=r_h/5\mathrm{pc}`$, then
$$M_V=7.5\sigma _{}^2r_h/G=8.7\times 10^5\sigma _{10}^2r_5\mathrm{M}_{}.$$
(1)
When using this formula to estimate star cluster masses from observed velocity dispersions and radii, a few points should be noted:
(i) Because the virial mass (1) is a global estimate, it is independent of velocity anisotropy. For example, shifting some stars to radial orbits while keeping the (spherical) potential fixed, will result in a larger central velocity dispersion but also lead to reduced velocities in the cluster halo. To maintain virial equilibrium these changes must add in just such a way that the global $`\sigma _{}`$ remains the same.
(ii) The dynamical evolution of star clusters leads to mass segregation and the formation of a halo of low-mass stars on preferentially radial orbits (see §3.2. below). For evolved clusters the measured half-light radius will therefore in general underestimate the half-mass radius, the observed velocity dispersion will underestimate the rms velocity dispersion, and eq. (1) will underestimate the mass.
(iii) Sometimes only the velocity dispersion for stars in the core is known, or the velocity dispersion from integrated light within some aperture. In these cases, a dynamical model is needed to convert this to the rms $`\sigma _{}`$. This introduces some uncertainty in the mass estimate because the derived $`\sigma _{}`$ depends on anisotropy.
With high-resolution spectra and HST photometry virial masses can be determined for some young ‘superclusters’ seen in starburst galaxies. Masses for three clusters in M82 are compared by Smith & Gallagher in these proceedings, spanning a range from $`3\times 10^5\mathrm{M}_{}2\times 10^6\mathrm{M}_{}`$.
### 2.2. Core Masses
Rood et al. (1972) gave a formula that is often used to determine core masses. This is based on the dynamics of King models (King 1966; Binney & Tremaine 1987) and assumes that the velocity distribution in the core is isotropic:
$$\left(\frac{M}{L}\right)=\frac{9\sigma _0^2}{2\pi GI_0r_c}.$$
(2)
Here $`\sigma _0`$ is the central velocity dispersion, $`I_0`$ the central surface brightness and $`r_c`$ the core radius. The product of the two last quantities is rather insensitive to errors caused by seeing. Eq. (2) is very accurate as long as the assumption of isotropy is met (Richstone & Tremaine 1986); but it can overestimate the mass by a factor $`2`$ if the system is actually radially anisotropic (Merritt 1988).
Core mass-to-light ratios can only be determined for well-resolved Galactic star clusters for which the core parameters $`r_c`$, $`I_0`$ and $`\sigma _0`$ can be estimated. Even with the resolving power of HST the cores of distant young clusters cannot be resolved.
### 2.3. Masses from Model Fitting
An alternative method of estimating cluster masses is fitting the photometric and kinematic data with dynamical models. A simple such scheme was used by Djorgovski et al. (1997) in their study of the M31 globular clusters mass-to-light ratios. They used structural and photometric parameters for these clusters obtained with HST and kinematic measurements in a rectangular aperture obtained with Keck and HIRES. They then estimated an aperture correction from King models to transform their measurements to central velocity dispersions, and used a formula analogous to eq. (1) to estimate masses with the constant again determined from models.
For some Galactic globular clusters large velocity samples are available and in such cases much more detailed model fitting is possible (Pryor et al. 1989). The masses of the Galactic globular clusters referred to in §1. (Pryor & Meylan 1993) have been determined by these techniques. A recent such study is Côté et al. (1995) who investigated the dynamics of the globular cluster NGC 3201, using a CCD surface brightness profile and a sample of 857 measured stellar radial velocities to trace the velocity dispersion profile to large radii.
In such work the data are fitted by single- or multi-mass King-Michie models. In the multi-mass models, a power-law mass function for the cluster stars is typically assumed, and for each mass bin $`m_i`$, a distribution function of King-Michie type (Michie 1963, Binney & Tremaine 1987) is used:
$$f_i(E,J)e^{\beta J^2}\left(e^{A_iE}1\right).$$
(3)
Here $`E`$, $`J`$ are the specific energy and angular momentum of a star in the (spherical) star cluster potential, and $`\beta `$ can be thought of as specifying an anisotropy radius. In the core of the cluster, stars of different masses are assumed to be in equipartition (§3.3.), so that $`A_im_i`$. In the fitting procedure the free parameters are the radius, velocity and luminosity scale, the cluster’s concentration parameter, the anisotropy radius, and the index of the mass function. These parameters are determined from fitting to the measured surface brightness and velocity dispersion profile. This leads to a determination of the mass-to-light ratio profile and anisotropy profile, rather than just a single $`M/L`$ as for the previously described techniques. However, the fit is non-unique in the sense that adding even fairly large numbers of faint low-mass stars in the halo (expected there from mass-segregation and evaporation, see §3.2. below) have little effect on the observed profiles. In their study, Côté et al. (1995) find a steady rise in M/L with distance from the cluster center, as expected from dynamical evolution theory, and a global $`M/L_BM/L_V2.0\pm 0.2`$.
### 2.4. Masses from Proper Motions
For nearby Galactic globular clusters, it is possible to measure stellar proper motions in addition to radial velocities. Proper motion measurements give information about the velocity dispersions in two directions on the sky (radial along projected $`R`$, and tangential), and for a spherically symmetric cluster they are therefore in principle sufficient to determine the velocity ellipsoid as a function of radius, and thus the mass profile free of assumptions about anisotropy. The projected proper motion dispersions $`\sigma _R(R)`$ and $`\sigma _T(R)`$ are related to the intrinsic velocity dispersions $`\sigma _r(r)`$ and $`\sigma _t(r)`$ by Abel integral equations and can thus be inverted (Leonard & Merritt 1989). Moreover, from the inferred $`\sigma _r(r)`$ and $`\sigma _t(r)`$ one can predict the line-of-sight velocity dispersions $`\sigma _{}(R)`$ and compare with independent radial velocity data. This provides a check on the modelling and also can be used to determine the cluster distance. In terms of global velocity dispersions, $`\sigma _{}^2=(\sigma _R^2+\sigma _T^2)/2`$ for a spherical cluster and the correct distance.
In an early study along these lines Leonard et al. (1992) investigated radial velocity and proper motion data for the globular cluster M13. They concluded that the mean anisotropy of this cluster $`\beta =3(\sigma _R^2\sigma _T^2)/(3\sigma _R^2\sigma _T^2)0.3`$ and that the effect of the anisotropy on the mass determination is $`20\%`$. Much more detailed modelling will be possible with the large proper motion surveys currently in progress.
### 2.5. Non-Parametric Cluster Mass Distributions
With large samples of stellar velocities at hand, radial velocities or proper motions, it is possible to infer the mass distribution of the cluster without making specific assumptions like King-Michie stellar distribution functions. This requires solving the Jeans and projection equations for the intrinsic density and velocity dispersions under some smoothness constraint, given the data. I do not give the equations here, but refer to the papers mentioned below.
When the data consist of several hundred radial velocities, some assumption about the anisotropy is still needed. Gebhardt & Fisher (1995) describe such a non-parametric analysis of radial velocity data for four Galactic globular clusters, assuming isotropy of the stellar orbits. With a few hundred stellar velocities in each case the results are still noisy, but indicate radially increasing $`M/L`$-profiles as expected. Merritt, Meylan & Mayor (1997) describe a similar analysis of the cluster $`\omega `$ Centauri, assuming that it is oblate and seen edge-on, and that it is described by a meridionally isotropic two-integral model. They find that the mass distribution cannot be strongly constrained by their data, but appears to be slightly more extended than the luminosity distribution.
As discussed above, proper motion data result in two independent velocity dispersions in the plane of the sky, and thus, within a spherical model, they contain sufficient information to determine the anisotropy of the stellar orbits. With sufficiently large data sets it will therefore be possible to model the anisotropy profile and mass distribution of a spherical cluster non-parametrically.
## 3. Dynamical Evolution Processes and their Time-Scales
### 3.1. Relaxation
On dynamical time-scales large star clusters ($`N>>100`$) evolve collisionlessly. That is, for some time after birth they are described by a quasi-equilibrium phase-space distribution function (df) which is a function of the integrals of motion (or of the stellar orbits) in the mean field potential. On longer times-scales, however, the graininess of the distribution of stars becomes important, and the dynamical evolution is no longer collisionless. Over a relaxation time two-particle interactions then deflect the cluster stars from the orbits they would otherwise have followed in the mean gravitational potential.
In the approximation of a homogeneous distribution of equal-mass stars, with density $`\rho _0`$ and isotropic Maxwellian velocity distribution with dispersion $`\sigma _0`$, the two-body relaxation time is (Spitzer & Hart 1971)
$$t_{r0}=0.34\frac{\sigma _0^3}{G^2m_{}\rho _0\mathrm{ln}\mathrm{\Lambda }}=1.8\times 10^8\sigma _{10}^3n_4^1m_,^2(\mathrm{ln}\mathrm{\Lambda })_{10}^1\mathrm{yr}.$$
(4)
Here $`n=10^4n_4\mathrm{pc}^3`$ is the number density of stars, $`\sigma _0=10\sigma _{10}\mathrm{km}\mathrm{s}^1`$ the velocity dispersion, $`m_{}=m_,\mathrm{M}_{}`$ is the mean mass per star, $`\mathrm{\Lambda }`$ is of order the number of stars, and $`(\mathrm{ln}\mathrm{\Lambda })_{10}=(\mathrm{ln}\mathrm{\Lambda })/10`$. Note that $`t_{r,0}`$ is inversely proportional to the stellar phase space density. Central relaxation times in globular cluster cores evaluated with eq. (4) are $`10^710^9\mathrm{yr}`$.
The relaxation time often varies by large factors between the central and outer parts of a stellar system. It is then useful to define a half-mass relaxation time. For a virialized star cluster this is obtained from eq. (4) by replacing $`\rho _0`$ with the mean density inside the system’s half-mass radius $`r_h=5r_5\mathrm{pc}`$ and $`\sigma _0^2`$ by one third of the rms $`V^2`$, and then using the virial theorem to express $`V^2`$ through $`r_h`$ and the total cluster mass $`M=10^5M_5\mathrm{M}_{}`$. The result is (Spitzer & Hart 1971)
$$t_{rh}=\frac{0.14N}{\mathrm{ln}0.4N}\left(\frac{r_h^3}{GM}\right)^{\frac{1}{2}}=\frac{N}{26\mathrm{log}0.4N}t_d=\mathrm{\hspace{0.17em}7.2}\times 10^8M_5^{1/2}r_5^{3/2}m_,^1(\mathrm{ln}\mathrm{\Lambda })_{10}^1\mathrm{yr},$$
(5)
where $`N`$ is the total number of stars in the system and
$$t_dr_h/V=1.58(r_h^3/GM)^{1/2}=8.3\times 10^5r_5^{3/2}M_5^{1/2}\mathrm{yr}$$
(6)
is the dynamical time. For comparison with the local formula eq. (4), the fiducial values used in eq. (5) correspond to a one-dimensional virial velocity dispersion $`\sigma 3.4\mathrm{km}\mathrm{s}^1`$ and a mean density $`n96\mathrm{pc}^3`$.
On short time-scales cluster evolution is still collisionless, so long as $`t_dt_{rh}`$ (requiring $`N100`$); for example, during the violent relaxation at formation. The resulting quasi-equilibrium df subsequently evolves slowly in response to collisions, which will tend to drive the system towards an isothermal energy distribution. One aspect of such slow evolution would be a decrease of ellipticity with dynamical age (Fall & Frenk 1985). This could be the reason for the significantly rounder globular clusters in M31 and the Milky Way as compared with the LMC and SMC clusters (Han & Ryden 1994).
### 3.2. Evaporation and Core Collapse
Collisions between single stars modify the stellar df in two ways. The rarer process is ejection, in which a single close encounter leads one of the stars to acquire a velocity greater than the local escape velocity $`v_e`$ and to escape from the cluster. The time-scale for this is $`t_{ej}N/(dN/dt)1.1\times 10^3\mathrm{ln}0.4Nt_{rh}`$ $`10^4t_{rh}`$ (Hénon 1969). The more important process of evaporation is caused by the cumulative effect of many weak encounters, which gradually increase a star’s energy until $`vv_e`$. It is easy to show that the rms escape velocity of the cluster is just twice its rms virial velocity. Thus, on average, a particle with $`v2V=\sqrt{12}\sigma _{}`$ will escape. For a Maxwellian velocity distribution, a fraction $`ϵ0.74\%`$ of stars have $`v2V`$; these stars will escape in one dynamical time, after which the high-velocity tail is repopulated only in $`t_{rh}`$. Thus one expects the evaporation time scale of the cluster to be $`(ϵ/t_{rh})^1`$; detailed calculations (Spitzer & Thuan 1972) show that
$$t_{ev}N(dN/dt)^1\mathrm{\hspace{0.33em}300}t_{rh}.$$
(7)
Because the evaporation is dominated by weak encounters, escaping stars leave the cluster with only very small positive energy; thus the total energy of the remaining cluster is nearly constant, but must be shared among a shrinking number of stars. In virial equilibrium $`N^2/r_h`$ const. and thus $`\rho N/r_h^31/N^5`$ and $`t_{rh}Nr_h^{3/2}/M^{1/2}N^{7/2}`$. So as the cluster becomes denser, evaporation accelerates and the system contracts to negligible mass and radius in finite time.
This evaporation model, however, neglects the fact that the evolution of the stellar cluster is not homologous and that the rate of evolution is much faster in the dense core than in the system’s outer parts. Stars gaining energy towards evaporation build up an extended halo where the time scale for further energy gain increases strongly, so that these stars may not in fact escape during the age of the cluster. On the other hand, the dense core loses stars to the halo on the much faster central relaxation time, and may collapse to very high densities before $`M_{\mathrm{tot}}`$ and $`r_h`$ can change much.
This phenomenon of core collapse may be understood as a consequence of the fact that self-gravitating star clusters have negative specific heat (Lynden-Bell & Wood 1968): In virial equilibrium the total energy $`E=T`$, where $`T`$ is the total kinetic energy, which is proportional to the virial temperature $`T/M=V^2`$. As energy is withdrawn from the cluster, its kinetic energy increases and so does the virial temperature. Since $`r_hGM^2/(2|E|)`$, the cluster thereby contracts. Vice-versa, an energy production mechanism (e.g., from binary stars) causes the cluster to cool and expand. Now the dense core of the cluster may be approximately regarded as a virialized system in thermal contact with the rest of the cluster. It is normally hotter than its surroundings and therefore loses energy to them through stellar encounters. As a result it shrinks and becomes yet hotter, loses still more energy to the surrounding stars, and contracts to formally zero radius in finite time.
For a single-mass star cluster the late stages of core collapse are self-similar (Lynden-Bell & Eggleton 1980, Cohn 1980). As the core radius
$$r_c3\sigma _c/\sqrt{4\pi G\rho _c}$$
(8)
shrinks, the central density $`\rho _c`$ and velocity dispersion $`\sigma _c`$ increase and the core mass $`M_c`$ decreases according to
$$\rho _cr_c^{2.23},\sigma _c(\rho _cr_c^2)^{1/2}r_c^{0.11},M_c\rho _cr_c^3,r_c^{0.77}$$
(9)
until the core radius and mass formally shrink to zero at time $`t_{cc}`$. Moreover, the density profile of the cluster outside the collapsing core has the same exponent: $`\rho r^{2.23}`$ for $`r_c(t)rr_c(0)`$.
The time-scale for core collapse is proportional to the central relaxation time; for a single mass cluster it is $`t_{cc}330t_{rc}`$ once the collapse is in the self-similar phase (Cohn 1980, Heggie & Stevenson 1988). The total time until core collapse in Cohn’s (1980) model is $`16`$ half-mass relaxation times or $`60`$ initial $`t_{rc}`$. As Goodman (1993) has emphasized, the former number depends on the mass distribution of the cluster, and $`t_{cc}/t_{rh}`$ will be less than 16 for clusters more centrally concentrated than Plummer models. By noting that $`r_c^1dr_c/dtt_{rc}^1\rho _c/\sigma _c^3r_c^{1.89}`$, one can solve for the asymptotic time-dependence of the collapse:
$$r_c(t_{cc}t)^{0.53},\rho _c(t_{cc}t)^{1.18},\sigma _c(t_{cc}t)^{0.06},M_c(t_{cc}t)^{0.41}.$$
(10)
In summary, the collapse of a single mass cluster occurs in two stages (Cohn 1980). The longer part of the evolution is an evaporative phase, during which stellar collisions simultaneously populate a halo and make the core shrink and become denser. Only towards the end does the evolution accelerate and enter the gravothermal instability phase of self-similar collapse.
### 3.3. Equipartition, Mass Segregation, and Multi-Mass Core Collapse
When the cluster contains different stellar masses, energy can flow not only from the core to the halo, but also between stars of different masses. Stellar collisions drive the system towards equipartition of energy $`m_iv_i^2`$ = const. As eq. (4) shows, relaxation proceeds faster for more massive stars. The equipartition time-scale measures the rate at which a group of heavy stars with masses $`m_2`$ loses energy to lighter stars of mass $`m_1`$ (Spitzer 1969):
$$t_{eq}=\frac{(v_1^2+v_2^2)^{3/2}}{8(6\pi )^{1/2}G^2m_1m_2n_1\mathrm{ln}\mathrm{\Lambda }}=1.2\frac{m_1}{m_2}t_{r0}(m_1)$$
(11)
where we have assumed equal temperatures $`v_1^2=v_2^2`$ and used eq. (4). Initially, $`v_i^2`$ is independent of stellar mass; thus the massive stars lose kinetic energy and sink to the center, while lighter stars gain kinetic energy in collisions and move outwards, a process called mass segregation. Moreover, eq. (11) shows that mass segregation of the massive stars occurs before relaxation of the cluster as a whole becomes significant.
However, equipartition may never be reached. A simple case considered by Spitzer (1969) is one with two mass groups such that the heavy stars are much more massive than the light stars, $`m_2m_1`$, but the total mass in the cluster core is dominated by the light stars: $`M_2\rho _1r_{c1}^3`$. In this case equipartition causes the formation of a small subsystem of heavy stars ($`M_2`$, $`m_2`$) in the core of the distribution of lighter stars. Applying the virial theorem to the subsystem of heavy stars gives
$$v_2^2\frac{0.4GM_2}{r_{h2}}+\frac{4\pi G\rho _{c1}}{3}k^2r_{h2}^2,$$
(12)
where the first term describes the self-energy of the subsystem $`M_2`$ and the second term its interaction with the system of light stars ($`k`$ is a constant of order unity). Spitzer (1969) noticed that the right-hand-side of eq. (12) has a minimum when regarded as a function of $`r_{h2}`$. An equilibrium can therefore exist only if $`v_2^2`$ is greater than this minimum, that is, assuming equipartition, if
$$\frac{M_2}{\rho _1r_{c1}^3}\mathrm{\hspace{0.33em}4.0}k^1\left(\frac{m_1}{m_2}\right)^{3/2}.$$
(13)
In other words, if its mass is too large, the subsystem of heavy stars remains a dynamically independent stellar system with mean square velocity greater than the equipartition value. It continues to lose energy to the lighter stars, becoming denser and hotter, and evolving away from equipartition all the time (mass stratification instability). Fokker-Planck calculations (Inagaki & Wiyanto 1984, Cohn 1985) show that in the end the subsystem of heavy stars core collapses independently from the cluster of light particles, just like a single mass system.
The evolution to core collapse with a spectrum of stellar masses has been considered by Inagaki & Saslaw (1985) and Chernoff & Weinberg (1990). The detailed evolution occurs in several phases: First, collisions trying to establish equipartition of energy lead to mass segregation and the formation of a heavy mass core. Then this core undergoes the gravothermal instability, i.e., contracts while remaining hotter than the mean temperature of the system and conducting energy outwards. This collapse accelerates towards core collapse, and finally goes over into a single-component collapse which reaches formally infinite central density. The time scale for this multi-mass core collapse evolution is faster than that for core collapse in any single component cluster, typically a factor of a few faster than for a cluster composed of the heaviest mass alone.
Deep in collapse, the density slopes of all mass groups $`m_k`$ are characterized by separate power laws in the region where the heaviest component dominates the potential, such that approximately (Cohn 1985, Chernoff & Weinberg 1990)
$$d\mathrm{ln}\rho _k/d\mathrm{ln}r1.89(m_k/m_u)0.35,$$
(14)
where $`m_u`$ is the mass of the heaviest species in the cluster. The overall mass profile is then not self-similar.
A multi-mass core collapse, however, may be strongly influenced by the stellar evolution of the more massive stars. This has two main effects: First, the mass loss from massive stars through winds may lead to an overall mass loss from the cluster, and thus cause an adiabatic expansion. Secondly, the finite stellar life-time $`t_{MS}`$ limits the time $`t_{cc}`$ during which they can core collapse, such that $`t_{cc}(m_{})\mathrm{}<t_{MS}(m_{})`$. Both effects greatly increase the overall core collapse times; compared to a system of point masses within the range $`(0.415)m_{}`$, Weinberg & Chernoff (1989) find an increase by about a factor of $`3060`$ in their globular cluster models, including the expansion effects.
A reasonable approximation for the stellar lifetime of all but the most massive stars is $`t_{MS}910^9(m_{}/m_{})^{2.6}\mathrm{yr}`$. Thus if the most massive stars leave black hole remnants of $`3\mathrm{M}_{}`$, these together with tight binaries will dominate the evolution after $`510^8\mathrm{yr}`$, while if the most massive remnants are $`1.4\mathrm{M}_{}`$ neutron stars, they and the binaries will dominate after $`410^9\mathrm{yr}`$. The latter time scale approaches the time expected for core collapse in typical Milky Way globulars.
### 3.4. Reversing core collapse
A number of energy source mechanisms can stop core collapse (e.g., Goodman 1993): (i) Processes that generate kinetic energy in the core directly, such as binary stars transferring energy to the field stars in collisions. (ii) Mass loss processes that heat the core indirectly above its virial temperature, including: normal stellar evolution, accelerated stellar evolution by the formation of massive stars in mergers, and ejection of stars through binaries. In all cases the net result is adiabatic expansion and cooling of the core.
Only hard binaries contribute to field star heating. Binaries are hard if their binding energy $`E_b=Gm_1m_2/a`$ (with $`a`$ their semi-major axis) exceeds the mean kinetic energy, $`E_b>3m_{}\sigma ^2`$; those with $`E_b<3m_{}\sigma ^2`$ are called soft binaries. Heggie’s law (Heggie 1975, Hut 1983) states that, on average, hard binaries get harder by collisions with field stars, and soft binaries get softer. Essentially, the orbital velocity of a hard binary is on average greater than the velocity of an incoming field star, and the tendency towards energy equipartition therefore results in a net transfer of energy to the field star. The opposite is true for soft binaries, which gain net energy and eventually dissolve. The binary behaves like a mini-system with negative specific heat: as energy is withdrawn from it, the orbit shrinks, the orbital velocity increases, and the binary hardens. When the binary becomes sufficiently hard, the typical recoil from a collision with a single star becomes large, and the binary will eventually be kicked out of the cluster. Just like in the Sun, the binaries providing the nuclear energy source will eventually be ’burned’.
Binaries can be formed by a close gravitational interaction of three stars (‘three-body binaries’), by dissipational tidal capture, or at the time of star formation (‘primordial binaries’). To be effective in reversing core collapse, binaries must have orbital semi-major axes
$$a<Gm_{}/3\sigma ^2=\mathrm{\hspace{0.17em}3}\sigma _{10}^2m_,\mathrm{AU}.$$
(15)
The formation of a hard three-body binary requires a close encounter between two stars ($`\delta vv`$) with a third star in the immediate vicinity, such that one of the three stars acquires additional energy, leaving the other two as a bound pair. Thus the time-scale is (Goodman 1984, Binney & Tremaine 1987)
$$t_3(np^2v)^1(np^3)^1\frac{\sigma ^9}{n^2(Gm_{})^5}N_c^2\mathrm{ln}N_ct_{r0},$$
(16)
where $`pGm_{}/v^2`$, $`v=O(\sigma )`$ because low relative velocities dominate, and we have used eq. (4) to express Goodman’s result in terms of the central relaxation time and the total number of stars in the core, $`N_c`$. This implies that about $`1/N_c\mathrm{ln}N_c`$ three-body binaries form per central relaxation time. In other words, three-body binaries become important if the final core collapse is driven by fewer than 100 of the largest mass stars.
### 3.5. Tidal field and tidal shocks
A steady tidal field lowers the energy threshold beyond which stars are no longer bound to the cluster. It thus increases the mass loss rates from evaporation, both because the fraction of stars in the velocity distribution that escape in a dynamical time increases, and because the decreasing number of stars in the cluster leads to shorter relaxation times. The Quintuplet and Arches young clusters (Figer et al. 1999) in the inner Galactic bulge are two clusters for which these tidal effects are very important (Kim et al. 1999).
In reality, the tidal field is not stationary in the frame of the cluster stars. This complicates the escape process, but more importantly it leads to a new dynamical process in cluster evolution, referred to as gravitational shocking (Ostriker, Spitzer & Chevalier 1972). The tidal field acting on the cluster may suddenly increase in strength when the cluster passes through the disk of its host galaxy, or when it comes close to the high-density inner bulge near peri-galacticon of its orbit. In both cases, the perburbations to the stellar orbits caused by the tidal shock lead to an effective energy input in the cluster which makes the cluster less bound and accelerates mass loss from internal processes.
A detailed recent discussion of this process is given by Kundić & Ostriker (1995) and Gnedin, Lee & Ostriker (1999). For stars in the outer parts of the cluster, the tidal perturbation can be approximated as impulsive because of the short time-scale of passage through the galactic disk. In the cluster’s central parts, on the other hand, the stellar orbital time-scales are short and adiabatic invariance reduces the effects of the perturbation. Traditionally these effects of the tidal shock were described by a first-order term $`(\mathrm{\Delta }E)_{ts}`$, which denotes the net energy gain averaged over stellar orbits at a given position in the cluster. Kundić & Ostriker (1995) noticed that the second-order term $`(\mathrm{\Delta }E)_{ts}^2`$ is typically even more important and competes with two-body relaxation near the half-mass radius in driving evolution of the cluster’s internal structure. This can speed up core collapse by a factor of three (Gnedin, Lee & Ostriker 1999). Cluster destruction is accelerated; recent modelling of the evolution of the Milky Way’s globular cluster system shows that the typical time to destruction becomes comparable to the typical age of the Galactic globulars (Gnedin & Ostriker 1997).
### 3.6. Dynamical friction and merging
As already noted by Tremaine, Ostriker & Spitzer (1975), massive star clusters experience dynamical friction against field stars as they move along their orbits through the host galaxy. Because of the frictional drag the cluster loses orbital energy and spirals into the galaxy center, where the tidal field becomes ever stronger and will eventually dissolve the cluster.
The time-scale for dynamical friction for a cluster on a circular orbit at initial radius $`r_i=2r_{i,2}\mathrm{kpc}`$ in a singular isothermal sphere with circular velocity $`v_c=250v_{250}\mathrm{km}\mathrm{s}^1`$ is (Chandrasekhar 1943, Binney & Tremaine 1987)
$$t_{\mathrm{df}}=\frac{1.17r_i^2v_c}{\mathrm{ln}\mathrm{\Lambda }GM}=2.64\times 10^{11}r_{i,2}^2v_{250}M_5^1(\mathrm{ln}\mathrm{\Lambda })_{10}^1\mathrm{yr}$$
(17)
where $`M_5`$ is again the cluster mass in units of $`10^5\mathrm{M}_{}`$. The friction time-scale thus scales with the square of the cluster’s initial radius in the potential, and is inversely proportional to its mass. It is the inner, most massive clusters which are affected first.
If we continue to model the inner parts of the host galaxy as an isothermal sphere with $`M_G(r)=v_c^2r/G`$ and use the virial theorem \[eq. (1)\] for the cluster, we can determine the radius at which the incoming cluster will dissolve as
$$r_{\mathrm{dis}}r_h\left(\frac{M_G(r_{\mathrm{dis}})}{M}\right)^{\frac{1}{3}}=\frac{r_hv_c}{\sqrt{7.5}\sigma _{}}=46r_5\sigma _{10}^1v_{250}\mathrm{pc}.$$
(18)
Young clusters formed in the high-density regions of starburst galaxies would thus contribute to the build-up of the nuclear bulge after being dragged inwards by dynamical friction and tidally shredded by the tidal field.
In some circumstances it may be possible that several young clusters are born close enough to eachother to tidally interact and even merge. To quantify this we use an approximate merging criterion fitted by Aarseth & Fall (1980) to the results of N-body merger simulations. For the escape velocity of the clusters at pericenter $`p`$ of their relative orbit we take an approximate expression assuming two overlapping Plummer spheres, $`v_e^2(p)=27.6\sigma _{}^2/(1+p^2/1.2r_h^2)^{1/2}`$ (see also the discussion in Gerhard & Fall 1983). Then the criterion for merging becomes
$$\left(\frac{p}{4r_h}\right)^2+\left(\frac{v_p}{6\sigma _{}}\right)^2\left(1+\frac{p^2}{1.2r_h^2}\right)^{\frac{1}{2}}\mathrm{}<1$$
(19)
where $`v_p`$ is the relative velocity at pericenter. Here we have used the virial equation (1), and $`\sigma _{}`$ is again the one-dimensional rms velocity dispersion of the cluster. For head-on collisions, this formula predicts merging for $`v_p\mathrm{}<6\sigma _{}=60\sigma _{10}\mathrm{km}\mathrm{s}^1`$ (slightly more than $`\sqrt{2}`$ times the rms escape velocity from each cluster), or $`\mathrm{\Delta }v=\sqrt{3627.6}\sigma _{}3\sigma _{}30\sigma _{10}\mathrm{km}\mathrm{s}^1`$ for their relative velocity at large separations. It also shows that merging requires the two clusters to come within several half-mass radii of eachother for merging to occur, at correspondingly smaller approach velocities. The most likely situation for this to happen would be when two clusters are born from the same giant molecular cloud complex.
### 3.7. Evolution of globular cluster systems
The evolution of globular cluster systems has recently been modelled in a number of studies, among others by Gnedin & Ostriker (1997), Murali & Weinberg (1997), Baumgardt (1998) and Vesperini (1998). These models combine assumptions about the initial cluster mass function and cluster locations with evolutionary models for individual clusters. In the models the various processes described above are considered, and treated in some studies by parametrized mass loss rates or analytic approximations to the results of N-body simulations, in others as diffusion terms in Fokker-Planck models. Some of the main results of these studies are:
(i) Globular cluster systems in galaxies evolve significantly. In the Milky Way the typical cluster destruction time is of order the age of the system, and about half of the present globulars will be destroyed in the next Hubble time.
(ii) Clusters in the inner regions of their host galaxy are disrupted most rapidly. Similarly, clusters on eccentric orbits are preferentially destroyed over clusters on tangential orbits.
(iii) Low-mass and high-concentration clusters are disrupted by evaporation, loosely bound clusters and those on central or eccentric orbits by tides, and massive inner clusters by dynamical friction and tides.
(iv) Low-mass clusters are destroyed most efficiently and initial power-law mass distributions tend to become transformed towards approximately log-normal mass distributions.
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# GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER
## 1 INTRODUCTION
The cold dark matter plus cosmological constant model ($`\mathrm{\Lambda }`$CDM) appears to be remarkably successful in describing cosmology on scales much larger than a Mpc (e.g. Bahcall, Ostriker, Perlmutter & Steinhardt 1999). However recent observations on smaller scales find $`\mathrm{\Lambda }`$CDM wanting: $`\mathrm{\Lambda }`$CDM predicts dwarf galaxy halos (e.g. Moore et al. 1999) that are denser and more concentrated than observed (Dalcanton & Bernstein 2000; Firmani et al. 2000b). The persistence of bars in galaxies like our own, suggests that the dark matter is not centrally concentrated (Sellwood 2000). In addition, the intensively studied CL0024+1654 is nearly spherical (Colley, Tyson & Turner 1996; Tyson, Kochanski & Dell’antonio 1998), at odds with $`\mathrm{\Lambda }`$CDM predictions of triaxial systems. Spergel & Steinhardt (2000) have recently suggested that these and other discrepancies (see for example Dav$`\stackrel{´}{\mathrm{e}}`$, Spergel, Steinhardt & Wandelt 2000; Wandelt et al. 2000 and references therein) could be resolved if cold dark matter were weakly self interacting with a large scattering cross-section (SIDM). Several authors (Burkert 2000; Yoshida, Springel, White & Tormen 2000; Dav$`\stackrel{´}{\mathrm{e}}`$ et al. 2000) have since undertaken N-body simulations of halos with self interactions in the regime suggested by Spergel & Steinhardt (2000). These works find that SIDM halos produce flatter, smoother and more spherical cores than their $`\mathrm{\Lambda }`$CDM counterparts.
We are interested in the statistics of multiple imaging due to a generalized NFW profile (Zhao 1996) which may be used to describe both CDM and SIDM halos. We concentrate our efforts on lensing by clusters of galaxies, since single galaxies contain a significant mass in baryons, distributed to produce a flat rotation curve. This makes SIDM in galaxies difficult to probe using gravitational lensing. Conversely, clusters of galaxies have a much larger mass to light ratio, and baryons which are less centrally concentrated (particularly in the absence of a central dominant galaxy). We assume spherically symmetric profiles and concentrate on the statistics of the optical depth to multiple imaging. The total magnification, and the maximum image splitting are also considered.
In Secs. 2 and 3 we discuss the properties of dark matter halos and introduce the formalism for computing the lensing statistics of the Zhao profile. In Secs. 4, 5 and 7 we discuss the statistics obtained, contrast the parameter degeneracies present in the lens statistics and in profile fitting, and present some implications for SIDM cluster halos. In Sec. 8 we discuss the important implications for lens statistics of the distribution of halo profile-parameters. We also discuss two potentially important caveats which turn out not to effect our conclusions: we discuss magnification bias (Sec. 6), and give a preliminary account of the effect of massive central galaxies on the multiple imaging rate (Sec. 9). Throughout the paper we have assumed a standard $`\mathrm{\Omega }_o=1`$, $`\mathrm{\Lambda }=0`$ and $`h=\frac{H_o}{100kmsec_1}=0.7`$ filled beam cosmology.
## 2 DARK MATTER PROFILES
Navarro, Frenk & White (1997) claimed the NFW profile (with an inner density $`\rho (r)`$ varying with radius $`r`$ as $`\rho (r)r^1`$) to be universal. More recently Moore et al. (1999) and Ghingna et al. (2000) have claimed a different universal profile, having $`\rho (r)r^{\frac{3}{2}}`$ in the innermost regions. However Jing & Suto (2000) find a systematic correlation of inner profile slope with mass. Moreover, they find some scatter in the slope at a fixed mass. On the other hand, from their convergence study, Klypin, Kravtsov, Bullock & Primack (2000) find that while there is some scatter, in general halos are equally well fit in the region where convergence is achieved by either an NFW or a Moore profile, independent of the halo mass. They interpret the correlation found by Jing & Suto (2000) as a trend in halo concentration rather than in asymptotic slope. Dav$`\stackrel{´}{\mathrm{e}}`$, Spergel, Steinhardt & Wandelt (2000) present a comprehensive suite of simulations of both CDM and SIDM dwarf galaxy halos. They find typical profiles for CDM halos in the inner most regions of $`\rho (r)r^{\frac{3}{2}}`$, with some scatter. However they found SIDM halos to be less cuspy, having central profiles ranging from $`\rho (r)r^1`$ to $`\rho (r)r^{\frac{1}{2}}`$. The uncertainty as to the optimal halo profile as well as the variation introduced by SIDM motivate an exploration of lensing by profiles that differ from the analytically treatable NFW form.
Following the introduction of the Navarro-Frenk-White profile (NFW) (Navarro, Frenk & White 1995,1996,1997), Bartelmann (1996) was the first to consider its lensing implications, showing that the profile is capable of producing radial arcs like those observed in several clusters (e.g. Mellier, Fort & Kneib 1993; Smail et al. 1995). There have been many subsequent efforts contrasting resultant lensing predictions with observation. Bartelmann et al. (1997) found the optical depth for large arcs for lensing by NFW profiles and found significant dependence on cosmology, though no models could reproduce the observed number. Wright & Brainerd (2000) and Asano (2000) have contrasted weak lensing mass measurements and 2-image system image positions flux ratios based on the NFW with those for a singular isothermal sphere (SIS). Maoz, Rix, Gal-Yam & Gould (1997) carried out a survey for wide separation lensed quasars, and compared their results with a thorough statistical study of the number of expected cluster lensed quasars. Modeling clusters using the NFW profile, they found that their null result for observed lensed quasars was consistent with the lensing statistics. Molikawa, Hattori, Kneib & Yamashita (1999) investigated giant arc statistics and found that the NFW profile can reproduce the number of observed lensed arcs. However their study finds some discrepancy between virial and X-ray temperature for some clusters. Broadhurst, Huang, Frye & Ellis (2000) have found that an NFW profile (singular core) can explain strong lensing in CL0024+1654, contrary to earlier claims of a flat core (Tyson et al. 1998). However Shapiro & Iliev (2000) find that the velocity dispersion implied by the NFW fit is a factor of two larger than that observed. The shallower cores resulting from SIDM simulations may resolve this discrepancy. On the other hand, Smith et al. (2000) find that radial arcs imply a core in A383 which has a steeper logarithmic slope than an NFW profile. Williams, Navarro & Bartelmann (1999) find that arc properties for the most massive clusters are consistent with NFW profiles derived from N-body simulations, but that clusters with velocity dispersions of $`1000kmsec^1`$ require the presence of a massive central galaxy to reconcile the models with observation. Fox, & Pen (2000) have considered lensing by galaxy groups modeled as NFW profiles in the Hubble Deep Field. They find that the null result for lenses is consistent with lensing probability. In addition, they compare the probability for multiple imaging due to an NFW profile with the probability for an SIS, and find that the probability is significantly reduced.
Recent complimentary studies have also considered lensing statistics for Zhao profiles. Molikawa & Hattori (2000) find that the ratio of radial to tangential arcs is a sensitive function of the inner slope of the density profile. Li & Ostriker (2000) consider a semi-analytic treatment of lensing by the Zhao profile. They find that lensing probability is sensitive to the profile considered, as well as cosmology. In addition, they find a probability for large splittings which is inconsistent with the observed distribution.
## 3 GRAVITATIONAL LENSING BY GENERALIZED NFW PROFILES
The NFW halo density profile can be generalized to describe a profile $`\rho (r)`$ as a function of radius $`r`$ with an arbitrary power-law shaped central cusp $`\rho (r)r^\beta `$ and outer regions that fall off as $`r^3`$ (Zhao 1996),
$$\rho (r)=\frac{\delta _c\rho _c}{\left(\frac{r}{r_s}\right)^\beta \left(1+\frac{r}{r_s}\right)^{3\beta }}$$
(1)
Navarro, Frenk and White (1996) defined the concentration parameter
$$C_{NFW}\frac{r_{200}}{r_s},$$
(2)
where
$$r_{200}\left(\frac{3M_{200}}{800\pi \rho _c}\right)^{\frac{1}{3}}$$
(3)
is the radius enclosing a mass $`M_{200}`$ with an over-density of 200 above the critical density. The characteristic over-density $`\delta _c`$ in Eqn. 1 is
$$\delta _c=\frac{200C_{NFW}^3}{3F(C_{NFW})}$$
(4)
where
$$F(y)_0^yx^{2\beta }(1+x)^{\beta 3}𝑑x.$$
(5)
Therefore, given a cluster mass $`M_{200}`$, the Zhao profile has two free parameters $`\beta `$ and $`C_{NFW}`$. $`C_{NFW}`$ is quoted for a profile at red-shift zero. Our lensing calculations assume a non-evolving cluster-lens population. The concentration parameter $`C_{NFW}(z)`$ and characteristic over-density $`\delta _c(z)`$ of a halo placed at red-shift $`z`$ are therefore the solution of
$$C_{NFW}(z):C_{NFW}^3(z)=\frac{F(C_{NFW}(z))}{F(C_{NFW})}\frac{C_{NFW}^3}{(1+z)^2(1+\mathrm{\Omega }_oz)}$$
(6)
and
$$\delta _c(z)=\frac{\delta _c}{(1+z)^2(1+\mathrm{\Omega }_oz)}.$$
(7)
These transformations leave $`r_s`$ and $`\rho _c\delta _c`$ independent of $`z`$. From the study of the evolution of a large sample of $`\mathrm{\Lambda }`$CDM halos, Bullock et al. (2000) find $`C_{NFW}(\beta =1)(1+z)^1`$, consistent with Eqn. 6. Note however, that the inclusion of a finite self-interaction cross-section should result in a different profile evolution.
The Zhao profile is spherically symmetric, while real clusters, and N-body simulations tend to be elliptical or tri-axial. Blandford & Kochanek (1987) and Kochanek & Blandford (1987) find that the the introduction of ellipticities smaller than $`0.2`$ into nearly singular profiles has little effect on the cross-section and image magnification. Furthermore, the images remain on opposite sides of the potential so that the variation in image separations is similar to the axis ratio of the potential. However profiles containing a finite core and ellipticity show a significant increase in cross-section and qualitative differences in image positions.
We make the thin screen approximation for gravitational lensing by extended bodies. The surface mass density $`\mathrm{\Sigma }(\xi )`$ at a radius $`\xi `$ may be written
$$\mathrm{\Sigma }(\xi )=2\delta _c\rho _cr_sx^{1\beta }_0^{\frac{\pi }{2}}sin\theta (sin\theta +x)^{\beta 3}𝑑\theta ,$$
(8)
where $`x\frac{\xi }{r_s}`$. The mass $`M(\xi )`$ enclosed within the cylinder of radius $`\xi `$ and the bend angle $`\alpha (\xi )`$ of a ray with impact parameter $`\xi `$ are
$$M(\xi )=2\pi r_s^2_0^xx^{}\mathrm{\Sigma }(x^{})𝑑x^{}.$$
(9)
and
$$\alpha (\xi )=\frac{4GM(\xi )}{c^2\xi }.$$
(10)
If $`\beta =1`$, Eqns. 8, 9 and 10 reduce to the expressions computed by Bartelmann (1996).
The lens equation describes the location $`\xi `$ of an image in the lens plane given a source position $`\eta `$ (both measured from the observer source line of sight). The angular diameter distances from the observer to a lens at red-shift $`z_d`$, from the observer to a source at $`z_s`$ and from the lens to the source are denoted by $`D_d`$, $`D_s`$ and $`D_{ds}`$. With these, the lens equation may be written (e.g. Schneider, Ehlers & Falco 1992)
$$0=\theta \beta \alpha (\theta )\frac{D_{ds}}{D_s},$$
(11)
where $`\theta \frac{\xi }{D_d}`$ and $`\beta \frac{\eta }{D_s}`$ are the angular positions of the image and source respectively. The point source magnification $`\mu _i`$ of an image at $`\theta _i`$ is
$$\mu _i=\frac{\theta _i}{\beta }\left(\frac{1}{1\frac{D_{ds}}{D_s}\frac{d\alpha (\theta _i)}{d\theta }}\right).$$
(12)
The Zhao profile has a critical radius $`\eta _{crit}`$ separating regions of the source plane that produce 1 or 3 images. The corresponding critical impact parameter $`\xi _{crit}`$ is the solution of
$$\xi _{crit}:\frac{4G}{c^2}\left(2\pi \mathrm{\Sigma }(\xi )\frac{M(\xi )}{\xi }\right)\frac{D_s}{D_dD_{ds}}=0$$
(13)
and the critical angle for multiple imaging is
$$\beta _{crit}=\frac{\xi _{crit}}{D_d}\alpha (\xi )\frac{D_{ds}}{D_s}.$$
(14)
For later comparison, the bend angle for a SIS is
$$\alpha ^{SIS}=\frac{2\pi GM_{200}}{c^2r_{200}}.$$
(15)
In this case $`\xi _{crit}^{SIS}=0`$ so that
$$\beta _{crit}^{SIS}=\frac{2\pi GM_{200}}{c^2r_{200}}\frac{D_{ds}}{D_d}.$$
(16)
For the computation of lens statistics for an unevolving population of lenses with constant co-moving number density $`n_d=n_o(1+z)^3`$ we find (Turner, Ostriker & Gott 1984) that for an $`\mathrm{\Omega }=1`$ universe the differential and total cross-section may be written
$$d\tau (z_d,z_s)=\pi n_o\frac{c}{H_o}D_d^2(z_d)\beta _{crit}^2(z_d,z_s)\sqrt{1+z_d}dz_d$$
(17)
and
$$\tau (z_s)=\pi n_o\frac{c}{H_o}_0^{z_s}D_d^2(z_d)\beta _{crit}^2(z_d,z_s)\sqrt{1+z_d}𝑑z_d.$$
(18)
Turner, Ostriker & Gott (1984) defined a dimensionless quantity $`F_{SIS}`$ that measures the efficiency of cosmologically distributed SISs for producing multiple images.
$$F_{SIS}4\pi ^3n_o\left(\frac{c}{H_o}\right)^3\left(\frac{GM_{200}}{r_{200}c^2}\right)^2.$$
(19)
We use $`F_{SIS}`$ to normalize $`\tau `$. More realistic calculations of $`\tau (z_s)`$ for the NFW ($`\beta =1`$) profile have previously been made within the contexts of wide separation quasars (Maoz, Rix, Gal-Yam & Gould 1997), giant arc statistics (Bartelmann et al. 1997; Molikawa, Hattori, Kneib & Yamashita 1999) and distant supernovae (Porciani & Madau 2000). However in this work we are interested in relative values of $`\tau `$ produced by different profiles, and so have neglected quantitatively important effects such as formation history.
One image is present for all source positions. We denote this image to be at $`\theta _1`$. Note that $`\theta _1>\beta _{crit}`$ for all $`\beta `$. For $`\beta <\beta _{crit}`$ there are two additional images with positions $`\theta _2<\beta _{crit}`$ and $`\theta _3>\beta _{crit}`$. We label the corresponding point source image magnifications as $`\mu _1`$, $`\mu _2`$ and $`\mu _3`$ respectively. If $`\beta _{crit}>0`$ then $`\theta _1>0`$, $`\theta _1,\theta _2<0`$, $`\mu _1,\mu _3>0`$ and $`\mu _2<0`$. In addition $`|\mu _3||\mu _2|`$. For $`\beta <\beta _{crit}`$, we define the average image separation of a multipli-imaged source at $`z_s`$ due to a deflector at $`z_d`$ to be
$$\mathrm{\Delta }\theta (z_d,z_s)=\frac{2}{\beta _{crit}^2}_0^{\beta _{crit}}\beta (\theta _1(\beta ,z_d,z_s)\theta _3(\beta ,z_d,z_s))𝑑\beta .$$
(20)
Similarly the average magnification of a multipli-imaged source at $`z_s`$ due to a deflector at $`z_d`$ is
$$\mu (z_d,z_s)=\frac{2}{\beta _{crit}^2}_0^{\beta _{crit}}\beta (|\mu _1(\beta ,z_d,z_s)|+|\mu _2(\beta ,z_d,z_s)|+|\mu _3(\beta ,z_d,z_s)|)𝑑\beta .$$
(21)
It follows that the average image splittings $`\mathrm{\Delta }\theta (z_s)`$ and total magnifications $`\mu (z_s)`$ for a source at $`z_s`$ are
$$\mathrm{\Delta }\theta (z_s)=\frac{1}{\tau (z_s)}_0^{z_s}\mathrm{\Delta }\theta (z_d,z_s)\frac{d\tau (z_d,z_s)}{dz_d}𝑑z_d$$
(22)
and
$$\mu (z_s)=\frac{1}{\tau (z_s)}_0^{z_s}\mu (z_d,z_s)\frac{d\tau (z_d,z_s)}{dz_d}𝑑z_d.$$
(23)
## 4 CROSS-SECTION, IMAGE SPLITTINGS AND MAGNIFICATION
In this section we present the differential and total cross-sections (Eqns. 17 and 18), distributions of image separations, and distributions of magnifications (Eqns. 22 and 23) for Zhao profiles having different parameters $`\beta `$ and $`C_{NFW}`$.
Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER shows $`\rho (r)`$ (top) and $`\mathrm{\Sigma }(\xi )`$ (bottom) for three different Zhao profiles. These are labeled profile $`I`$ (solid lines), $`II`$ (dashed lines) and $`III`$ (dotted lines), have parameters of $`(\beta ,C_{NFW})=(1.5,20)`$, $`(1.0,15)`$, and $`(0.5,10)`$ respectively, and are increasingly steep cusps in their central regions in addition to being increasingly centrally concentrated. The model clusters have an $`M_{200}`$ equal to that of a SIS ($`M_{200}^{SIS}`$) with a velocity dispersion corresponding to a large cluster $`\sigma _{SIS}^{cstr}=1000kmsec^1`$:
$$M_{200}^{SIS}=(\sigma _{SIS}^{cstr})^3\left(\frac{3}{800\pi \rho _c}\right)^{\frac{1}{2}}\left(\frac{2}{G}\right)^{\frac{1}{3}}.$$
(24)
The solid light lines show $`\rho _{SIS}(r)`$ and $`\mathrm{\Sigma }_{SIS}(\xi )`$ of the SIS for comparison. The horizontal dashed lines show the critical surface mass density
$$\mathrm{\Sigma }_{crit}=\frac{c^2}{4\pi G}\frac{D_s}{D_dD_{ds}}$$
(25)
for a source at $`z_s=3`$ and lenses at $`z_d=0.5`$, $`1.0`$ and $`1.5`$. Comparison of $`\mathrm{\Sigma }_{crit}`$ with $`\mathrm{\Sigma }(\xi )`$ indicates a significant variation of the lensing probability between profiles having the same $`M_{200}`$ but different profile parameters. Profiles $`I`$ and $`II`$ reach the critical density over a much larger area than profile $`III`$, which is only just supercritical at $`z_d=1.0`$, and sub-critical at $`z_d=2.0`$.
The variation of cross-section with the profile parameters is quantified by Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER which shows $`\tau (z_s)`$ for Zhao profiles $`IIII`$. The cross-section is given in units of $`F_{SIS}`$ (Eqn. 19). For comparison, the cross-section for a SIS ($`\tau _{SIS}(z_s)`$, light line, Turner, Ostriker & Gott 1984) is also shown in Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER. Examples of the effect of different cosmologies on $`\tau _{SIS}(z_s)`$ (and on $`\frac{1}{\tau (z_s)}\frac{d\tau (z_s,z_d)}{dz_d}`$ in Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER) can be found in Turner et al. (1984) and Turner (1990). There is a striking variation in $`\tau `$ over these three profiles. Profiles $`II`$ and $`III`$ produce $`\tau `$s which differ by an order of magnitude, while there are 4 orders of magnitude between $`\tau `$ for profiles $`I`$ and $`III`$. More concentrated profiles are much more efficient for producing multiple images due to their higher central densities. For larger source red-shifts the SIS has a similar cross-section to profile $`II`$, while at smaller red-shifts the cross-section is similar to that of profile $`I`$, indicating that different profiles are efficient lenses at different red-shifts.
In Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER we have plotted $`\frac{1}{\tau (z_s)}\frac{d\tau (z_s,z_d)}{dz_d}`$ for profiles $`I`$ (solid lines), $`II`$ (dashed lines) and $`III`$ (dotted lines) with $`z_s=1.5`$ (upper panel) and $`z_s=3.0`$ (lower panel). The light lines show the corresponding distributions for a SIS. Less concentrated profiles have distributions that are more sharply peaked, with modes at lower red-shift (c.f. Kochanek & Blandford 1987). The lower cross-section to multiple imaging is therefore due to a combination two effects. Firstly, the halo reaches the critical density to lensing only over a small region in the center due to the low concentration or central slope. In addition, there is a decreased redshift range in which the central surface density exceeds the critical value anywhere. Note that the assumption in Eqns. 17 and 18 of an unevolving population of SIDM halos is aided in the case of low-concentration profiles since these preferentially lens at low red-shift.
To further investigate the variation of cross-section with $`\beta `$ and $`C_{NFW}`$ we have plotted contours of $`\frac{\tau (z_s)}{\tau _{SIS}(z_s)}`$ with $`z_s=3`$, over a parameter space of $`\beta `$ and $`C_{NFW}`$ (Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER). Two additional dark lines are shown. The upper dark line is a contour of unity and marks the parameters that produce an cross-section equal to that of a SIS. The lower dark line marks the boundary of the region of parameter space where no multiple imaging is produced (on this and subsequent plots this line is the approximated by the $`10^5`$ contour). The contours describe sets of parameters $`\beta `$ and $`C_{NFW}`$ that are equally efficient for producing multiple images. Profiles $`I`$, $`II`$ and $`III`$ lie on a line in $`\beta C_{NFW}`$ space that is approximately orthogonal to the contours of $`\tau (z_s)`$. However, figure GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER demonstrates that dramatic variability in $`\tau (z_s)`$ may be obtained by variation of either $`\beta `$ or $`C_{NFW}`$.
We have calculated $`\frac{\tau (z_s)}{\tau _{SIS}(z_s)}`$, $`\mathrm{\Delta }\theta (z_s)`$ and $`\mu (z_s)`$ over a small grid of $`\beta `$ and $`C_{NFW}`$ (see Tab. 1). We find that profiles which produce a small cross-section to multiple imaging also produce smaller image splittings, although only by a factor of a few compared with more highly concentrated profiles. However, the image magnification is significantly larger for the less concentrated profiles. Distributions of image separations ($`\frac{dP}{d\mathrm{\Delta }\theta (z_s)}`$) for the Zhao profile are presented in Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER for the values of $`\beta `$ and $`C_{NFW}`$ in Tab. 1. For comparison, the light line shows the distribution of splitting angles for a SIS ($`\frac{dP_{SIS}}{d\mathrm{\Delta }\theta (z_s)}`$). There is a trend for the distributions with lower $`\mathrm{\Delta }\theta (z_s)`$ to also have a narrower range of $`\mathrm{\Delta }\theta (z_s)`$. In contrast to the SIS distribution, the Zhao profile distributions are asymmetric about the mode, having a longer large splitting tail. Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER shows magnification distributions $`\frac{dP}{d\mu }`$ for the values of $`\beta `$ and $`C_{NFW}`$ in Tab. 1. The distribution for an SIS,
$$\frac{dP^{SIS}}{d\mu }=\frac{8}{\mu ^3},$$
(26)
is shown for comparison (light line). From Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER we see that the high mean magnifications shown in Tab. 1 result from a higher minimum magnification of multiple images. Lenses with low concentration and small $`\beta `$ result in a low $`\tau (z_s)`$. Very close alignment is therefore required to produce multiple imaging, and the close alignment results in high magnification. More concentrated lenses do not require such high alignment for multiple imaging. This results in a larger cross-section and reduced mean magnification.
The basic lensing properties of the Zhao profile may therefore be summarized as follows. Profiles that have a small concentration and flat central profile produce significantly less multiple imaging than a profile with either a high concentration or steep central profile (or both). However, these few multiple images are significantly magnified. In any lens survey, the increased $`\mu `$ will partially offset the decreased $`\tau (z_s)`$ in determining the lens frequency. This magnification bias is discussed in Sec. 6. For clusters of galaxies, the reduced splittings for low density profiles are $`10^{\prime \prime }`$ and should not pose a significant selection bias for well defined surveys.
## 5 THE $`C_{NFW}\beta `$ DEGENERACY
A property of the Zhao profile (Eqn. 1) is the ability to describe a high central density either by a steep central cusp, or by a large concentration parameter (so that the turnover separating the core regions from the outer halo occurs closer to the center). This degeneracy is present between other profiles exhibiting a power-law core and a characteristic radius. Klypin, Kravtsov, Bullock & Primack (2000) have investigated the relative suitability of Moore et al. (1999) and NFW ($`\beta =1`$) profiles as descriptions of CDM halos. For galaxy size halos, they find differences which are small at radii above $`0.01r_{200}`$. Larger differences are found for lower concentration cluster halos as the profiles central slope begins to approach its asymptotic value. Moreover, using a very general 4-parameter NFW profile, Klypin et al. (2000) obtained equally good fits to an N-body halo for several different sets of parameters. Our interest lies in the relationship between the $`C_{NFW}\beta `$ degeneracy found in in the contexts of multiple imaging statistics and profile fits to simulated halos.
Our approach is to select a collection of profiles $`\rho _o(r)`$ having $`\beta =0.5`$ but different concentrations $`C_{NFW,o}`$. Keeping $`M_{200}`$ constant (for comparison with our strong lensing calculations), we fit a new profile $`\rho _{fit}(r)`$ with $`C_{NFW}`$ as a function of $`\beta `$. Fits are obtained by minimizing each of three different quantities which emphasize different parts of the halo:
$$\chi ^2=_{\frac{r_{200}}{100}}^{r_{200}}r^2(\rho (r)\rho _o(r))^2𝑑r,$$
(27)
$$\mathrm{\Delta }\rho _{max}=max(|log(\rho (r))\mathrm{log}(\rho _o(r))|)$$
(28)
(Klypin, Kravtsov, Bullock & Primack 2000), and
$$\chi _{rel}^2=_{\frac{r_{200}}{100}}^{r_{200}}r^2\left(\frac{\rho (r)\rho _o(r)}{\rho _o(r)}\right)^2𝑑r.$$
(29)
Minimization is performed over the range $`0.01r_{200}r_{200}`$.
Some examples of the fitted profiles and their fractional residuals are presented in Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER (light lines). The dark lines show the profile being fitted. In each case we quote the parameters $`\beta _o`$, $`C_{NFW,o}`$ and $`\beta `$, $`C_{NFW}`$, as well as the average absolute relative difference over the fitting volume
$$\frac{\mathrm{\Delta }\rho }{\rho }=\frac{3}{r_{200}^3\left(\frac{r_{200}}{100}\right)^3}_{\frac{r_{200}}{100}}^{r_{200}}r^2\frac{|\rho _{fit}(r)\rho _o(r)|}{\rho _o(r)}𝑑r.$$
(30)
The central density dominates minimization of $`\chi ^2`$ (top panels) so that the fitted profile has the correct central density at the cost of error over large regions around $`r=\frac{r_{200}}{10}`$. Alternatively if we set $`C_{NFW}`$, and adjust $`\beta `$, the fitted profile has the correct central slope and density, but is systematically high or low over the outer regions. The initial and fitted halos have average absolute relative differences of a few to $`15\%`$. As noted by Klypin, Kravtsov, Bullock & Primack (2000), minimization of $`\mathrm{\Delta }\rho _{max}`$ (central panels) often does a much better job of reproducing the density at small $`r`$. If $`\beta `$ is the fitted parameter rather than $`C_{NFW}`$, a similar result is obtained, the central slope adjusts so as to give the correct density at small $`r`$. Minimization of $`\mathrm{\Delta }\rho _{max}`$ produces halo fits with similar average absolute relative differences to $`\chi ^2`$. Finally, minimizing $`\chi _{rel}^2`$ (lower panels) obtains good fits beyond $`r0.1r_{200}`$ at the expense of an incorrect central density. The mean relative error is substantially reduced, with average absolute relative differences of a fraction of a percent.
The dot-dashed lines in the upper, central and lower plots of Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER are depictions of the $`\beta C_{NFW}`$ degeneracies corresponding to the minimization quantities $`\chi ^2`$, $`\mathrm{\Delta }\rho _{max}`$ and $`\chi _{rel}^2`$. Since $`\chi ^2`$ produces good fits to the core region, which is known to be most important for strong lensing (e.g. Turner, Ostriker & Gott 1984), the lines of degeneracy are nearly parallel to those describing the cross-section to multiple imaging over much of the parameter space. However for combinations of low or high concentration and low inner slope, the approximately degenerate profiles do not provide a accurate description of the lensing statistics. Minimization of $`\mathrm{\Delta }\rho _{max}`$ results in profile degeneracies similar to those for strong lensing unless the profile has large $`C_{NFW}`$ and small $`\beta `$. Finally, minimization of $`\chi _{rel}^2`$ leads to degeneracies that differ significantly from those for strong lensing. Fits that are weighted towards the central regions, and which result in highly concentrated profiles may therefore be converted to the analytically treatable NFW ($`\beta =1`$) profile while approximately preserving the lensing statistics. This is particularly true for fits based on minimization of $`\mathrm{\Delta }\rho _{max}`$. On the other hand, lensing calculations based on profiles which best describe the density over the largest volume (out to $`r_{200}`$) may severely under or over-estimate the strong lensing statistics. For example, Zhao halos having the same $`M_{200}`$, but different parameters ($`\beta ,C_{NFW}`$) = (1.5,6.0) and (1.0,10.0) are degenerate to fitting by the minimization of $`\chi _{rel}^2`$, but differ in $`\tau `$ by an order of magnitude. Thus Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER may be considered a quantification (for Zhao profiles) of the statement that lensing is sensitive to the size of the halo core.
## 6 MAGNIFICATION BIAS
Statistics generated from surveys for gravitational lenses are subject to several classes of bias. Kockanek (1991) has given a detailed discussion of selection biases in optical, imaging, galaxy mass gravitational lens surveys. Two primary sources of bias against finding a lens are selection against small separations due to the finite angular resolution of the survey, and the obscuration or differential reddening of one or both images due to dust in the lens galaxy. When considering lensing by clusters, the splittings are nearly always larger than the survey resolution. In addition, the lensed images are viewed through the halo rather than the outskirts of a galaxy. These biases should not therefore operate on any significant level in the analysis of cluster lensing statistics. The magnification of a lensed image has the potential to boost the observed flux of a source which is intrinsically faint above the survey detection limit. Lensed images therefore sample fainter members of the source population. This is magnification bias (Gott & Gunn 1974; Turner 1980), and it a significant factor in any lens survey.
As mentioned in Sec. 4, lenses with low concentration and a shallow cusp result in a low cross-section. Very close alignment is therefore required to produce multiple imaging, and this results in a high minimum magnification. More concentrated lenses do not require such high alignment to produce multiple imaging, thus the cross section is larger, and the minimum magnification not so high. In the latter case, although the mean is lower, high magnifications are still present where the alignment is high. As a result, the mean magnification is not necessarily a fair indication of the magnification bias.
Magnification bias operates differently in surveys for galaxy and cluster lenses. The typical separation of a galaxy lens is smaller than most optical survey resolution limits, and the total magnification must be used to compute the magnification bias. However, the total magnification is not the important quantity for cluster lensing surveys since the separate images are generally resolved (see also Chiba & Yoshii 1997). There are different approaches to making a wide separation gravitational lens survey. Firstly, candidate lenses may be selected from galaxy or cluster catalogues having one image of a source in good alignment with a potential lens. Follow up imaging can then be used to look for a second, fainter image. This is the approach taken by systematic surveys to date (e.g. Maoz et al. 1997; Phillips, Browne & Wilkinson 2000). Here the magnification bias which results in the inclusion of the brighter image into the initial survey is important (though the overall bias for the survey must also consider the detectability of the faint image in follow up observations).
A second method looks for pairs of lensed images in existing data. In this instance, the magnification bias which also includes the fainter image (the second brightest for three image lenses, i.e. image 3), and thus both images into the survey is the relevant quantity.
To compute magnification bias, a number magnitude relation for the sources is required. We use the function described by Kochanek (1996), based on data from Boyle, Shanks & Peterson (1988) and Hartwick & Schade (1990). He uses a broken-power-law form (Turner 1990). To simplify our calculations and to keep the conclusions as general as possible, we assume that the survey depth is deeper than the break B-magnitude at $`z_s=3`$ of $`m_B19`$. The magnitude-number relation may then be written as
$$\frac{dN}{dm}=N_o10^{\alpha _mm}$$
(31)
where $`N_o`$ is a normalizing constant and $`\alpha _m`$ is the logarithmic slope. For a survey with a limit $`m_{lim}`$, the magnification bias is defined as (Fukugita & Turner 1991)
$$B\frac{_0^{\mathrm{}}\frac{dP}{d\mu }N(<m_{lim}+\frac{5}{2}log_{10}(\mu ))d\mu }{N(<m_{lim})}.$$
(32)
In application to a real lens survey, the magnification bias may be written more conveniently as
$$B\frac{_0^{\mathrm{}}_{\mu _{min}}^{\mu _{max}}d\mu _b\frac{dP}{d\mu _b}N(<m_{lim}+\frac{5}{2}log_{10}(\mu _b))d\mu _b}{N(<m_{lim})},$$
(33)
where the limits $`\mu _{min}`$ and $`\mu _{max}`$ are the range of the bright image magnification $`\mu _b`$ such that at least one other image is visible. Note that only the magnification distribution for the bright image is directly relevant. In this paper we adopt the simpler and non-specific case of magnification bias for bright and faint images based on Eqn. 32.
Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER shows the magnification distributions for the bright (left hand panel) and faint (right hand panel) images for Zhao profiles. The light lines show the distributions for a SIS:
$$\frac{dP^{SIS}}{d\mu _f}=\frac{2}{(\mu _f+1)^3}for\mu _f>0$$
(34)
and
$$\frac{dP^{SIS}}{d\mu _b}=\frac{2}{(\mu _b1)^3}for\mu _b>2$$
(35)
for comparison. The image magnification distributions for the Zhao profiles exhibit similar behavior to the SIS, and there is a high probability that the two images will have different magnifications. However unlike the SIS, the Zhao profiles have finite $`\frac{d\alpha (\theta )}{d\theta }`$ as $`\theta 0`$. The faint image therefore has a finite minimum magnification.
The magnification bias described by Eqn. 32 describes the additional number of lensed pairs. If the quantity of interest is the density of lensed images on the sky, then an additional factor of $`\mu ^1`$ must be inserted under the integral to account for the smaller region of source plane sampled per unit area of lens plane. Inserting Eqn. 31 into Eqn. 32 we find
$$B_b=_0^{\mathrm{}}\frac{dP}{d\mu _b}\mu _b^{\frac{5}{2}\alpha _m}𝑑\mu _b$$
(36)
for magnification bias due to bright images, and
$$B_f=_0^{\mathrm{}}\frac{dP}{d\mu _f}\mu _f^{\frac{5}{2}\alpha _m}𝑑\mu _f$$
(37)
for magnification bias due to faint images. Note that the magnification bias is equal to the mean magnification for $`\alpha _m=0.4`$. Since the high magnification tail of any magnification distribution has
$$\frac{dP}{d\mu }\mu ^3,$$
(38)
the integral in Eqns. 36 and 37 is finite for $`\alpha _m<\frac{4}{5}`$, and in this range no low-magnitude cutoff need be assumed.
Using Eqns. 34 and 35 the magnification biases for the faint and bright images of the SIS are
$$B_b^{SIS}=_2^{\mathrm{}}\frac{2}{(\mu _b1)^3}\mu _b^{\frac{5}{2}\alpha _m}𝑑\mu _b.$$
(39)
and
$$B_f^{SIS}=_0^{\mathrm{}}\frac{2}{(\mu _f+1)^3}\mu _f^{\frac{5}{2}\alpha _m}𝑑\mu _f.$$
(40)
Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER shows the absolute magnification bias (left hand panel), and the magnification bias relative to the SIS (right hand panel) based on the bright image as a function of $`\alpha _m`$. Plots are shown for the values of $`\beta `$ and $`C_{NFW}`$ in Tab. 1. The light line in the left hand panel shows the magnification bias curve for the SIS. Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER shows the corresponding results for the faint image. Values of absolute and relative (to the SIS) magnification bias based on the faint and bright image are presented in Tabs. 3, 3, 5 and 5 for $`\alpha _m=0.2,0.4,0.6`$. For quasars Kockanek (1996) finds $`0.27\pm 0.07`$.
The absolute magnification biases for the bright images range up to a few, a few tens and a few hundred for $`\alpha _m=0.2,0.4`$ and $`0.6`$ respectively. However, as discussed at the top of this section, the quantity of interest is the relative magnification bias. The values in Tab. 3 demonstrate the reduced amplitude of the relative (to the SIS) bias, which is similar for $`\alpha _m=0.2`$, but reduced by a factors of $`3`$ and $`5`$ for $`\alpha _m=0.4`$ and $`0.6`$. Similar behavior is seen in magnification bias of faint images. The absolute magnification biases for the faint images produced by Zhao profiles are slightly reduced from those computed for the bright image. However the magnification biases for the faint images relative to the SIS are increased by a factor of two over the corresponding bright image case. This apparent contradiction arises because the faint image produced by a SIS may be arbitrarily faint.
The magnification bias results in additional sources being available for lensing. Since the relative magnification bias is greater than 1 where the ratio of cross-sections is less than one, the magnification bias counteracts the dramatic dependence of $`\tau `$ on $`\beta `$ and $`C_{NFW}`$ seen in Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER. The relevant quantity is therefore the optical depth to multiple imaging $`\tau (z_s)B_b`$ or $`\tau (z_s)B_f`$. Tabs. 7 and 7 show values of the optical depth relative to the SIS for biases computed using bright and faint images. Values are shown assuming $`\alpha _m=0.2`$, $`0.4`$ and $`0.6`$. The magnification bias reduces the difference between cross-sections for different profiles, though the variation with $`\beta `$ and $`C_{NFW}`$ is still very dramatic. However, profiles with no cross-section to multiple imaging are not subject to magnification bias. Magnification bias therefore increases the dependence of the cross-section to multiple imaging on $`\beta `$ and $`C_{NFW}`$ near the multiple-imaging / no-multiple-imaging boundary.
## 7 IMPLICATIONS FOR CDM AND SIDM PROFILES
Miralda-Escude (2000) argued that strong lensing was a powerful test of SIDM models. Assuming a cross-section independent of velocity, he pointed out that if SIDM were responsible for the cores in dwarf galaxies, then scaling arguments lead to clusters cores that are larger than observed. While Miralda-Escude used a scaling law that is not a good fit to the SIDM simulations (Dav$`\stackrel{´}{\mathrm{e}}`$, Spergel, Steinhardt & Wandelt 2000), his basic argument, that lensing is a powerful test of the model remains valid. Furthermore, Miralda-Escude (2000) points out that the mis-alignment of the images in the cluster lens MS2137-23 cannot be reconciled with a spherical lens implied by SIDM halos.
Recently, Yoshida, Springel, White & Tormen (2000) have produced simulations of a cluster with $`M_{200}=7.4\times 10^{14}M_{}`$. They model the cluster using collision-less cold dark matter (labeled S1), as well as weakly self interacting cold dark matter having collision cross-sections of $`\sigma ^{}=0.1`$, 1.0 and $`10.0cm^2gm^1`$ (labeled S1W-a, S1W-b, and S1W-c). Halo density profiles are presented at $`z=0`$ for all 4 realizations. In addition, they present halo profiles for the simulation S1W-c at a range of different epochs. During the evolution of the halo, the central densities are most extreme when the cosmological expansion factor $`a`$ (normalized to its present value) takes the values $`a=0.78`$ and $`a=1.72`$. The 6 halo profiles are reproduced in Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER (dots). To estimate the lensing rate due to these different profiles we have fitted each with the Zhao profile (Eqn. 1). The fitting was achieved via minimization of $`\chi ^2`$ (solid lines) and $`\mathrm{\Delta }\rho _{max}`$ (dashed lines). The constraint that $`\beta 0.2`$ was imposed and the halos S1W-c ($`a`$=0.78) and S1W-c ($`a`$=1.72) were fitted assuming an $`r_{200}`$ (and hence $`M_{200}`$) corresponding to that at $`a`$=1. Yoshida et al. (2000) note that the halo approximately doubled its $`M_{200}`$ between $`a=0.73`$ and $`a=1.72`$. If half this gain was made before (after) $`a=1`$, then $`C_{NFW}`$ for the halo S1W-c ($`a`$=0.78) (S1W-c ($`a`$=1.72)) should be revised downwards (upwards) by $`10\%`$. Both the fits and fractional residuals are plotted in Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER. The Zhao profile provides an adequate fit to halos S1 and S1W-a. However for the larger interaction cross-sections, the core density is overestimated by the fitted profile. The corresponding lensing rate may therefore also be overestimated in these cases.
The estimates for the profile parameters $`C_{NFW}`$ and $`\beta `$ of the Yoshida, Springel, White & Tormen (2000) halos are plotted over contours of $`\frac{\tau (z_s)}{\tau _{SIS}(z_s)}`$ in Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER. Each parameter set is labeled and the upper and lower plots show points corresponding to the fits obtained by minimization of $`\chi ^2`$ and $`\mathrm{\Delta }\rho _{max}`$ respectively. From Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER we find that $`\tau `$ falls off very rapidly for small $`\beta `$ and $`C_{NFW}`$, the region of interest for SIDM halos. The halos S1W-b and S1W-c at the current epoch do not produce any multiple images. The halo S1W-a (at the current epoch) lenses at a significantly reduced rate with respect to the CDM halo S1. The variation in the density profile at different epochs is due to successive mergers (Yoshida et al. 2000). The spread in parameters obtained from profile fitting is representative of the range of lensing strengths expected. A population of S1W-c ($`a=1.72`$) halos would yield a finite lensing rate, comparable with that of S1W-a halos, while at the other epochs considered ($`a=0.78`$ and $`a=1.0`$) the halos would produce no multiple images.
For comparison, we have also included in Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER, points corresponding to parameters found by two complimentary studies of CDM halos. Jing & Suto (2000) have estimated of the spread in the distribution of concentration parameter for CDM halos by fitting the Zhao form to 4 cluster mass halos ($`M_{200}35\times 10^{14}M_{}`$) using both $`\beta =1`$ and $`\beta =1.5`$. Their results are plotted in the lower panel of Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER (labeled 1-4 as per Jing & Suto (2000)). This range is in agreement with the 68% range for $`C_{NFW}`$ (extrapolated to $`M_{200}10^{15}M_{}`$) found by Bullock et al. (2000) (plotted in the upper panel of Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER). The halo S1 has a $`C_{NFW}`$ which is consistent with the range of values found by these authors. It is important to note that the range of parameters describe halos that produce lensing cross-sections differing by several orders of magnitude. The lensing statistics will therefore be dominated by the more concentrated halos, indicating that the medians of profile parameter distributions may not provide a reasonable estimation of the lensing rate for a particular dark matter model. We return to this point in the following section.
More detailed modeling including evolution, non-spherical mass distributions, and a range of formation red-shifts, as well as observational details including magnification bias etc. (e.g. Turner, Ostriker & Gott 1984; Kochanek 1991) appropriate for surveys such as that being undertaken by the SDSS collaboration (e.g. York et al. 2000), is required before quantitative constraints on SIDM models can be derived. In addition, since lensing probes the central regions of the cluster, the contribution to the potential of a centrally dominant galaxy may be significant and must also be considered. We offer a preliminary discussion of the contribution of central galaxies in Sec. 9, however, we conclude that the sensitivity of $`\tau `$ to $`\beta `$ and $`C_{NFW}`$ indicates that cluster lensing statistics will provide a powerful probe of the self-interaction cross-section.
While lensing constrains the dark matter properties for dark matter collisions with velocities of $`v10^3kmsec^1`$, observations of dwarfs constrain dark matter properties with $`v30kmsec^1`$. If the dark matter particle is a composite particle, analogous to a nucleon with a low energy resonance, then its cross-section scales as $`\sigma ^{}v=const.`$ For fixed cross-sections on the dwarf scale, a cross-section that is inversely proportional to the characteristic velocity will produce significantly smaller, less spherical cores in clusters and large cores in dwarfs (Yoshida, Springel, White & Tormen 2000), consistent with observations of cluster cores like CL0024+1654 and CL0016+16 (Avila-Reese, Firmani, D’Onghia & Hernandez 2000; Firmani et al. 2000). Alternatively, if the dark matter can self annihilate at late times (Kaplinghat, Knox & Turner 2000) then $`\sigma ^{}v=const.`$, this leads to low density tri-axial cores and differing lensing predictions.
## 8 DEPENDENCE OF OPTICAL DEPTH ON THE HALO PROFILE DISTRIBUTION
Recent studies of N-body simulations have described the distribution of halo concentrations (Bullock et al. 2000; Jing & Suto 2000; Dav$`\stackrel{´}{\mathrm{e}}`$, Spergel, Steinhardt & Wandelt 2000). Moreover, the simulations of Yoshida, Springel, White & Tormen (2000) demonstrate large variation in the concentration of a SIDM halo during its evolution. An important quantity which has hitherto been neglected in parametric cluster lens analyses is the distribution of halo profiles.
The upper panels of Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER show cross-sections $`\frac{\tau (z_s)}{\tau _{SIS}(z_s)}`$ (thick light lines) and optical depths $`\frac{B_f}{B_f^{SIS}}\frac{\tau (z_s)}{\tau _{SIS}(z_s)}`$ (thick dark lines) to multiple imaging with respect to an SIS as a function of $`C_{NFW}`$ for $`\beta =1.0`$ (left hand panel) and $`\beta =0.5`$ (right hand panel). We have assumed $`\alpha _m=0.2`$, appropriate for the faint end of the quasar luminosity function. The light curves are sections through Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER. The corresponding curves for $`\frac{B_f}{B_f^{SIS}}\frac{\tau (z_s)}{\tau _{SIS}(z_s)}`$ show the decreased slope of $`\frac{B_f}{B_f^{SIS}}\frac{\tau (z_s)}{\tau _{SIS}(z_s)}`$ with $`C_{NFW}`$ for large $`C_{NFW}`$, and the increased slope near the no-multiple-imaging cutoff discussed at the end of Sec. 6. The upper panels also show hypothetical $`C_{NFW}`$ distributions having means of $`C_{NFW}=6.1`$ (thin solid lines) and $`C_{NFW}=10.9`$ (thin dashed lines). In each case one delta function and one Gaussian distribution (half-width of $`\sigma _{C_{NFW}}=2`$) are plotted. The lower panels show the dependence on $`\sigma _{C_{NFW}}`$ of convolutions of $`\frac{\tau (z_s)}{\tau _{SIS}(z_s)}`$ (light lines) and $`\frac{B_f}{B_f^{SIS}}\frac{\tau (z_s)}{\tau _{SIS}(z_s)}`$ (dark lines) with Gaussian distributions of $`C_{NFW}`$. The solid and dashed lines correspond to convolutions with $`C_{NFW}`$ distributions having $`C_{NFW}=6.1`$ and $`10.9`$ respectively. The convolutions in the lower left and right panels assumed $`\beta =1.0`$ and $`\beta =0.5`$.
The optical depth shows significant dependence on the range of $`C_{NFW}`$. For example, the simulations of Bullock et al. (2000) yield $`\beta =1`$, $`C_{NFW}6`$ and $`\sigma _{C_{NFW}}2`$ for CDM halos. The spread of $`\sigma _{C_{NFW}}2`$ results in a 2 fold increase in optical depth to multiple imaging. Much more dramatic examples are obtained where $`C_{NFW}`$ lies close to, or below the no-multiple-imaging boundary, and the significant portion of optical depth is produced by outliers in the concentration population. This is the region of parameter space important for SIDM halos. We use the S1W-b halo of Yoshida, Springel, White & Tormen (2000) as an example. Using the $`C_{NFW}\beta `$ degeneracy lines in Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER, we find that this halo may be approximated by a Zhao profile with $`\beta =0.5`$ and $`C_{NFW}6`$. These parameters place the S1W-b halo in the no-multiple imaging regime. However, if halos have concentrations distributed about $`C_{NFW}6.1`$, the optical depth is finite. Taking the example of $`C_{NFW}=6`$, we find that 2 orders of magnitude separate the optical depth due to distributions of halos with $`\sigma _{C_{NFW}}=1`$ and $`\sigma _{C_{NFW}}=2`$, and that 3 orders of magnitude separate the optical depth between distributions having $`\sigma _{C_{NFW}}=1`$ and $`\sigma _{C_{NFW}}=4`$. Note that the S1W-c halo of Yoshida et al. (2000) has a concentration that differs by up to $`\mathrm{\Delta }C_{NFW}8`$ from the present day value during its evolution. These results demonstrate that knowledge of the distribution of halo concentrations is an important consideration for cluster lensing studies.
## 9 THE CONTRIBUTION OF CLUSTER GALAXIES
The statistics of multiple imaging due to clusters of galaxies are sensitive to the inner most regions of the mass distribution, where the contribution of a central galaxy may be important. One expects that this is particularly true where the halo is less concentrated or or has a shallower cusp, as is the case for SIDM halos. The presence of an isothermal central galaxy renders the cross-section to multiple imaging finite for all halos.
Williams, Navarro & Bartelmann (1999) have investigated the effect of including a massive central galaxy. They find that arc properties of the most massive clusters are consistent with predictions of the NFW profile, but that observed separations are too large in the lower mass cluster lenses. They demonstrate that a centrally dominant galaxy with a dark halo can resolve the discrepancy. However there are several caveats. As noted by Williams et al. (1999), their cluster sample over represents the expected fraction of massive clusters (based on the algorithm of Press & Schechter 1974), possibly indicating a selection bias towards large splittings and highly elongated arcs. Furthermore, the higher resolution of simulations of Moore et al. (1998) show a steep core $`\rho (r)r^{\frac{3}{2}}`$, indicating that the concentrations found by NFW may be a factor of $`2`$ too low. This reduces the required mass of the centrally dominant galaxy. More importantly, Bullock et al. (2000) find significant scatter in concentration for a halo of given mass. High concentration halos are more likely to lens, produce larger image splittings, and so are more likely to be observed as cluster lenses. In the previous section this was shown to be a significant effect and might give the appearance that the profiles found from N-body simulations cannot account for observations of cluster lenses. On the other hand, Bullock et al. (2000) also find a concentration that decreases faster with redshift than originally thought (Navarro, Frenk & White 1997). This will exasperate the problem pointed out by Williams et al. (1999) for cluster lenses above a redshift of $`1`$.
On the other hand, a cluster need not have a single centrally dominant galaxy (e.g. CL0024+1654). Flores, Maller & Primack (2000) have analyzed the inclusion of cluster galaxies on the statistics of giant arc properties. They find that observationally based constraints suggest that there are not enough massive galaxies in a cluster to significantly alter arc statistics. In a complementary study, Menghetti et al. (2000) artificially added galaxies to N-body simulations. When the properties of short arcs are excluded, their analysis suggests that the contribution of cluster galaxies is negligible.
A proper account of the cluster density profile in the presence of a centrally dominant galaxy requires a full treatment of galaxy formation at the center of the cluster. The studies mentioned artificially add intrinsically or extrinsically truncated parametric galactic halos to the cluster halo. Williams, Navarro & Bartelmann (1999) (following Navarro, Frenk & White (1997)) take the additional step of modifying the cluster halo adiabatically in response to the formed central galaxy. Profiles with large and small $`\beta `$ describe $`CDM`$ and $`SIDM`$ halos respectively. The response of the dark matter to the formation of a central galaxy cannot be commonly approximated in these two cases.
In CDM halos, the central galaxy may form adiabatically (see Blumenthal, Faber, Flores & Primack 1984), resulting in additional contraction of the cluster dark matter halo. In this scenario, it is possible that the central galaxy may retain its own halo if it is not tidally striped during formation. Note however that N-body simulations of cluster halos are expected to properly describe the halo of the central galaxy, since the cusp in the cluster density profile is formed from the infall of smaller halos.
To reconcile observations with CDM NFW halo models, Williams et al. (1999) find that an additional central cluster galaxy of mass $`3\times 10^{12}M_{}`$is required. With a stellar velocity dispersion of $`\sigma _{}=300kmsec^1`$, this implies that the central galaxy must retain its halo during formation. To find an approximate contribution to the multiple imaging cross-section we have taken the most simple approach and superimposed the dark matter halo profile of Brainerd, Blandford & Smail (1996),
$$\rho (r)=\frac{\sigma _{}^2r_o^2}{4\pi Gr^2(r^2+r_o^2)},$$
(41)
onto the center of the Zhao profile. Here $`\sigma _{}`$ is the central velocity dispersion and $`r_o`$ is the characteristic outer scale on the center of the Zhao profile. The profile has a total mass
$$M_g=\frac{\pi r_o\sigma _{}^2}{2G}.$$
(42)
We have calculated cross-sections to multiple imaging for Zhao profiles having $`M_{200}10^{15}M_{}`$ in the presence of an additional central galaxy of mass $`M_g=3\times 10^{12}M_{}`$ with a central stellar velocity dispersion of $`\sigma _{}=300kmsec^1`$. The resulting cross-section calculations are applicable to Zhao profiles with large $`\beta `$ and $`C_{NFW}`$ appropriate for CDM halos. The left hand panels of Tabs. 9 and 9 show values of the cross-section to multiple imaging $`\frac{\tau _{d+b}(z_s)}{\tau _{SIS}(z_s)}`$, and the fractional increase in cross-section $`\frac{\tau _{d+b}(z_s)}{\tau (z_s)}`$ over a small grid of $`\beta `$ and $`C_{NFW}`$.
In the self interacting scenario, the response of the halo to the formation of the central galaxy should be larger. The contraction heats up the dark matter, after which interactions transfer the heat outwards, resulting in cooling and additional halo contraction. An isothermal approximation is therefore appropriate for the response of the SIDM halo to the formation of a central galaxy. Consider an initial dark matter density profile resulting from an isothermal distribution function:
$$\rho _d(r)\mathrm{exp}\frac{\mathrm{\Phi }_d(r)}{kT}$$
(43)
where $`k`$ is Boltzmann’s constant and $`T`$ the temperature. The ratio of the dark matter density before and after the addition of baryons is
$$\frac{\rho _{d+b}}{\rho _d}=\mathrm{exp}\frac{\mathrm{\Phi }_d\mathrm{\Phi }_{d+b}}{kT}\mathrm{exp}\frac{2\sigma _{}}{\sigma _{cstr}}.$$
(44)
Thus the density increase of the cluster halo with $`\sigma _{cstr}1000kmsec^1`$ and central galaxy velocity dispersion $`\sigma _{}300kmsec^1`$, the density increase is of order 20%. This density increase is negligible with respect to the orders of magnitude which separate the densities of different profiles in the inner most regions (see Fig. GRAVITATIONAL LENS STATISTICS FOR GENERALIZED NFW PROFILES: PARAMETER DEGENERACY AND IMPLICATIONS FOR SELF-INTERACTING COLD DARK MATTER).
If the cluster has a SIDM halo, the central galaxy cannot have retained its own halo during formation. The additional galactic mass should therefore consist only of stars, and is thus more confined. We have calculated cross-sections to multiple imaging for Zhao profiles having $`M_{200}10^{15}M_{}`$ in the presence of an additional central of mass $`M_g=3\times 10^{11}M_{}`$ (appropriate for the mass in stars of a typical central galaxy) with a central stellar velocity dispersion of $`\sigma _{}=300kmsec^1`$. The calculation is applicable to all the Zhao profiles assuming that the central dark-matter distribution is described solely by the Zhao profile. The right hand panels of Tabs. 9 and 9 show the resulting values of the cross-section to multiple imaging $`\frac{\tau _{d+b}(z_s)}{\tau _{SIS}(z_s)}`$, and the fractional increase in cross-section $`\frac{\tau _{d+b}(z_s)}{\tau (z_s)}`$ where a central galaxy is included over a small grid of $`\beta `$ and $`C_{NFW}`$. $`\frac{\tau _{d+b}(z_s)}{\tau (z_s)}`$ is of order 1 for all but the shallowest halo profiles.
The addition of the extra mass increases the cross-section of all profiles with respect to the SIS. We have not re-normalized the mass of the Zhao-profile-plus-galaxy as we are interested in the relative rather than the absolute cross-sections of different Zhao-profile-plus-galaxies. The addition of the central isothermal galaxy sets a lower limit on the cross-section, and reduces the dependence of multiple imaging on the cluster profile parameters, though the differences are still orders of magnitude. The values in Tab. 9 may be compared directly with the magnification bias (Tabs. 3 and 5). These show that the contribution of the central galaxy to the optical depth is not as important as that of magnification bias. Note that while magnification bias and the addition of a central galaxy both preferentially increase the cross-section of less concentrated clusters, the addition of a central galaxy reduces the magnification bias due to the increased cross-section. Thus, while a self-consistent calculation is yet to be done, we conclude that centrally dominant galaxies will not prevent lens statistics from being used as a probe of SIDM halos.
## 10 CONCLUSIONS
We have calculated the differential and total optical depths to multiple imaging, the average image splitting and the total magnification for a constant co-moving number density of unevolving generalized NFW (Zhao) profile cluster mass gravitational lenses. We find that the number of expected strongly lensed quasars is a very sensitive function of the profile parameters. Profiles whose central density is low either due to a shallow central cusp, or to a scale radius that is a reasonable fraction of the virial radius have cross-sections to multiple imaging that are reduced by a significant factor. Moreover, the separation of multiple images is reduced (by a factor of a few), although the total magnification is significantly enhanced. We find that the resulting magnification bias does not alter our conclusions.
The Zhao profile exhibits degeneracies between profile parameters with respect to lensing statistics. Similarly, profile fits have parameter degeneracies which are a function of the minimization quantity adopted. If the properties of a profile core are accurately reproduced by an approximately degenerate profile, then the parameter degeneracies are nearly equivalent to those obtained from strong lensing statistics. However, a profile that attempts to fit the entire halo will introduce serious uncertainty into the inferred lensing rates (up to an order of magnitude). This is particularly true in the region of parameter space where both $`C_{NFW}`$ and $`\beta `$ are small and will be an important consideration for detailed calculations of the lensing rate for SIDM halos based on parametric results from N-body simulations.
The lensing rate is a powerful probe of SIDM in clusters of galaxies, particularly in the absence of a centrally dominant galaxy. We have obtained profile parameters for halos around clusters of galaxies by fitting Zhao profiles to the simulations of Yoshida, Springel, White & Tormen (2000). We find that the optical depth to multiple imaging seriously constrains the SIDM self-interaction cross-section. In particular, an interaction cross-section of $`10.0cm^2gm^1`$ is rarely capable of producing multiple images during its evolution. In the most highly concentrated phases, a halo composed of interacting dark matter with this cross-section has a rate of multiple imaging 1 to 2 orders of magnitude lower than the corresponding typical CDM halo. Preliminary calculations show that centrally dominant galaxies increase the multiple imaging cross-section, but that the increase is small with respect to the variation in cross-section among different profiles. The presence of centrally dominant galaxies should not therefore inhibit the use of the multiple imaging rate as a probe of SIDM.
An important result from this study, with implications for cluster lensing studies is that the scatter in CDM/SIDM profile parameters obtained from N-body studies describe a very large range of optical depth. The lensing statistics will therefore be dominated by the more concentrated members of the population, rather than by the typical halo. As a result, an estimate of the distribution of halo profiles must be included in future studies of cluster lensing statistics.
The authors would like to thank Romeel Dav$`\stackrel{´}{\mathrm{e}}`$, Bartosz Pindor and Daniel Mortlock for helpful and stimulating discussions. We would also like to thank the anonymous referee whose comments led to the improvement of this work. This research was supported by NSF grant AST98-02802 to ELT. JSBW acknowledges the support of an Australian Postgraduate award and a Melbourne University Overseas Research Experience Award.
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# FLOW EQUATIONS IN THE LIGHT-FRONT QCD
## 1 Introduction
The perturbative aspects of non-abelian gauge theories were underestood many years ago, and the perturbative calculations provided convincing proof that QCD is the theory of strong interactions. However nonperturbative QCD phenomena have been difficult to analyze mainly because calculational techniques are still lacking, even though the qualitative features have been more or less understood.
In particular, it is widely believed that pure Yang-Mills theory, with no dynamical quarks, posseses the features like asymptotic freedom, mass generation through the transmutation of dimensions, and confinement: linear rising potential between static (probe) quarks. Adding dynamical quarks chiral symmetry is broken spontaneously. The ultimate aim of this study is to understand these nonperturbative mechanisms in a Hamiltonian framework, solving flow equations for canonical QCD Hamiltonian in the few lowest sectors selfconsistently for dynamical qluons and quarks, and their effective interactions. In this work dynamics of quarks has been excluded to disentangle the complexity of chiral symmetry breaking.
In the past few years there were several studies addressing the issue of confinement and generation of mass gap in the framework of the Schrödinger picture $`^\mathrm{?}`$, $`^\mathrm{?}`$. In the relatively recent works $`^\mathrm{?}`$ one have been using the special ansatz for a vacuum wave functional suggested by Kogan and Kovner, and integrating over all possible gauge configurations. To mention a few early works, see refs. $`^\mathrm{?}`$. The calculational technique in this approach is still rather complicated, and allows to solve a field theory problem only in $`1+1`$ dimensions $`^\mathrm{?}`$, but treats only ground states in $`3+1`$ dimensions $`^\mathrm{?}`$.
In alternative studies of nonperturbative problems in Hamiltonian framework one considers directly the QCD Hamiltonian quantized in a special gauge, in particular the light-front gauge, $`A^+=0`$ $`^\mathrm{?}`$. There is a belief that the light-front gauge may be the most suitable frame to study the nonperturbative QCD $`^\mathrm{?}`$. Previous studies of confinement and bound states in the light-front frame have been done using the methods of similarity renormalization $`^\mathrm{?}`$ and transverse lattice $`^\mathrm{?}`$, and are based on the fact that the light-front QCD in $`3+1`$ already has a confining interaction in the form of the instantaneous four-fermion interaction, $`1/q^{+2}`$, which is the confining interaction in $`1+1`$ along the light-front direction, $`x^{}`$ . However, since instantaneous interaction does not provide confinement mechanism for quarks and gluons in $`3+1`$ dimensions, the task of maintenance rotational symmetry becomes difficult to achieve when trying to extend light-front confinement to $`3+1`$.
Some time ago Wilson proposed a formalism to construct a confining light-front quark-gluon Hamiltonian for light-front $`QCD_{3+1}`$ $`^\mathrm{?}`$. Wilson suggested that a starting point for analyzing the full QCD with confinement in light-front coordinates is the light-front infrared divergences. Based on light-front power counting, the counterterms for the light-front infrared divergences can involve the color charge densities and involve unknown nonlocal behavior in transverse direction that become a possible source for transverse confinement. However, the analysis was not complete and a scheme for practical calculation has not been developed.
A subtle point in the light-front field theory is that the light-front vacuum is just empty space. Therefore it seems a problem how confinement can occur in the light-front frame, and what quantity sets up a scale for a dynamical mass gap and the string tension. The infrared longitudinal cutoff properties of light-front theory suggest a fundamental role for the light-front counterterms, solving this paradox. Namely, the longitudinal infrared cutoff in light-front dynamics makes it impossible to create particles from a bare vacuum by a translationally invariant Hamiltonian and in addition the number of constituents in a given eigenstate is limited. When introducing infrared cutoff physics below the cutoff is missed, and to restore it one may introduce appropriate counterterms. Light-front counterterms to the longitudinal infrared cutoff dependence provide a nonzero amplitude of particle creation, and become therefore a possible alternative source for features normally associted in standard equal-time dynamics with a nontrivial vacuum structure, including confinement and spontaneous symmetry breaking. Note, that small light-front $`x`$ correspond to high light-front energies. Therefore in order to remove small $`x`$, one should use renormalization group.
To be more specific we adopt the following model suggested by Susskind and Burkardt in the context of chiral symmetry breaking in the light-front frame $`^\mathrm{?}`$. In the parton model one pictures a fast moving hadron as being some collection of constituents with relatively large momentum, such that when one boosts the system, doubles its momentum, all these partons double their momenta and so forth. Therefore one can formulate an effective field theory on the axis of the light-front momentum $`x`$ (or on rapidity axis, which is logarithm of $`x`$). Partons that form a hadron are at positive, finite $`x`$ and according to Feynman and Bjorken fill the $`x`$-axis in a way which gets denser and denser as one goes to smaller $`x`$; and the vacuum is at $`x=0`$. The fundamental property of light-front Hamiltonians, that under a rescaling of the light-front momentum, $`x\lambda x`$, the light-front Hamiltonian scales like $`HH/\lambda `$, can be interpreted as a dilatation symmetry along the $`x`$-axis, if one thinks of the $`x`$-axis as a spacial axis. This symmetry holds on a classical level and it is broken on a quantum level by anomali. As one approaches small $`x`$, interaction between partons gets stronger, contributing divergent matrix elements. A natural cutoff is provided by $`\delta x=\epsilon ^+x`$, a minimal spacing between constituents, which plays the role of UV-regulator. As long as $`\delta x`$ (or $`\epsilon ^+`$) is finite, i.e. as long as the density of partons on the $`x`$-axis is not infinite, one obtains finite matrix elements. Cutoff $`\delta x`$ breaks the dilatation symmetry along the $`x`$-axis and gauge symmetry, and sets up an energy scale in effective light-front field theory formulated on $`x`$-axis. In terms of effective theory a generated mass gap defines a strength of effective interactions between quarks, in this case string tension in quark-antiquark potential. Formation of the $`q\overline{q}`$ bound state through breaking an internal symmetry can be viewed analogously to the creation of Cooper pairs in superconducter.
The dilatation symmetry reflects some underlying scale invariance of the light-front field theory formulated on $`x`$-axis. Introducing the cutoff, breaks this symmetry. Because physics should remain independent from the cutoff, one must be looking for a fixed point of the renormalization group. Therefore the right tool for studing such a system is the renormalization group, which is provided by the method of flow equations for Hamiltonians $`^\mathrm{?}`$.
Incorporating effects from small $`x`$ into an effective light-front Hamiltonian is equivalent to integrating out high light-front energy modes in asymtotically free domain. In terms of renormalization group, regulating small $`x`$ introduces a mass gap, which together with asymptotic freedom leads to a renormalization group invariant scale and dimensional transmutation along $`x`$-axis. Mass gap and a possible, more singular than a Coulomb, confining potential between quark and antiquark are direct consequences of dimensional transmutation in the effective light-front field theory, namely the light-front QCD, formulated on the light-front $`x`$-axis.
In the main part of the paper, basing on the QCD Hamiltonian in the light-front gauge, flow equations for an effective gluon mass and effective quark-antiquark interaction are formulated in the light-front frame and solved selfconsistently within the leading iteration. By discussing solutions of these equations in concluding section, it seems that flow equation method is superior to perturbation theory and perturbative similarity approach.
## 2 Flow equations in the light-front QCD
We apply flow equations to the light-front QCD Hamiltonian in order to eliminate the minimal quark-gluon interaction, namely to decouple matter and gauge degrees of freedom in the leading order. In the two lowest Fock sectors of the effective QCD Hamiltonian we obtain coupled differential equations which correspond to renormalization of the light-front gluon mass (renormalization of the quark mass is not considered) and generation of an effective quark-antiquark interaction. These flow equations are solved selfconsistently in the sense that the influence of the gluon mass renormalization on the elimination of the quark-gluon coupling and the induced quark-antiquark interaction is taken into account. As a result, the original gauge field coupling is completely eliminated, even when the states connected by this interaction are degenerate. Furthermore, in the degenerate case where effective $`q\overline{q}`$ interaction obtained within perturbation theory is not defined, we obtain more singular $`1/q^4`$ behavior at small gluon momenta.
### 2.1 Gluon gap equation
Coupled system of the light-front equations for the effective quark and gluon masses as functions of a cut-off $`\lambda `$ was derived first by Glazek $`^\mathrm{?}`$. We decouple this system of equations by neglecting the cut-off dependence of the quark mass, i.e. $`m(\lambda )=m`$ with $`m`$ current quark mass. The light-front gluon gap equation reads
$`{\displaystyle \frac{d\mu ^2(\lambda )}{d\lambda }}=`$ $``$ $`2T_fN_fg^2{\displaystyle _0^1}{\displaystyle \frac{dx}{x(1x)}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d^2k_{}}{16\pi ^3}}{\displaystyle \frac{1}{Q_2^2(\lambda )}}{\displaystyle \frac{df^2(Q_2^2(\lambda );\lambda )}{d\lambda }}`$ (1)
$`\times `$ $`\left({\displaystyle \frac{k_{}^2+m^2}{x(1x)}}2k_{}^2\right)`$
$``$ $`2C_ag^2{\displaystyle _0^1}{\displaystyle \frac{dx}{x(1x)}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d^2k_{}}{16\pi ^3}}{\displaystyle \frac{1}{Q_1^2(\lambda )}}{\displaystyle \frac{df^2(Q_1^2(\lambda );\lambda )}{d\lambda }}`$
$`\times `$ $`\left(k_{}^2(1+{\displaystyle \frac{1}{x^2}}+{\displaystyle \frac{1}{(1x)^2}})\right),`$
with
$`Q_1^2(\lambda )`$ $`=`$ $`{\displaystyle \frac{k_{}^2+\mu ^2(\lambda )}{x(1x)}}\mu ^2(\lambda ),Q_2^2(\lambda )={\displaystyle \frac{k_{}^2+m^2}{x(1x)}}\mu ^2(\lambda ),`$ (2)
where in the integral kernel gluon couples to the quark-antiquark pairs and pairs of gluons with the bare strong coupling $`g`$; $`(x,k_{})`$ is the light-front momentum in the loops. Group factors for $`SU(N_c)`$ are $`T_f\delta _{ab}=\mathrm{Tr}(T^aT^b)=\frac{1}{2}\delta _{ab}`$ and the adjoint Casimir $`C_a\delta _{ab}=f^{acd}f^{bcd}=N_c\delta _{ab}`$, $`N_c`$ is the number of colors (i.e., $`N_c=3`$).
Solution of this gap equation defines an effective gluon mass at zero gluon momentum and within the leading iteration reads (for details see $`^\mathrm{?}`$)
$`\mu ^2(\lambda )`$ $`=`$ $`\mu _0^2+\delta \mu _{PT}^2(\lambda )+\delta \mu ^2(\lambda ,\lambda _0).`$ (3)
where $`\mu _0`$ is the ’physical’ mass parameter, which fixes the resulting effective gluon mass at the scale $`\lambda _00`$ as $`\mu _{ren}^2(\lambda =\lambda _0)=\mu _0^2`$ (see below); the perturbative term
$`\delta \mu _{PT}^2(\lambda )={\displaystyle \frac{g^2}{4\pi ^2}}\lambda ^2\left(C_a(\mathrm{ln}{\displaystyle \frac{u^2}{\mu _0^2}}{\displaystyle \frac{11}{12}})+T_fN_f{\displaystyle \frac{1}{3}}\right).`$ (4)
reproduces the result of the light-front perturbation theory $`^\mathrm{?}`$. Renormalizing the effective Hamiltonian through the second oder in coupling $`O(g^2)`$, the perturbative mass correction is absorbed by the mass counterterm, $`m_{CT}^2(\mathrm{\Lambda }_{UV})=\delta \mu _{PT}^2(\mathrm{\Lambda }_{UV})`$ with $`\mathrm{\Lambda }_{UV}\mathrm{}`$, and the renormalized effective gluon mass reads $`\mu _{ren}^2(\lambda )=\mu ^2(\lambda )+m_{CT}^2(\lambda )`$ for $`\lambda =\mathrm{\Lambda }_{UV}\mathrm{}`$. Though the perturbative renormalization is applied at large cut-off scales, $`\mathrm{\Lambda }_{UV}`$, we assume that the leading cutoff dependence is absorbed by the mass counterterm for all $`\lambda `$. Therefore the resulting effective mass, renormalized in the second order, is given
$`\mu _{ren}^2(\lambda )`$ $`=`$ $`\mu _0^2+\delta \mu ^2(\lambda ,\lambda _0)=\mu _0^2+\sigma (\mu _0,u)\mathrm{ln}{\displaystyle \frac{\lambda ^2}{\lambda _0^2}}`$ (5)
$`\sigma (\mu _0,u)`$ $`=`$ $`{\displaystyle \frac{g^2}{4\pi ^2}}\mu _0^2\left(C_a({\displaystyle \frac{u^2}{\mu _0^2}}+\mathrm{ln}{\displaystyle \frac{u^2}{\mu _0^2}}{\displaystyle \frac{5}{12}})+T_fN_f({\displaystyle \frac{1}{3}}+{\displaystyle \frac{m^2}{\mu _0^2}})\right),`$
where scale $`u`$ has been introduced $`^\mathrm{?}`$, $`^\mathrm{?}`$ to regulate small light-front $`x`$ divergences, $`x0`$ and $`x1`$, which correspond to high light-front energies. In asymtotic free theories, such as QCD, the regulating scale can be related to the gauge invariant scale, using Callan-Semanzchik type equation. Namely scale $`u`$ can be expressed through $`\mathrm{\Lambda }_{QCD}`$, solving the third order flow equations for the strong coupling constant $`\alpha _s(\lambda )`$. Some calculations have been recently done in this direction for the asymptotic free toy model $`^\mathrm{?}`$ and for the three-gluon vertex in QCD $`^\mathrm{?}`$. However, here we do not perform these calculations and assume that the value $`u`$ is given by the hadron scale, $`u\mathrm{\Lambda }_{hadron}`$.
The resulting effective Hamiltonian is defined at the scale $`\lambda 0`$, therefore an effective gluon mass equals the ’physical’ mass, $`\mu _{ren}^2(\lambda =\lambda _0=0)=\mu _0^2`$. In particular, when the ’physical’ mass is set to zero, $`\mu _0=0`$, the effective QCD Hamiltonian has zero mass gauge fields, therefore our unitary transformaion does not violate gauge invariance (though at the intermediate stages for finite $`\lambda `$ gauge invariance is broken by the cutoffs). From Eq. (5) one has in the limit $`\mu _00`$
$`\sigma ^{}=\underset{\mu _00}{lim}\sigma (\mu _0,u)=u^2{\displaystyle \frac{g^2C_a}{2\pi ^2}},`$ (6)
and, as shown below, $`\sigma ^{}`$ plays the role of the string tension between quark and antiquark. In Eq. (5) the squared of the light-front cutoff, $`u^2`$, defines the rate how fast the effective gluon mass approaches asymtotically $`\lambda \lambda _00`$ (from above $`\lambda \lambda _0`$) the ’physical’ value $`\mu _0`$. Asymptotic behavior of the effective gluon mass near the renormalization point $`\lambda \lambda _0`$ Eq. (5) is important to take into account when solving flow equations for effective quark interactions at vanishing gluon momenta. In the next section this dependence, Eq. (5), is used to find an effective quark potential, generated by flow equations.
### 2.2 Effective quark-antiquark interaction
Eliminating the quark-gluon coupling generates an effective interaction in the quark-antiquark sector. In the light-front frame an effective quark-antiquark interaction is given
$`V_{q\overline{q}}=4\pi \alpha _sC_f\gamma ^\mu \gamma ^\nu \underset{(\mu _0,\lambda _0)0}{lim}B_{\mu \nu },`$ (7)
which includes dynamical interactions generated by flow equations and the instananeous term present in the original light-front gauge Hamiltonian. The gluon renormalization mass parameter (’physical’ gluon mass) $`\mu _0`$ and the renormalization point $`\lambda _0`$ are set to zero at the end of calculations. Here the current-current term in the exchange channel reads
$`\gamma ^\mu \gamma ^\nu ={\displaystyle \frac{\left(\overline{u}(k_{},(1x))\gamma ^\mu u(k_{},x)\right)\left(\overline{v}(k_{}^{},x^{})\gamma ^\nu v(k_{}^{},(1x^{}))\right)}{\sqrt{xx^{}(1x)(1x^{})}}},`$ (8)
where helicities of quarks are suppressed for simplicity; and the interaction kernel is given in full analogy with an effective electron-positron interaction in the light-front frame $`^\mathrm{?}`$, except for keeping the cutoff dependence in the four-momentum transfers, as
$`B_{\mu \nu }=g_{\mu \nu }\left(I_1+I_2\right)+\eta _\mu \eta _\nu {\displaystyle \frac{\delta Q^2}{q^{+2}}}\left(I_1I_2\right),`$ (9)
where $`g_{\mu \nu }`$ is the light-front metric tensor, and the light-front unity vector $`\eta _\mu `$ is defined as $`\eta k=k^+`$. The cutoff dependence of four-momentum transfers along quark and antiquark lines is accumulated in the factor
$`I_1={\displaystyle _0^{\mathrm{}}}𝑑\lambda {\displaystyle \frac{1}{Q_1^2(\lambda )}}{\displaystyle \frac{df(Q_1^2(\lambda );\lambda )}{d\lambda }}f(Q_2^2(\lambda );\lambda ),`$ (10)
$`I_2`$ is obtained by the interchange of indices $`1`$ and $`2`$; $`f(z)`$ is a similarity function; and the light-front four-momentum transfers are defined
$`Q_1^2(\lambda )`$ $`=`$ $`Q_1^2+\mu _{ren}^2(\lambda )`$
$`Q_2^2(\lambda )`$ $`=`$ $`Q_2^2+\mu _{ren}^2(\lambda ),`$ (11)
with
$`Q_1^2`$ $`=`$ $`{\displaystyle \frac{(x^{}k_{}xk_{}^{})^2+m^2(xx^{})^2}{xx^{}}}`$
$`Q_2^2`$ $`=`$ $`Q_1^2|_{x(1x);x^{}(1x^{})},`$ (12)
where $`\mu _{ren}`$ is given in Eq. (5); and the average momenta are $`Q^2=(Q_1^2+Q_2^2)/2`$ and $`\delta Q^2=(Q_1^2Q_2^2)/2`$. Calculating an effective kernel with explicit similarity functions, one has (for details see $`^\mathrm{?}`$)
$`\underset{(\mu _0,\lambda _0)0}{lim}B_{\mu \nu }=g_{\mu \nu }\left({\displaystyle \frac{1}{Q^2}}+{\displaystyle \frac{\sigma ^{}}{Q^4}}\right)+O\left(({\displaystyle \frac{\delta Q^2}{Q^2}})^n\right),`$ (13)
where $`n=2`$ for smooth and $`n=1`$ for sharp cut-off functions, and the first term does not depend on the cut-off function. The four-momenta can be represented in the ’mixed’ light-front $`(x,k_{})`$ and instant $`\stackrel{}{k}=(k_z,k_{})`$ frames as
$`Q^2`$ $`=`$ $`\stackrel{}{q}^2(2x1)(2x^{}1)(M_1M_2)^2/4`$
$`\delta Q^2`$ $`=`$ $`(xx^{})(M_1^2M_2^2)/2,`$ (14)
where $`\stackrel{}{q}=\stackrel{}{k}\stackrel{}{k}^{}=(q_z,q_{})`$ is the gluon three-momentum transfer, and $`M_1^2=4(\stackrel{}{k}^2+m^2)`$ and $`M_2^2=4(\stackrel{}{k}^2+m^2)`$ are the total energies of the initial and final states. Therefore in the limit of vanishing gluon momentum $`\stackrel{}{q}0`$, which defines mainly the bound state spectrum because then the effective $`q\overline{q}`$ interaction is singular, the four-momenta are $`Q^2\stackrel{}{q}^2`$ and $`\delta Q^20`$, and the effective interaction is given
$`V_{q\overline{q}}=\gamma ^\mu \gamma _\mu \left(C_f\alpha _s{\displaystyle \frac{4\pi }{\stackrel{}{q}^2}}+\sigma {\displaystyle \frac{8\pi }{\stackrel{}{q}^4}}\right),`$ (15)
where a new constant $`\sigma `$ is introduced instead of $`\sigma ^{}`$, given by Eq. (6), as $`\sigma =\sigma ^{}\alpha _sC_f/2`$. One recovers the standard Coulomb and linear rising confining potentials, $`C_f\alpha _s/r+\sigma r`$, in configuration space. It is remarkable, that, though calculations are done in the light-front frame, the result for the leading quark-antiquark effective interaction, Eq. (15), is rotationally invariant. Confining term in Eq. (15) with singular behavior like $`1/\stackrel{}{q}^4`$, arises from elimination of the quark-gluon coupling at small gluon momenta, that is governed by the asymtotic behavior of the effective gluon mass.
As was shown in $`^\mathrm{?}`$ with $`\sigma =0`$, solving QED effective interaction for positronium spectrum numerically, rotational symmetry holds with high accuracy for smooth cut-off functions and even for a sharp cutoff if the collinear singular part is subtracted. Based on our analyses here, we anticipate to see in numerical calculations, that also in the QCD effective interaction rotationally nonsymmetric part contributes negligible and meson spectrum manifests rotational invariance.
## 3 Conclusions and outlook
An effective QCD Hamiltonian in the light-front gauge has been obtained, solving flow equations for the two lowest Fock sectors selfconsistently. It has been shown that it is possible to eliminate the minimal quark-gluon interaction by using continuous unitary transformation. This elimination causes the renormalization of the coupling functions of the Hamiltonian described by the flow equations. In the two lowest Fock sectors this change of the couplings corresponds to the renormalization of the one-particle energies and to the generation of effective interactions between quarks, in particular quark-antiquark interaction. In oder to set up these differential equations the generated new interactions, with more than three intermediate states, have to be neglected. Truncating in number of particles participating in intermediate states is not equivalent to perturbation theory in coupling constant, but rather is close to Tamm-Dancoff approach.
Our approach has the following advantages: (i) The original gauge field coupling can be completely eliminated, even when the states connected by this interaction are degenerate. The continuous transformation is chosen in such a way that the transformed Hamiltonian does not contain any interations between one (anti)quark and the creation or annihilation of one gluon. These interactions, connecting the states with energy differences less than a cutoff scale $`|E_pE_q|\lambda `$, are still present in similarity approach due to nelecting the renormalization of single particle energies. They may cause mixing between the low and high Fock sectors, and it is problematic to include them perturbatively. Ignoring these low-energy interactions may break the gauge invariance, and in the light-front frame the rotational symmetry. (ii) Effective quark-antiquark interaction is rotationally symmetric at small gluon momenta $`q`$, while all collinear singular terms $`1/q^+`$ and $`1/q^{+\mathrm{\hspace{0.17em}2}}`$ cancel; and contains in addition to a perturbative term $`1/q^2`$, which can be obtained in the second order perturbation theory, also a more singular behavior of $`1/q^4`$ type. Our result for the induced $`q\overline{q}`$-interaction differes also from the result of similarity scheme, where the collinear singular part of the uncanceled instantaneous interaction $`1/q^{+\mathrm{\hspace{0.17em}2}}`$, produces a logarithmic potential, which is not rotationally symmetric. The origin of these differences lies in the fact that in our approach all couplings depend on $`l`$. In order to obtain properties (i) and (ii) the influence of the renormalization of the one-particle energies, in particular the gluon effective energy, on the elimination of the quark-gluon coupling has to be taken into account. By doing so the renormalization of the light-front gluon mass $`\mu (\lambda )`$ at zero gluon momentum in the asymptotic regime of small cutoff scales is described by an integro-differential equation. This equation can be solved, assuming the renormalization condition that the renormalized gluon mass is given at some small cutoff scale $`\lambda _0`$ by the ’physical’ value $`\mu _0`$, $`\mu _{ren}(\lambda =\lambda _0)=\mu _0`$. Renormalization is understood in perturbative sense, i.e. the renormalized effective gluon mass is obtained by absorbing the leading cutoff dependence into the second order mass counterterm. As a result the asymptotic behavior of the renormalized gluon mass $`\mu _{ren}(\lambda )`$ for small cutoffs $`\lambda `$, approaching the renormalization point from above $`\lambda \lambda _0`$, has been obtained. As a consequence of the properties of $`\mu _{ren}(\lambda )`$ the quark-gluon coupling is always eliminated even in the case of degeneracies, namely for vanishing gluon momenta. A similar argumentation was used by Kehrein, Mielke and Neu $`^\mathrm{?}`$ for the spin-boson model, where the authors argued that the coupling to the bosonic bath always is eliminated because of the renormalization of the tunneling frequency. Also, in a complete analogous to our problem of interacting electrons in BCS-theory, Lenz and Wegner $`^\mathrm{?}`$ showed that the elimination of electron-phonon coupling for all states is a direct consequence of the renormalization of phonon frequency.
Furthermore it is shown that due to the asymptotic behavior of the gluon energy, elimination of quark-gluon coupling at small gluon momenta gives rise to the enhancement of ’zero modes’ in the effective quark-antiquark interaction, i.e. in the infrared region $`q0`$ a more singular potential $`1/q^4`$, than in the perturbative case, is induced. By discussing the consequences of this asymptotic behavior it becomes clear that the approach of flow equations is superior to perturbative calculations, and probably also to (perturbative) similarity scheme which works in terms of bare unrenormalized fields. It can be seen that the flow of the coefficients (couplings) changes the generator $`\eta `$ of the unitary transformation. Even if the flow of couplings is obtained within the perturbative frame, the unitary transformation based on this generator $`exp(𝑑l\eta (l))`$, includes infinite many orders in perturbation theory, corresponding to (leading log) resummation of diagrams. It is worth mentioning, that Lenz and Wegner found in the system of interacting electrons $`^\mathrm{?}`$, that carrying out the $`l`$-ordering of $`\eta `$, the induced electron-electron interactions differ from the Fröhlich’s, where the unitary transformation based on the second order bound state perturbation theory is used. Moreover, this interaction is attractive in all momentum space, providing binding electrons in Cooper pairs. Kehrein and Mielke obtained similar modifications due to $`l`$-dependent generator by eliminating the hybridization term in the single impurity Anderson model by continuous unitary transformations $`^\mathrm{?}`$. The authors showed that their approach generates a spin-spin interaction which differs from the one obtained by the famous Schrieffer-Wolff transformation. Their induced interaction has the right high-energy cutoff, as compared to the Schrieffer-Wolff’s result. Summarizing, within flow equations approach it is possible to make statements which can not be obtained by perturbation theory.
Besides to this comparison with perturbative schemes, due to complete elimination of the quark-gluon coupling, flow equation for an effective quark-antiquark interaction can be integrated for all cutoffs down to $`\lambda =0`$. In similarity approach one removes couplings perturbatively untill the finite cutoff, below which perturbation theory breaks down. The choice of this cutoff depends on the problem considered, that might be ambiguous. In QCD treatment this cutoff introduces the nonzero scale in the theory, which breaks gauge and rotational invariance $`^\mathrm{?}`$. In our approach, the regulator of small light-front $`x`$ sets up a scale in the effective theory, in particular for the string tension in the effective quark-antiquark potential. Besides of the nonzero scale, the resulting renormalized gluon mass vanishes asymptotically, maitaining gauge invariance. As a consequence, the effective quark-antiquark interaction is rotationally symmetric. However, in this work the small light-front $`x`$ cutoff scale $`u`$ enters as an input parameter, and is fitted to the string tension from lattice calculations. To improve this, one should be looking for the fixed point of renormalization group and possible relate the cutoff $`u`$ with renormalization group invariant scale $`\mathrm{\Lambda }_{QCD}`$. In this way one should include higher Fock sectors in the internediate states. By doing so it is desirable to establish the independence on the regularization scheme, used to regulate small light-front momenta $`x`$.
The ultimate goal of the study is to solve the coupled chain of flow equations in different sectors selfconsistently. As has been shown, even an approximate solution of the gluon gap equation together with the flow equation for the effective interaction between probe quarks provides some information beyond the perturbation theory. The next step is to introduce dynamical quark degrees of freedom, formulating quark gap equation, and study the influence of the renormalization of the light-front quark mass on the effective interaction between quarks.
It seems, that it is possible to isolate in the light-front frame the degrees of freedom which are responsible for the long-range properties of QCD, and obtain some insight into the nonperturbative QCD phenomena. This suggests that probably the light-front formulation may be the most suitable frame to solve the system of flow equations selfconsistently on computer.
## Acknowledgments
The author would like to thank the organizers of the workshop for hospitality and support. The author is thanful to Stan Brodsky, Stan Glazek, Igor Klebanov, Lev Lipatov and Zvi Bern for helpful discussions. This work was supported by DOE grants DE-FG02-96ER40944, DE-FG02-97ER41048 and DE-FG02-96ER40947.
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# PARAMETER ESTIMATION IN ASTRONOMY WITH POISSON-DISTRIBUTED DATA. II. THE MODIFIED CHI-SQUARE-GAMMA STATISTIC
## 1 INTRODUCTION
The goodness-of-fit between an observation of $`N`$ data values, $`x_i`$, with errors, $`\sigma _i`$, and a model, $`m_i`$, can be determined by using the standard chi-square statistic:
$$\chi ^2\underset{i=1}{\overset{N}{}}\left[\frac{x_im_i}{\sigma _i}\right]^2.$$
(1)
The standard chi-square statistic is the appropriate statistic to determine the goodness-of-fit whenever the errors can be described by a normal (a.k.a. Gaussian) distribution.
Let us consider a more complicated situation where all the data values come from a pure counting experiment where each measurement<sup>1</sup><sup>1</sup>1 For example, X-ray photons, molecules, stars, galaxies, et cetera. , $`n_i`$, is a random integer deviate drawn from a Poisson (po1837 (1837), p. 205 et seq.) distribution,
$$P(k;\mu )\frac{\mu ^k}{k!}e^\mu ,$$
(2)
with a mean value of $`\mu `$. The use of equation (1) in analyzing Poisson-distributed data is technically never correct. While the Poisson distribution approaches the normal distribution as the Poisson mean approaches infinity, a Poisson distribution never actually becomes a normal distribution even at very large Poisson mean values. The normal distribution is always symmetric; the coefficient of skewness for the normal distribution is zero. Poisson distributions are always asymmetric; the coefficient of skewness for a Poisson distribution of mean $`\mu `$ is $`\mu ^{1/2}`$. While the Poisson distribution is almost symmetric about the mean for large mean values, its shape becomes progressively more asymmetric as the mean approaches zero. Thus the standard assumption that a Poisson distribution is approximately normally distributed is a good approximation only when the coefficient of skewness is negligible (e.g., $`\mu ^{1/2}1`$).
How does one then determine the goodness-of-fit with Poisson-distributed data? Historically, many $`\chi ^2`$ statistics have been proposed for the analysis of Poisson-distributed data. This paper will investigate the following four:
Pearson’s $`\chi ^2`$:
$$\chi _\mathrm{P}^2\underset{i=1}{\overset{N}{}}\frac{(n_im_i)^2}{m_i},$$
(3)
where the expectation value of the mean of the parent Poisson distribution of the $`i`$th data value is assumed to be equal to the Poisson deviate \[$`\mu _i=n_i`$\] and the square of the measurement error is assumed to be equal to the mean of the model Poisson distribution \[$`\sigma _i^2=m_i`$\];
the modified Neyman’s $`\chi ^2`$:
$$\chi _\mathrm{N}^2\underset{i=1}{\overset{N}{}}\frac{(n_im_i)^2}{\mathrm{max}(n_i,1)},$$
(4)
where the expectation value of the mean of the parent Poisson distribution of the $`i`$th data value is assumed to be equal to the Poisson deviate \[$`\mu _i=n_i`$\] and the square of the measurement error is assumed to be equal to Poisson deviate or one — whichever is greater \[$`\sigma _i^2=\mathrm{max}(n_i,1)`$\];
the Maximum Likelihood Ratio statistic for Poisson distributions:
$$\chi _\lambda ^22\underset{i=1}{\overset{N}{}}\left[m_in_i+n_i\mathrm{ln}\left(\frac{n_i}{m_i}\right)\right]$$
(5)
(see, e.g., Baker & Cousins baco1984 (1984) and references therein);
and the chi-square-gamma statistic (Mighell mi1999 (1999); hereafter PaperI):
$$\chi _\gamma ^2\underset{i=1}{\overset{N}{}}\frac{\left[n_i+\mathrm{min}(n_i,1)m_i\right]^2}{n_i+1},$$
(6)
where the expectation value of the mean of the parent Poisson distribution of the $`i`$th data value is assumed to be equal to the Poisson deviate plus a small correction factor of zero for zero deviates and one in all other cases \[$`\mu _i=n_i+\mathrm{min}(n_i,1)`$\] and the square of the measurement error is assumed to be equal to the Poisson deviate plus one \[$`\sigma _i^2=n_i+1`$\].
In PaperI, I demonstrated that the application of the standard weighted mean formula, $`\left[_in_i\sigma _i^2\right]/\left[_i\sigma _i^2\right]`$, to determine the weighted mean of data, $`n_i`$, drawn from a Poisson distribution, will, on average, underestimate the true mean by $``$$`1`$ for all Poisson mean values larger than $``$$`3`$ when the common assumption is made that the error of the $`i`$th observation is $`\sigma _i=\mathrm{max}(\sqrt{n_i},1)`$. This small, but statistically significant offset, explains the long-known observation that chi-square minimization techniques which use the modified Neyman’s $`\chi ^2`$ statistic \[eq. (4)\] to compare Poisson-distributed data with model values, $`m_i`$, will typically predict a total number of counts that underestimates the actual total by about $`1`$ count per bin (see, e.g., Bevington be1969 (1969), Wheaton et al. weet1995 (1995)).
Based on my finding that the weighted mean of data drawn from a Poisson distribution can be determined using the formula $`\left[_i\left[n_i+\mathrm{min}(n_i,1)\right]\left(n_i+1\right)^1\right]/\left[_i\left(n_i+1\right)^1\right]`$, I proposed that the chi-square-gamma statistic, $`\chi _\gamma ^2`$ \[eq. (6)\], should always be used to analyze Poisson-distributed data in preference to the modified Neyman’s $`\chi ^2`$ statistic. Following my own advice, I will not discuss the modified Neyman’s $`\chi ^2`$ statistic in the remainder of this article.
The chi-square distribution for $`\nu `$ degrees of freedom approaches a Gaussian distribution with a mean equal to $`\nu `$ (i.e., $`\mu \nu `$) and a variance equal to $`2\nu `$ (i.e., $`\sigma ^22\nu `$) as the number of degrees of freedom approaches infinity. Ideally, a $`\chi ^2`$ statistic for Poisson distributions for $`\nu `$ (independent) degrees of freedom would exhibit the same behavior as the number of degrees of freedom approaches infinity for all Poisson mean values (i.e., $`\mu >0`$).
Do the $`\chi _\mathrm{P}^2`$, $`\chi _\lambda ^2`$, and the $`\chi _\gamma ^2`$ statistics perform as expected for large Poisson mean values? These three $`\chi ^2`$ statistics are applied to the same data set in Fig. 1margin: $``$Fig1 (top to bottom, respectively). For this example, an ideal $`\chi ^2`$ statistic for Poisson-distributed data would have a cumulative distribution similar to that of the chi-square distribution for $`10^4`$ degrees of freedom which well approximated as the cumulative distribution function of a Gaussian distribution with a mean of $`10^4`$ and a variance of $`2\times 10^4`$. The results of the top and bottom panels are well matched to the expected cumulative distribution; the differences between the expected and measured mean and rms values are not statistically significant. The $`\chi _\mathrm{P}^2`$ and the $`\chi _\gamma ^2`$ statistics perform as expected with a Poisson mean value of 100. The cumulative distribution of the middle panel, however, clearly deviates from the expected cumulative distribution; the difference between the expected and measured mean and rms values, while small, is statistically significant. The $`\chi _\lambda ^2`$ does not perform like an ideal $`\chi ^2`$ statistic for Poisson distributions with a mean value of 100 — a level that is generally considered to be well above the low-count regime ($`\mu <25`$).
Let us continue the investigation of the performance of the Maximum Likelihood Ratio statistic for Poisson distributions with 1000 samples of $`10^4`$ Poisson deviates with Poisson mean values of 100, 10, 1, 0.1, and 0.001. Figure 2margin: $``$Fig2 confirms that the $`\chi _\lambda ^2`$ statistic does not perform like an ideal $`\chi ^2`$ statistic in the low-count regime. The average contribution by the $`i`$th deviate to an ideal $`\chi ^2`$ statistic for the analysis of Poisson-distributed data would be exactly one and the average contribution to its variance would be exactly two. Figure 3margin: $``$Fig3 expands the previous analysis of PaperI of the $`\chi _\lambda ^2`$ statistic over a wide range of Poisson mean values from 0.001 to 1000. The dashed lines of Fig. 3 show the results for an ideal $`\chi ^2`$ statistic; one can clearly see that the average contribution to $`\chi _\lambda ^2`$ is not equal to one and the average contribution to its variance is not equal to two for Poisson mean values $`<`$10. The poor performance of the $`\chi _\lambda ^2`$ statistic with low-count data may come as a surprise to many readers since it has historically been advocated as being one of the best $`\chi ^2`$ statistics for the analysis of Poisson-distributed data.
Chi-square statistics can serve (at least) two distinct purposes: (1) their functional forms can be utilized as the core of parameter estimation algorithms, and (2) their values can serve as a measure of the goodness-of-fit between a model and a data set.
$``$ While the functional form of the $`\chi _\lambda ^2`$ statistic can be successfully utilized for the purpose of parameter estimation with Poisson-distributed data in the low-count regime (see, e.g., PaperI), the Maximum Likelihood Ratio statistic for Poisson distributions \[eq. (5)\] should not be used to determine the goodness-of-fit with low-count data where the Poisson mean is $`<`$10.
In this work, I investigate the use of Pearson’s $`\chi ^2`$ statistic and the chi-square-gamma statistic for the determination of the goodness-of-fit between theoretical models and data derived from counting experiments. I develop a methodology in §2 which modifies Pearsons’s chi-square statistic for the purpose of improving its goodness-of-fit performance. This methodology is then be applied to modify the chi-square-gamma statistic (§3). The modified chi-square-gamma statistic is shown to perform (nearly) like an ideal $`\chi ^2`$ statistic for the determination of goodness-of-fit with low-count data. Simulated X-ray images are analyzed in §4 as a practical demonstration of the possible use of the modified chi-square-gamma statistic in experimental astrophysics. The summary of the paper is presented in §5.
## 2 THE MODIFIED PEARSON’S $`\chi ^2`$ STATISTIC
Let us continue the investigation of the performance of Pearson’s $`\chi ^2`$ statistic with 1000 samples of $`10^4`$ Poisson deviates with Poisson mean values of 100, 10, 1, 0.1, and 0.001. Figure 4margin: $``$Fig4 shows that the $`\chi _\mathrm{P}^2`$ statistic does not perform like an ideal $`\chi ^2`$ statistic for Poisson mean values $`<10`$. The average contribution by the $`i`$th deviate to an ideal $`\chi ^2`$ statistic for the analysis of Poisson-distributed data would be exactly one and the average contribution to its variance would be exactly two. Figure 5margin: $``$Fig5 expands the previous analysis of PaperI of Pearson’s $`\chi ^2`$ statistic over a wide range of Poisson mean values from 0.001 to 1000. The dashed lines of Fig. 5 show the results for an ideal $`\chi ^2`$ statistic; one can see that while the average contribution to $`\chi _\mathrm{P}^2`$ is one, the average contribution to its variance is not equal to two for Poisson mean values $`<10`$.
$``$ Pearson’s $`\chi ^2`$ statistic \[eq. (3)\] should not be used to determine the goodness-of-fit with low-count data where the mean of the parent Poisson distribution is $`<10`$.
The variance of Pearson’s $`\chi ^2`$ statistic is, by definition,
$$\sigma _{\chi _\mathrm{P}^2}^2\underset{i=1}{\overset{N}{}}\left[\frac{(n_im_i)^2}{m_i}\left(\frac{1}{\nu }\underset{j=1}{\overset{N}{}}\frac{(n_jm_j)^2}{m_j}\right)\right]^2,$$
(7)
where $`\nu =NM`$ is the number of independent degrees of freedom, N is the number of data values, and M is the number of free parameters. The variance of the reduced chi-square of a $`\chi ^2`$ statistic for a large number of observations should ideally be two. In the limit of a large number of observations of a single Poisson distribution with a mean value of $`\mu `$, the variance of the reduced chi-square of the Pearson’s $`\chi ^2`$ statistic is
$`\sigma _{\chi _\mathrm{P}^2/\mathrm{}}^2`$ $``$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{\sigma _{\chi _\mathrm{P}^2}^2}{\nu }}\right]`$ (8)
$``$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{\nu }}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{(n_im_i)^2}{m_i}}\left[{\displaystyle \frac{1}{\nu }}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{(n_jm_j)^2}{m_j}}\right]\right\}^2\right]`$
$``$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{\nu }}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{(n_im_i)^2}{m_i}}\left[{\displaystyle \frac{\chi _\mathrm{P}^2}{\nu }}\right]\right\}^2\right]`$
$``$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{\nu }}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{(n_im_i)^2}{m_i}}\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{\chi _\mathrm{P}^2}{\nu }}\right]\right\}^2\right]`$
$`=`$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{NM}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{(n_im_i)^2}{m_i}}1\right\}^2\right]\text{[see eq. (25) of }\text{PaperI}\text{]}`$
$`=`$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{N1}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{(n_i\mu _\mathrm{P})^2}{\mu _\mathrm{P}}}1\right\}^2\right]\text{[see eq. (5) of }\text{PaperI}\text{]}`$
$``$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{N1}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{\left(n_i\underset{N\mathrm{}}{lim}[\mu _\mathrm{P}]\right)^2}{\underset{N\mathrm{}}{lim}[\mu _\mathrm{P}]}}1\right\}^2\right]`$
$`=`$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{N1}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{(n_i\mu )^2}{\mu }}1\right\}^2\right]\text{[see eq. (7) of }\text{PaperI}\text{]}`$
$``$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{N1}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\left\{NP(k;\mu )\right\}\left\{{\displaystyle \frac{(k\mu )^2}{\mu }}1\right\}^2\right]`$
$`=`$ $`{\displaystyle \frac{1}{\mu ^2}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}P(k;\mu )\left[(k\mu )^2\mu \right]^2`$
$`=`$ $`2+{\displaystyle \frac{1}{\mu }}.`$
If we assume that Pearson’s $`\chi ^2`$ applied to a large number of observations of a single Poisson distribution with a mean value of $`\mu `$ always produces a normal distribution with a mean equal to the number of degrees-of-freedom ($`\nu `$) \[see eq. (25) of PaperI\] and a variance of $`\nu (2+\mu ^1)`$ \[see eq. (8)\], we can then attempt to create an ideal $`\chi ^2`$ statistic for the analysis of Poisson-distributed data by modifying Pearson’s $`\chi ^2`$ as follows:
$$\chi _{\mathrm{PM}}^2\underset{i=1}{\overset{N}{}}\left[\left\{\chi _{\mathrm{P}i}^2\chi _{\mathrm{P}i}^2\right\}\left[\frac{2}{\sigma _{\chi _{\mathrm{P}i}^2}^2}\right]^{1/2}+1\right]$$
(9)
where
$$\chi _{\mathrm{P}i}^2\frac{(n_im_i)^2}{m_i}$$
(10)
is the contribution of the $`i`$th data value to Pearson’s $`\chi ^2`$,
$$\chi _{\mathrm{P}i}^21$$
(11)
is the expectation value of $`\chi _{\mathrm{P}i}^2`$ \[see eq. (25) of PaperI\],
$$\sigma _{\chi _{\mathrm{P}i}^2}^22+m_i^1$$
(12)
is the variance of $`\chi _{\mathrm{P}i}^2`$ \[see eq. (8)\]. Translating the mathematical notation to English, we have (1) shifted the mean of the standard $`\chi _\mathrm{P}^2`$ distribution from $`\nu `$ times equation (11) to zero, (2) forced the variance of the shifted distribution to be exactly $`2\nu `$, and then (3) shifted the mean of the variance-corrected distribution from zero back to $`\nu `$. Thus, by definition, the modified Pearson’s chi-square statistic ($`\chi _{\mathrm{PM}}^2`$) will have a mean value of $`\nu `$ and a variance of $`2\nu `$ — in the limit of a large number of observations.
Let us now investigate the performance of the modified Pearson’s $`\chi ^2`$ statistic with 1000 samples of $`10^4`$ Poisson deviates with Poisson mean values of 100, 10, 1, 0.1, and 0.001. Figure 6margin: $``$Fig6 shows that $`\chi _{\mathrm{PM}}^2`$ results are significantly better than $`\chi _\mathrm{P}^2`$ results \[Fig. 4\] — especially for Poisson mean values less than 10. Figure 7margin: $``$Fig7 investigates the performance of the modified Pearson’s $`\chi ^2`$ statistic over a wide range of Poisson mean values from 0.001 to 1000. The dashed lines of Fig. 7 show the results for an ideal $`\chi ^2`$ statistic; one can see that while the average contribution to $`\chi _{\mathrm{PM}}^2`$ is 1, as expected, and the average contribution to its variance is equal to 2, as expected, the performance is not uniform for all Poisson mean values — the spread seen in the variance plot (bottom panel) increases as the Poisson mean approaches zero.
Figures 6 and 7 indicate the the modified Pearson’s $`\chi ^2`$ statistic works well in the perfect case where one has a priori knowledge of the true Poisson mean. In a real experiment, the true mean of the parent Poisson distribution is rarely (if ever) known and model parameters must be estimated from the observations. How well does the modified Pearson’s $`\chi ^2`$ statistic work with reasonable parameter estimates? Comparing Fig. 8margin: $``$Fig8 with Fig. 6 and Fig. 9margin: $``$Fig9 with Fig. 7, we see that the variances are significantly smaller when a realistic model (i.e., the sample mean) is used instead of a perfect model (i.e., the true mean). A statistic that fails with reasonable parameter estimates is not a very useful statistic for the analysis of real observations with low-count data.
$``$ The modified Pearson’s $`\chi ^2`$ statistic \[eq. (9)\] should not be used to determine the goodness-of-fit with low-count data where the mean of the parent Poisson distribution is $`<10`$.
## 3 THE MODIFIED CHI-SQUARE-GAMMA STATISTIC
Let us continue the investigation of the performance of the chi-square-gamma statistic with 1000 samples of $`10^4`$ Poisson deviates with Poisson mean values of 100, 10, 1, 0.1, and 0.001. Figure 10margin: $``$Fig10 confirms that the $`\chi _\gamma ^2`$ statistic does not perform like an ideal $`\chi ^2`$ statistic in the low-count regime. Figure 11margin: $``$Fig11 expands the previous analysis of PaperI of the $`\chi _\gamma ^2`$ statistic over a wide range of Poisson mean values from 0.001 to 1000. The $`\chi _\gamma ^2`$ statistic clearly does not perform like an ideal $`\chi ^2`$ statistic for Poisson mean values $`<10`$.
$``$ The chi-square-gamma statistic \[eq. (6)\] should not be used to determine the goodness-of-fit with low-count data where the mean of the parent Poisson distribution is $`<10`$.
The variance of the chi-square-gamma statistic is, by definition,
$$\sigma _{\chi _\gamma ^2}^2\underset{i=1}{\overset{N}{}}\left[\frac{[n_i+\mathrm{min}(n_i,1)m_i]^2}{n_i+1}\left(\frac{1}{\nu }\underset{j=1}{\overset{N}{}}\frac{[n_j+\mathrm{min}(n_j,1)m_j]^2}{n_j+1}\right)\right]^2,$$
(13)
where $`\nu =NM`$ is the number of independent degrees of freedom, N is the number of data values, and M is the number of free parameters. The variance of the reduced chi-square of a $`\chi ^2`$ statistic for a large number of observations should ideally be two. In the limit of a large number of observations of a single Poisson distribution with a mean value of $`\mu `$, the variance of the reduced chi-square of the chi-square-gamma statistic is
$`\sigma _{\chi _\gamma ^2/\mathrm{}}^2`$ $``$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{\sigma _{\chi _\gamma ^2}^2}{\nu }}\right]`$ (14)
$``$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{\nu }}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{\left[n_i+\mathrm{min}(n_i,1)m_i\right]^2}{n_i+1}}\left[{\displaystyle \frac{1}{\nu }}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\left[n_j+\mathrm{min}(n_j,1)m_j\right]^2}{n_j+1}}\right]\right\}^2\right]`$
$``$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{\nu }}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{\left[n_i+\mathrm{min}(n_i,1)m_i\right]^2}{n_i+1}}\left[{\displaystyle \frac{\chi _\gamma ^2}{\nu }}\right]\right\}^2\right]`$
$``$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{\nu }}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{\left[n_i+\mathrm{min}(n_i,1)m_i\right]^2}{n_i+1}}\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{\chi _\gamma ^2}{\nu }}\right]\right\}^2\right]`$
$`=`$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{NM}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{\left[n_i+\mathrm{min}(n_i,1)m_i\right]^2}{n_i+1}}\left[1+e^\mu \left(\mu 1\right)\right]\right\}^2\right]\text{[see eq. (29) of }\text{PaperI}\text{]}`$
$`=`$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{N1}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{\left[n_i+\mathrm{min}(n_i,1)\mu _\gamma \right]^2}{n_i+1}}\left[1+e^\mu \left(\mu 1\right)\right]\right\}^2\right]\text{[see eq. (18) of }\text{PaperI}\text{]}`$
$``$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{N1}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{\left[n_i+\mathrm{min}(n_i,1)\underset{N\mathrm{}}{lim}[\mu _\gamma ]\right]^2}{n_i+1}}\left[1+e^\mu \left(\mu 1\right)\right]\right\}^2\right]`$
$`=`$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{N1}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{{\displaystyle \frac{\left[n_i+\mathrm{min}(n_i,1)\mu \right]^2}{n_i+1}}\left[1+e^\mu \left(\mu 1\right)\right]\right\}^2\right]\text{[see eq. (19) of }\text{PaperI}\text{]}`$
$``$ $`\underset{N\mathrm{}}{lim}\left[{\displaystyle \frac{1}{N1}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\left\{NP(k;\mu )\right\}\left\{{\displaystyle \frac{\left[k+\mathrm{min}(k,1)\mu \right]^2}{k+1}}\left[1+e^\mu \left(\mu 1\right)\right]\right\}^2\right]`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}P(k;\mu )\left\{{\displaystyle \frac{\left[k+\mathrm{min}(k,1)\mu \right]^2}{k+1}}\left[1+e^\mu \left(\mu 1\right)\right]\right\}^2`$
$`=`$ $`\mu ^3e^\mu \left[\mathrm{Ei}(\mu )\gamma _{\mathrm{EM}}\mathrm{ln}(\mu )+4\right]\mu ^2\mu +e^\mu \left[2\mu ^2+2\mu +1\right]+e^{2\mu }[\mu ^2+2\mu 1],`$
where Ei$`(x)`$ is the exponential integral of $`x`$ $`[\text{Ei}(x)=_x^{\mathrm{}}\frac{e^t}{t}𝑑t=_{\mathrm{}}^x\frac{e^t}{t}𝑑t\text{ for }x>0]`$ and $`\gamma _{\mathrm{EM}}lim_n\mathrm{}\left[\left\{_{i=1}^n\frac{1}{n}\right\}\mathrm{ln}(n)\right]0.5772156649`$ is the Euler-Mascheroni constant. Equation (14) approaches the expected value of 2 for large Poisson mean values \[see the solid curve in the bottom panel of Fig. 11\].
If we assume that chi-square-gamma statistic is applied to a large number of observations of a single Poisson distribution with a mean value of $`\mu `$ always produces a normal distribution with a mean equal to the number of degrees of freedom ($`\nu `$) times see equation (29) of PaperI and a variance of $`\nu `$ times equation (14), we can then attempt to create an ideal $`\chi ^2`$ statistic for the analysis of Poisson-distributed data by modifying the chi-square-gamma statistic as follows:
$$\chi _{\gamma \mathrm{M}}^2\underset{i=1}{\overset{N}{}}\left[\left\{\chi _{\gamma i}^2\chi _{\gamma i}^2\right\}\left[\frac{2}{\sigma _{\chi _{\gamma i}^2}^2}\right]^{1/2}+1\right]$$
(15)
where
$$\chi _{\gamma i}^2\frac{\left[n_i+\mathrm{min}(n_i,1)m_i\right]^2}{n_i+1}$$
(16)
is the contribution of the $`i`$th data value to the chi-square gamma statistic,
$$\chi _{\gamma i}^21+e^{m_i}\left(m_i1\right)$$
(17)
is the expectation value of $`\chi _{\gamma i}^2`$ \[see equation (29) of PaperI\], and
$`\sigma _{\chi _{\gamma i}^2}^2`$ $``$ $`m_i^3e^{m_i}\left[\mathrm{Ei}(m_i)\gamma _{\mathrm{EM}}\mathrm{ln}(m_i)+4\right]m_i^2m_i+`$ (18)
$`e^{m_i}\left[2m_i^2+2m_i+1\right]+e^{2m_i}\left[m_i^2+2m_i1\right]`$
is the variance of $`\chi _{\gamma i}^2`$ \[see equation (14)\]. Translating the mathematical notation to English, we have (1) shifted the mean of the standard $`\chi _\mathrm{P}^2`$ distribution from $`\nu `$ times equation (17) to zero, (2) forced the variance of the shifted distribution to be exactly $`2\nu `$, and then (3) shifted the mean of the variance-corrected distribution from zero back to $`\nu `$. Thus, by definition, the modified chi-square statistic statistic ($`\chi _{\gamma \mathrm{M}}^2`$) will have a mean value of $`\nu `$ and a variance of $`2\nu `$ — in the limit of a large number of observations.
Let us now investigate the performance of the modified chi-square-gamma statistic with 1000 samples of $`10^4`$ Poisson deviates with Poisson mean values of 100, 10, 1, 0.1, and 0.001. Figure 12margin: $``$Fig12 shows that $`\chi _{\gamma \mathrm{M}}^2`$ results are significantly better than the $`\chi _\gamma ^2`$ results \[Fig. 10\] — especially for Poisson mean values less than 10. Figure 13margin: $``$Fig13 investigates the performance of the $`\chi _{\gamma \mathrm{M}}^2`$ statistic over a wide range of Poisson mean values from 0.001 to 1000. The dashed lines of Fig. 13 show the results for an ideal $`\chi ^2`$ statistic; one can see that while the average contribution to $`\chi _{\gamma \mathrm{M}}^2`$ is 1, as expected, and the average contribution to its variance is equal to 2, as expected, the performance is not uniform for all Poisson mean values. The bump seen in the bottom panel of Fig. 13 near the Poisson mean value of 10 is an artifact caused by the $`\mathrm{min}(\mathrm{n}_\mathrm{i},1)`$ offset in the numerator of the definition of the chi-square-gamma statistic \[eq. (6)\].
Figures 12 and 13 indicate the the modified chi-square-gamma statistic works well in the perfect case where one has a priori knowledge of the true Poisson mean. How well does the modified chi-square-gamma statistic work with reasonable parameter estimates? Comparing Fig. 14margin: $``$Fig14 with Fig. 12 and Fig. 15margin: $``$Fig15 with Fig. 13, we see that the results for the modified $`\chi _\gamma ^2`$ statistic with a realistic model (i.e., the sample mean) are nearly identical<sup>2</sup><sup>2</sup>2 The measured mean values appearing on the right side of top 3 panels of Fig. 14 are about 1 lower than the comparable value given in Fig. 12. This is the result of losing one degree-of-freedom due to the determination of the sample mean from the data (i.e., $`\nu `$ drops from 10000 to 9999). The scrambling caused by the modification of the $`\chi _\gamma ^2`$ statistic appears to have caused this expected loss of one degree-of-freedom to vanish in the very-low-count data regime ($`\mu <0.1`$). to those obtained with a perfect model (i.e., the true mean).
$``$ The modified chi-square-gamma \[eq. (15)\] statistic performs (nearly) like an ideal $`\chi ^2`$ statistic for the determination of the goodness-of-fit with low-count data. On average, for a large number of observations, the mean value of $`\chi _{\gamma \mathrm{M}}^2`$ statistic is equal to the number of degrees of freedom $`(\nu )`$ and its variance is $`2\nu `$ — like the $`\chi ^2`$ distribution for $`\nu `$ degrees of freedom.
## 4 SIMULATED X-RAY IMAGES
I now demonstrate the new modified chi-square-gamma statistic by using it to study simulated X-ray images. Cash (ca1979 (1979)) applied his $`C`$ statistic to the problem of determining the position of a weak source in a X-ray image. Let us use Cash’s Point Spread Function (PSF) but with a resolution of 100 pixels per unit area:
$$\varphi (x,y)\{\begin{array}{cc}\frac{\pi }{3}\left(1r\right)\hfill & \text{for }r1\text{,}\hfill \\ 0\hfill & \text{for }r>1\text{,}\hfill \end{array}$$
(19)
where $`r^2=(x/10)^2+(y/10)^2`$. This PSF has the volume integral of
$$\mathrm{\Phi }(x,y)\{\begin{array}{cc}6\left[\frac{r^2}{2}\frac{r^3}{3}\right]\hfill & \text{for }r1\text{,}\hfill \\ 1\hfill & \text{for }r>1\text{.}\hfill \end{array}$$
(20)
Figure 16margin: $``$Fig16 shows a simulated observation of a point source with an intensity of 40 X-ray photons on a background flux of 0.06 X-ray photons per pixel. This observation contains 2786 pixels with 0 photons, 204 pixels with 1 photons, and 10 pixels with 2 photons. There are a total of 56 photons found in the 317 pixels within a radius of 10 pixels of the center of the X-ray point source which is located at the $`(x,y)`$ position of $`(33,26)`$. This is clearly a marginal detection of a weak X-ray point source on a noisy background; the peak signal-to-noise ratio ($``$5.2) occurs at a radius of $``$8 pixels.
We will now use the modified chi-square-gamma statistic to answer the following questions about this X-ray image:
1. Is there an X-ray point source in the image?
2. If so, where is it located?
3. What is its total intensity?
The exact determination of the location and intensity of the X-ray point source in Fig. 16 is precluded by the fact that this particular observation contains only low-count data — we must be content with realistic estimates for the location and intensity based on a detailed statistical analysis of the data.
Our first objective is to determine if there is an X-ray point source in the observation. One way this can be done is to investigate the region(s) containing non-point-source pixels (data values). This approach requires knowledge of the background flux level — which we will henceforth assume is constant throughout the entire image. We begin by making a rough first estimate of the background flux level by dividing the total number of photon in the image by the total number of pixels: $`0.0747`$ $`(224/\mathrm{\hspace{0.17em}3000})`$ photons per pixel.
The background flux level estimate may be significantly improved with a bit more work. There are 18 photons in the 317 pixels within 10 pixels of the position $`(12,15)`$ of Fig. 16. Is the detection of 18 photons consistent with the expected value of 23.6799 $`(=0.0747\times 317)`$ photons? The upper and lower 99.9% single-sided confidence limits for 18 photons are 35.35 and 7.662, respectively \[see Tables 1 and 2 of Gehrels ge1986 (1986)\]. I conclude that the pixel at $`(12,15)`$ is a background pixel (shown as such with a gray box in Fig. 17margin: $``$Fig17 ) because the expected number of background photons lies within the range of the upper and lower 99.9% single-sided confidence limits of the observed number of photons (i.e., $`7.66223.679935.35`$). There are 56 photons in the 317 pixels within 10 pixels of the position $`(33,26)`$ of Fig. 16. The upper and lower 99.9% single-sided confidence limits for 56 photons are 83.1784 and 35.6834, respectively \[see eqs. (10) and (14) of Gehrels ge1986 (1986)\]. I conclude that the pixel at $`(33,26)`$ is not a background pixel because the expected number of background photons is less than the lower 99.9% single-sided confidence limit of the observed number of photons (i.e., $`23.6799<35.6834`$). This conclusion was expected since $`(33,26)`$ is the center of the X-ray source. The 2676 gray pixels in Fig. 17 have a total of 162 photons. We can now make a second estimate of the background flux: $`0.0605`$ $`(162/\mathrm{\hspace{0.17em}2676})`$ photons per pixel. Repeating this process once more yields the final estimate of the background flux: $`0.0601`$ $`(153/\mathrm{\hspace{0.17em}2547})`$ photons per pixel. The measurement error for this estimate is approximately 0.0049 $`(\sqrt{153+1}/\mathrm{\hspace{0.17em}2547})`$. The final estimate of the X-ray background flux, $`0.0601\pm 0.0049`$, is in excellent agreement with the true value of 0.06 \[see Fig. 18margin: $``$Fig18 \].
I conclude that Fig. 16 has at least one X-ray point source because the entire data set is not consistent with a X-ray background flux of 0.0601 photons per pixel for every pixel. Assuming that there is only one X-ray source, we can make the first rough estimate of its location by stating that it probably is located at a non-gray pixel location in Fig. 18. There are 453 non-background pixels in Fig. 18 with a total of 71 photons. This fact allows us to restrict the uncertainty of the location of the X-ray point source to about 15% $`(453/\mathrm{\hspace{0.17em}3000})`$ of the total image. Assuming a background flux of 0.06 photons per pixel, we expect that 27 $`(453\times 0.06)`$ photons of the total 71 photons found in the non-background pixels would be due to the background and not the point source. We can now make the first rough estimate of the intensity of the X-ray source: 44 $`(=7127)`$ photons.
There are 48 photons in the 317 pixels within 10 pixels of the position $`(24,26)`$ of Fig. 16. Is this photon sum consistent with a model of a 40 photon point source centered at that location on a background of 0.06 photons per pixel? The upper and lower 95% single-sided confidence limits for 48 photons is 61.05 and 37.20 photons, respectively. The model predicts that we should find 58.9802<sup>3</sup><sup>3</sup>3 These computations only included pixels within an aperture if the center of the pixel was within the given aperture radius; partial pixels whose center was just outside aperture boundary were rejected. The minimum number we would expect the model to predict is 58.8496 $`[40+(\pi \times 10^2\times 0.0600)]`$ photons. The maximum number we would expect the model to predict is 59.0200 $`[40+(317\times 0.0600)]`$ photons. The model prediction lies within these extremes: $`58.8496<58.9802<59.0200`$. photons within a radius of 10 pixels. The 48 photons found within 10 pixels of $`(24,26)`$ is consistent with the model (shown as such with a dark-gray circle in Fig. 18), because the expected number of photons lies within the range of the upper and lower 95% single-sided confidence limits of the observed number of photons (i.e., $`37.2058.980261.05`$). There are 217 circled pixels in Fig. 18. This fact allows us to further restrict the uncertainty of the location of the X-ray point source to $``$7.2% $`(217/\mathrm{\hspace{0.17em}3000})`$ of the total image.
Since the peak signal-to-noise ratio occurs near a radius of 8 pixels, we now investigate if we can improve our estimate of the location of the X-ray point source by considering photon sums within a smaller aperture with a radius of 8 instead of 10 pixels. There are 34 photons in the 197 pixels within a radius of 8 pixels of the position $`(26,26)`$ of Fig. 16. Is this photon sum consistent with a model of a 40 photon point source centered at that location with a background of 0.06 photons per pixel? The upper and lower 95% single-sided confidence limits for 34 photons is 45.27 and 25.01 photons, respectively. The model predicts that we should find 47.2664 photons within a radius of 8 pixels. The 34 photons found within 8 pixels of $`(26,26)`$ is not consistent with the model since the expected number of photons is greater than the upper 95% single-sided confidence limit of the observed number of photons (i.e., $`47.2664>45.27`$). However, the 39 photons found within 8 pixels of $`(27,26)`$ is consistent with a 40 photon point source centered at $`(27,26)`$ on a background of 0.06 photons per pixel (i.e., $`29.3347.266450.94`$). There are a total of 111 circled pixels in Fig. 19 margin: $``$Fig19 . This fact allows us to further restrict the uncertainty of the location of the X-ray point source to $``$3.7% $`(=111/\mathrm{\hspace{0.17em}3000})`$ of the total image.
One way to boost the data out of the low-count regime is to compare the cumulative radial distribution of the model with the cumulative radial distribution of the data. The modified chi-square-gamma statistic was used to compare the cumulative radial distribution of the model (10 1-pixel-wide bins $``$ 10 degrees-of-freedom) with the cumulative radial distribution of the data (similarly formatted). At the position of $`(27,26)`$ the value of $`\chi _{\gamma \mathrm{M}}^2`$ for the cumulative radial distributions was computed to be 13.7338 ($`\nu 10`$) for a model of a 40 photon point source centered at $`(27,26)`$ on a background of 0.06 photons per pixel. The 95th percentage point for the chi-square distribution with 10 degrees of freedom may be found in several standard references: 18.3 (CRC Handbook of Chemistry and Physics, Lide & Frederikse crc1994 (1995), p. A-106), 18.31 (Bevington be1969 (1969), p. 315) and 18.3070 (Abramowitz & Stegun abst1964 (1964), p. 985). I conclude that the image location $`(27,26)`$ is within the 95% confidence interval because the value of the modified chi-square-gamma statistic is less than the 95th percentage point for the chi-square distribution with 10 degrees of freedom (i.e. $`13.7338<18.3070`$). The probability that the observed chi-square value for a correct model should be less than a value of $`\chi ^2`$ for $`\nu `$ degrees of freedom is $`P(\chi ^2|\nu )=P(\frac{\nu }{2},\frac{\chi ^2}{2})`$ where the latter function is the incomplete gamma function \[ $`P1Q`$; see, e.g., the GAMMP routine in Numerical Recipes (Press et al. pret1986 (1986))\]. If we assume that $`\chi _{\gamma \mathrm{M}}^2`$ is distributed like the $`\chi ^2`$ distribution, then we can assign a probability for the modified chi-square-gamma value for 10 degrees of freedom: $`P(13.7338|10)=P(\frac{10}{2},\frac{13.7338}{2})=0.814517`$. There is thus a $``$81.5% chance that the observed modified chi-square-gamma statistic will be less than 13.7338 for 10 degrees of freedom. The contour in Fig. 19 shows the 95% confidence interval of the X-ray point source based on the $`\chi _{\gamma \mathrm{M}}^2`$ analysis of the cumulative radial distribution of the data. The value of $`\chi _{\gamma \mathrm{M}}^2`$ for the cumulative radial distribution at $`(26,26)`$ was computed to be 30.4707 giving a probability of $``$99.9%; this location in Fig. 19 lies outside the 95% confidence interval.
We can further use the modified chi-square-gamma statistic with the cumulative radial distribution to determine the 95% confidence limits of the intensity of the X-ray source in the image (see Fig. 20 margin: $``$Fig20 ). The upper and lower single-sided 95% confidence limits for the intensity of an X-ray point source at $`(33,26)`$ in Fig. 16 is 54.5 and 28.0, respectively. The true intensity of the X-ray source is 40 photons. Given a background flux uncertainty of $`\sigma _B=0.0049`$ photons per pixel (see above), we can approximate the theoretical rms measurement error for a 40 photon point source spread over 317 pixels $`(A=317`$ px$`{}_{}{}^{2})`$ as $`\sigma \sqrt{40+1}+A\sigma _B8.0`$ photons. The difference between the upper and lower 95% single-sided confidence limits is approximately 3.3 standard deviations of the normal probability function. This fact can be used to approximate an rms measurement error for our intensity estimate of $`\sigma _I8.0`$ \[$`(54.528.0)/(2\times 1.65)`$\] photons. The $`\chi _{\gamma \mathrm{M}}^2`$ analysis using the cumulative radial distribution has yielded an excellent intensity estimate.
The analysis presented in Figures 19 and 20 is predicated on the assumption that the modified chi-square-gamma statistic is distributed like $`\chi ^2`$. But is this assumption valid? Figure 21 margin: $``$Fig21 shows that the analysis of $`10^4`$ simulated X-ray observations like Fig. 16 yields modified chi-square-gamma values that are distributed like the chi-square distribution for $`\nu 317`$ degrees of freedom: a Gaussian distribution with a mean of $`\nu `$ and a variance of $`2\nu `$. The above analysis has assumed that the probability that the observed modified chi-square-gamma value for a correct model should be less than a value of $`\chi ^2`$ for $`\nu `$ degrees of freedom can be given as $`P(\chi _{\gamma \mathrm{M}}^2|\nu )`$. Assuming that the predicted probability from $`P(\chi _{\gamma \mathrm{M}}^2|\nu )`$ is an accurate prediction of the true probability, then the predicted probability of the 9500th simulated observation of a total of 10000 (sorted by $`\chi _{\gamma \mathrm{M}}^2`$ value) should be very close to 95%; Fig. 22 margin: $``$Fig22 indicates that this is indeed the case (i.e., the probability for the $`\chi _{\gamma \mathrm{M}}^2`$ value of the 9500th simulated observation is 94.8666%). Since $`\chi _{\gamma \mathrm{M}}^2`$ statistic is distributed (nearly) like the $`\chi ^2`$ distribution, the usage of the incomplete gamma function to predict probabilities for modified chi-square-gamma values appears to be justified in practical analysis problems.
## 5 SUMMARY
I investigated the use of Pearson’s chi-square statistic \[eq. (3)\], the Maximum Likelihood Ratio statistic for Poisson distributions \[eq. (5)\], and the chi-square-gamma statistic \[eq. (6)\] for the determination of the goodness-of-fit between theoretical models and low-count Poisson-distributed data. I concluded that none of these statistics should be used to determine the goodness-of-fit with data values of 10 or less.
I modified Pearson’s chi-square statistic for the purpose of improving its goodness-of-fit performance. I demonstrated that modified Pearson’s $`\chi ^2`$ statistic \[eq. (9)\] works well in the perfect case where one has a priori knowledge of the correct (true) model. In a real experiment, however, the true mean of the parent Poisson distribution is rarely (if ever) known and model parameters must be estimated from the observations. I demonstrated that the modified Pearson’s $`\chi ^2`$ statistic has a variance that is significantly smaller than that of the $`\chi ^2`$ distribution when realistic models, defined as having parameters estimated from the observational data, are compared with Poisson-distributed data. Any statistic that fails with models based on reasonable parameter estimates is not a very practical statistic for the analysis of astrophysical observations. I concluded that the modified Pearson’s $`\chi ^2`$ statistic should not be used to determine the goodness-of-fit with low-count data values of 10 or less.
I modified the chi-square-gamma statistic for the purpose of improving its goodness-of-fit performance. I demonstrated that the modified chi-square-gamma statistic \[eq. (15)\] performs (nearly) like an ideal $`\chi ^2`$ statistic for the determination of goodness-of-fit with low-count data. On average, for correct (true) models, the mean value of the modified chi-square-gamma statistic is equal to the number of degrees of freedom $`(\nu )`$ and its variance is $`2\nu `$ — like the $`\chi ^2`$ distribution for $`\nu `$ degrees of freedom.
An ideal $`\chi ^2`$ statistic for the determination of goodness-of-fit with low-count data should fail in a predictable manner. Hypothesis testing of low-count Poisson-distributed data with the modified Pearson’s $`\chi ^2`$ statistic will produce the peculiar and undesirable result that correct models are more likely to be rejected than realistic models \[cf. Fig. 6 with Fig. 8\]. The modified chi-square-gamma statistic is a practical statistic to use for hypothesis testing of astrophysical data from counting experiments because it performs (nearly) like an ideal $`\chi ^2`$ statistic for realistic and correct models in the low-count and the high-count data regimes; accurate and believable probabilities for $`\chi _{\gamma \mathrm{M}}^2`$ goodness-of-fit values can be calculated with the incomplete gamma function \[Figs. 21 and 22\]. A lot of nothing can tell you something — as long as there are some observations with signal in them.
Vincent Eke sent me an e-mail asking if I had an expression for the variance of the $`\chi _\gamma ^2`$ statistic which described the mysterious second hump of the solid curve of Fig. 3 of PaperI. After a rapid exchange of email with me over the period of a week, he was the first to derive an analytical formula for $`\sigma _{\chi _\gamma ^2/\mathrm{}}^2`$ \[eq. (14)\]. The knowledge that the variance of $`\chi _\gamma ^2`$ could in fact be expressed explicitly as an analytical expression turned out to be the breakthrough that I had needed in order to complete the development of the modified $`\chi _\gamma ^2`$ statistic. It is a pleasure to acknowledge his contribution to this research. I would like to thank Mike Merrill for the use of his copy of Mathematica which I used to check some of the arithmetic of the critical last step in the derivation of Eq. (14). Special thanks are due to Mary Guerrieri, the NOAO librarian, who has greatly facilitated this research effort by finding and securing loans for many a quaint and curious volume of forgotten lore. I was supported by a grant from the National Aeronautics and Space Administration (NASA), Order No. S-67046-F, which was awarded by the Long-Term Space Astrophysics Program (NRA 95-OSS-16). This research has made use of NASA’s Astrophysics Data System Abstract Service which is operated by the Jet Propulsion Laboratory at the California Institute of Technology, under contract with NASA.
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# The Gromov norm and foliations
## 1. Introduction
The study of group actions on trees and tree–like objects has for a long time been an important tool in $`3`$–manifold topology. J. Stallings pioneered this approach with a topological proof of Grushko’s theorem , and more generally it is by now a standard observation that a decomposition of a $`3`$–manifold along an incompressible surface is “dual” to some action of $`\pi _1(M)`$ on a tree. M. Culler and P. Shalen used the action of $`\pi _1(M)`$ for $`M`$ a hyperbolic manifold on the Bruhat–Tits tree of $`SL(2,F)`$, where $`F`$ is the function field of a curve in the $`SL(2,)`$ representation variety of $`\pi _1`$, to obtain striking topological results about $`M`$. More generally, a sequence of hyperbolic (or merely negatively curved) structures on a fixed manifold which are not precompact in the Gromov–Hausdorff topology may be rescaled and filtered to give in the limit an action of $`\pi _1(M)`$ on an $``$–tree (see e.g. F. Paulin ). In a $`2`$–dimensional setting, this idea is implicit in Thurston’s compactification of Teichmüller space by projective measured laminations .
However, the tree–like structures on which $`\pi _1(M)`$ acts in all these cases admit some kind of invariant measure structure. For a taut foliation or an essential lamination, the existence of such a transverse structure is rare, and leads to strong topological conditions on the underlying manifold. Consequently many “naturally occurring” actions of fundamental groups of $`3`$–manifolds on non–Hausdorff simply connected $`1`$–manifolds and order trees admit no invariant measure. Nevertheless, one would like to quantify the amount of branching of such trees in some natural way. In this paper, we introduce a norm on the homology of a foliated manifold, which is a refinement of the usual Gromov norm on homology, where one restricts the admissible chains representing a homology class to those which are transverse — that is, each singular map $`\sigma :\mathrm{\Delta }^iM`$ in the support of an admissible chain must induce a standard foliation on $`\mathrm{\Delta }^i`$, one which is topologically conjugate to an affine foliation. For a hyperbolic manifold $`M^n`$, the Gromov norm of $`[M]`$ is proportional to the volume of $`M`$, and for $`n3`$, a chain whose norm is close to the infimum actually “detects” the geometry of $`\stackrel{~}{M}=^n`$ (this is just a restatement of Mostow’s rigidity theorem). For a foliated manifold, the tension between the geometry of the ambient manifold and the local affine structure determined by the foliation can be used to show that the foliated norm differs from the usual norm in certain cases, which reflect the topology and the geometry of the foliation.
In particular, we have the following theorems:
Theorem 2.2.10. Let $``$ be a foliation of $`M^n`$ whose universal cover is topologically conjugate to the standard foliation of $`^n`$ by horizontal $`^{n1}`$’s. Then
$$[M]_G=([M],)_{FG}$$
Theorem 2.5.9. Suppose that $``$ is a taut foliation with one–sided branching. Then there is an equality of norms
$$[M]_G=([M],)_{FG}$$
Theorem 2.4.5. Suppose $`M^n`$ is hyperbolic and $``$ is asymptotically separated. Then
$$[M]_G<([M],)_{FG}$$
Here we say that a foliation $``$ of a hyperbolic manifold is asymptotically separated if for some leaf $`\lambda `$ of $`\stackrel{~}{}`$, there are a pair of open hemispheres $`H^+,H^{}^n`$ in the complement of $`\lambda `$ which are separated by $`\lambda `$. We point out that a standard conjecture would imply that a taut foliation of a hyperbolic $`3`$–manifold is asymptotically separated iff $`\stackrel{~}{}`$ has two–sided branching.
It should be mentioned that a norm for foliations with transverse measures was defined by Connes (unpublished) and developed in and . This norm uses generalized simplices which are simplicial in the tangential direction and measure–theoretical in the transverse direction. It has the usual proportionality properties for foliations whose leaves are all locally isometric to a space of constant curvature. The “fundamental cycle” on which this norm is evaluated is really attached to the measured foliation, and not to the underlying manifold. By contrast, our definition is closer in spirit to norms for stratified or decorated spaces.
With our definition, the norm on a homology class is upper semi–continuous as a function of the underlying foliation, in the geometric topology. Since the norm is defined topologically, this gives obstructions for the existence of a family of isotopies of a fixed topological foliation to converge geometrically to some other foliation. The leaf space of the universal cover of a taut foliation is a (typically non–Hausdorff) simply–connected $`1`$–manifold. The non–Hausdorffness comes from branching of the leaf space. This branching can occur in both directions, in only a single direction, or not at all. The taut foliation is said in these three cases to have branching in both directions, to have branching in only one direction, or to be $``$–covered. We show
Corollary 3.1.5. Let $``$ with branching on at most one side and $`𝒢`$ with two–sided branching, be taut foliations of $`M^3`$. Then there is no sequence of isotopies $`𝒢_i`$ of $`𝒢`$ which converges geometrically to $``$.
Using similar techniques, we show
Theorem 3.3.2. Let $`M`$ be a hyperbolic $`3`$–manifold, and $``$ any taut foliation with $`2`$–sided branching. Then there is a geodesic triangulation $`\tau `$ of $`M`$ which cannot be made transverse to $``$. Furthermore, $`\tau `$ cannot be made transverse to $``$ in any finite cover (i.e. $`\tau `$ is not virtually fine).
Problem 3.16 in Kirby’s problem book asks for a reasonable real–valued function on the set of $`3`$–manifolds which measures the complexity of $`\pi _1(M)`$ and behaves appropriately under finite covers and positive degree maps. One may translate this problem into the foliated context, where one considers pairs $`(M,)`$, finite covers and transverse positive degree maps (a map $`f:(M,)(N,𝒢)`$ between foliated manifolds is transverse if every transversal to $``$ is mapped to a transversal to $`𝒢`$), in which context our norm seems like a “reasonable” solution.
I would like to thank I. Agol, A. Casson and W. Thurston with whom I had some interesting discussions about this material. In particular, I. Agol’s work on volumes of hyperbolic $`3`$–manifolds with boundary was particularly inspiring. Furthermore, I received partial support from a Sloan Dissertation Fellowship and from the Clay Mathematical Institute while carrying out work on this paper.
## 2. Foliated norms
### 2.1. Foliations
We give the basic definitions of various kinds of foliations of $`3`$–manifolds. More details are to be found in .
###### Definition 2.1.1.
Let $`M`$ be the subspace of $`^3`$ for which $`z0`$, minus the origin. $`M`$ has a foliation $`\stackrel{~}{}`$ by leaves of the form $`z=\text{const.}`$ which are all planes, except for the leaf $`z=0`$ which is a punctured plane. This foliation is preserved by the dilation $`(x,y,z)(2x,2y,2z)`$ and so descends to a foliation of the solid torus. This is called the Reeb foliation of the solid torus.
###### Definition 2.1.2.
A codimension $`1`$ foliation of a $`3`$–manifold is Reebless if it has no solid torus subsets foliated with a Reeb foliation.
For Reebless foliations, every leaf is incompressible and the ambient manifold is irreducible or covered by $`S^2\times S^1`$ foliated by horizontal spheres. Moreover, every loop transverse to a Reebless foliation is homotopically essential. It follows that a Reebless foliation of a manifold pulls back in a finite covering to a co–orientable foliation, one for which there is a choice of orientation on transversals which is invariant under leaf–preserving isotopy. Equivalently, there is a $`\pi _1(M)`$–invariant orientation on the leaf space in the universal cover.
###### Definition 2.1.3.
A codimension $`1`$ foliation of a $`3`$–manifold is taut if there is a circle in the manifold transverse to the foliation which intersects every leaf.
Every taut foliation is Reebless. The induced foliation $`\stackrel{~}{}`$ of the universal cover $`\stackrel{~}{M}`$ of a Reebless foliation is (topologically) a foliation of $`^3`$ by planes, which is the product of a foliation of $`^2`$ by lines with $``$. Consequently, every plane is properly embedded and separates $`^3`$ into two topological half–spaces.
The leaf space $`L`$ of the universal cover of a taut foliation is a (possibly non–Hausdorff) simply connected $`1`$–manifold. The orientation on this leaf space induces a partial order on the leaves: $`\lambda >\mu `$ iff there is a positively oriented transversal in $`\stackrel{~}{M}`$ from $`\mu `$ to $`\lambda `$. The absence of loops in the leaf space make this partial order well–defined. For readers unfamiliar with the topology of non–Hausdorff $`1`$–manifolds, consult . The “non–Hausdorffness” arises from branching of the leaf–space: an embedded half–open arc in the leaf space might have a countably infinite collection of limiting endpoints. Moreover, this branching might take place at a dense set of points.
###### Definition 2.1.4.
A taut foliation is $``$–covered if the pulled back foliation of the universal cover is topologically conjugate to the standard foliation of $`^3`$ by horizontal planes.
For $`M`$ not finitely covered by $`S^2\times S^1`$, the leaf space of $`\stackrel{~}{}`$ is Hausdorff exactly when $``$ is $``$–covered, for taut $``$.
### 2.2. Gromov norms
###### Definition 2.2.1.
A singular $`n`$ chain in $`M`$ is a finite $``$–linear combination of singular $`n`$–simplices, where a singular $`n`$–simplex is a map $`\sigma :\mathrm{\Delta }^nM`$ from the standard affine $`n`$–simplex into $`M`$. The support of a chain, denoted $`\text{supp}(C)`$ is the set of singular $`n`$–simplices with non–zero coefficients in $`C`$.
Notice that our convention is to assume that the coefficients in our chains are allowed to be in $``$. We assume this without comment in the sequel.
The classical Gromov norm is defined in as follows:
###### Definition 2.2.2.
For $`M`$ an orientable $`n`$–manifold, let $`[M]`$ denote the fundamental class of $`M`$ in $`H_n(M;)`$. The Gromov norm of $`M`$ is the infimum of the $`L_1`$ norm on the singular cycles representing $`[M]`$. That is,
$$[M]_G=\underset{[_ir_i\sigma _i]=[M]}{inf}\underset{i}{}|r_i|$$
###### Definition 2.2.3.
For $`M`$ an orientable $`n`$–manifold and $``$ a codimension $`m`$ foliation, call a singular cycle $`_ir_i\sigma _i`$ transverse if the foliation on the $`n`$–simplex $`\mathrm{\Delta }^n`$ induced from each singular map $`\sigma _i:\mathrm{\Delta }^nM`$ by pulling back $``$ is topologically conjugate to some affine foliation of $`\mathrm{\Delta }^n`$: that is, the foliation by preimages of points obtained from some affine map $`\mathrm{\Delta }^n^{nm}`$. The foliated Gromov norm of $`[M]`$ with respect to $``$ is defined to be
$$([M],)_{FG}=\underset{[_ir_i\sigma _i]=[M];\sigma _i\text{ transverse}}{inf}\underset{i}{}|r_i|$$
###### Remark 2.2.4.
Notice that any map from the vertices of a simplex to $``$ can be extended to an affine map of the simplex to $``$.
###### Theorem 2.2.5.
Let $`M`$ be a $`3`$–manifold and let $``$ be a taut foliation of $`M`$.
1. $`([M],)_{FG}<\mathrm{}`$.
2. $`[M]_G([M],)_{FG}`$.
3. There is a $`K(M)<\mathrm{}`$ such that for any taut foliation $`𝒢`$ of $`M`$, $`([M],𝒢)_{FG}<K(M)`$. If $`K(M)`$ is the infimum of such, define $`[M]_{FG}=K(M)`$.
Proof: Fact 2. is obvious from the definition of the norms. Fact 1. follows from the existence of a triangulation of $`M`$ in which the foliation $``$ can be put into normal form. Fact 3. follows from the stronger theorem of D. Gabai that on any $`3`$–manifold there is a triangulation with respect to which every taut foliation of $`M`$ can be put in normal form. ()
###### Remark 2.2.6.
Of course, one can define any number of norms on homology by restricting the class of singular maps which are deemed admissible. The particular restriction of transversality seems suitable for studying foliations, since it behaves nicely with respect to many of the usual constructions of foliations, e.g. branched covers.
###### Lemma 2.2.7.
Let $`f:N^nM^n`$ be a degree $`d`$ branched cover where the branch locus $`\gamma `$ is transverse to a foliation $``$ of $`M^n`$. Let $`𝒢`$ be the foliation obtained from $``$ by pullback. Then
$$([N],𝒢)_{FG}d([M],)_{FG}$$
Proof: Let $`C=_ir_i\sigma _i`$ be a transverse chain representing $`[N]`$. Then $`f_{}C=_ir_i(f\sigma _i)`$ is a transverse chain representing $`d[M]`$.
It is a well–known fact that the simplex is distinguished amongst all affine polyhedra by the fact that any total ordering of its vertices is induced from an affine map of the entire simplex to $``$. In the context of foliations, this fact has the following generalization:
###### Lemma 2.2.8.
Let $``$ be a foliation of $`M^n`$ and let $`C=_ir_i\sigma _i`$ be a cycle representing $`[M]`$. Suppose that the leaf space of $`\stackrel{~}{}`$ is an acyclic $`1`$–manifold and that the leaves of $`\stackrel{~}{}`$ are all $`^{n1}`$. Suppose further for each $`i`$ that $`\sigma _i`$ lifts to $`\stackrel{~}{\sigma _i}:\mathrm{\Delta }^n\stackrel{~}{M}`$ such that the images of the vertices of $`\mathrm{\Delta }^n`$ inherit a total order from the natural partial order on $`L`$. Then $`C`$ can be “straightened” to $`C^s=_ir_i\sigma _i^s`$ where $`\stackrel{~}{\sigma _i}`$ and $`\stackrel{~}{\sigma _i^s}`$ have the same endpoints in $`\stackrel{~}{M}`$, and $`C^s`$ is a transverse cycle.
Proof: Let $`S_n=\text{supp}(C)`$ and inductively define $`S_i=\text{supp}(S_{i1})`$. Lift each $`\sigma S_i`$ to $`\stackrel{~}{\sigma }:\mathrm{\Delta }^i\stackrel{~}{M}`$ and call the union of some choice of lifts $`\stackrel{~}{S_i}`$.
Each $`\stackrel{~}{\sigma }\stackrel{~}{S_1}`$ maps its vertices to leaves in $`\stackrel{~}{}`$ which inherit a total order from the partial order on $`L`$. (Here we actually require the images of distinct vertices of $`\mathrm{\Delta }^n`$ to lie on distinct leaves of $`\stackrel{~}{}`$). It follows we can replace each $`\sigma S_1`$ by a transverse $`\sigma ^s`$ with the same endpoints.
Let $`\sigma :\mathrm{\Delta }^i\stackrel{~}{}`$ be transverse. Then the induced foliation of $`\mathrm{\Delta }^i`$ is the standard foliation of the unit sphere $`S^{i1}`$ in $`^i`$ by its intersection with horizontal planes. The leaves of this foliation are $`(i2)`$–spheres away from the top and bottom vertex, with respect to the partial ordering on $`L`$. Since leaves of $`\stackrel{~}{}`$ are just $`^{n1}`$’s, this family of maps of $`(i2)`$–spheres extends to a family of maps of $`(i1)`$–balls converging at the top and bottom to the image of the top and bottom vertices respectively. This family of maps gives a transverse map $`\sigma :\mathrm{\Delta }^i\stackrel{~}{}`$ agreeing with $`\sigma `$ on $`\mathrm{\Delta }_i`$.
Thus the straightening procedure can be performed inductively, as required.
###### Remark 2.2.9.
The “nondegeneracy” assumption — that distinct vertices of a simplex get mapped to distinct leaves of $`\stackrel{~}{}`$ under lifts of singular maps $`\stackrel{~}{\sigma }`$ in the support of $`C`$ is not really necessary, since our definition of “transverse” cycle includes degenerate affine maps. In any case, any finite chain can be perturbed chain homotopically to a nondegenerate chain without violating the total ordering assumption.
A taut foliation of $`M^3`$ has the properties required by lemma 2.2.8.
###### Theorem 2.2.10.
Let $``$ be a foliation of $`M^n`$ whose universal cover is topologically conjugate to the standard foliation of $`^n`$ by horizontal $`^{n1}`$’s. Then
$$[M]_G=([M],)_{FG}$$
Proof: Since $`\stackrel{~}{}`$ is the standard foliation of $`^n`$ by horizontal planes, the leaf space of $`\stackrel{~}{}`$ is totally ordered. As remarked earlier, we can perturb a chain $`C`$ by a chain homotopy to a nearby chain $`C^{}`$ so that for each $`\sigma `$ in the support of $`C`$, $`\stackrel{~}{\sigma }`$ maps the vertices of $`\mathrm{\Delta }^n`$ to distinct leaves of $`\stackrel{~}{}`$. It follows from lemma 2.2.8 that any chain can be straightened to a transverse chain.
In particular, this result holds for $``$ an $``$–covered taut foliation of some $`3`$–manifold $`M`$.
###### Remark 2.2.11.
One easily extends this argument to see that the foliated Gromov norm agrees with the usual Gromov norm for foliations whose universal covering foliations are standard foliations of $`^n`$ by horizontal $`^{nm}`$’s — that is, for product–covered foliations.
### 2.3. Measurable chains and equivariant straightening
Theorem 2.2.10 says that finite chains can be straightened with respect to an $``$–covered foliation. An interesting question is whether the same is true of infinite chains. We make this question more precise.
###### Definition 2.3.1.
Let $`\sigma :\mathrm{\Delta }^iM`$ lift to $`\stackrel{~}{\sigma }:\mathrm{\Delta }^i\stackrel{~}{M}`$ which can be projected to $`\tau :\mathrm{\Delta }^iL`$, the leafspace of $`\stackrel{~}{}`$. $`\sigma `$ is monotone if the stratification of $`\mathrm{\Delta }^i`$ by preimages of points in $`L`$ is homotopy equivalent to an affine foliation of $`\mathrm{\Delta }^i`$. If $`\sigma `$ is monotone, every $`\sigma ^{}`$ in the support of $`\sigma `$ is also monotone.
###### Definition 2.3.2.
For each $`i`$, let $`\mathrm{\Sigma }_i`$ denote the space of singular maps $`\sigma :\mathrm{\Delta }^iM`$ with the compact–open topology. Let $`\mathrm{\Sigma }_i^t`$ denote the subspace of transverse singular maps. Let $`\mathrm{\Sigma }_i^m`$ denote the subspace of monotone singular maps.
In we prove the following result:
###### Theorem 2.3.3.
Let $``$ be an $``$–covered foliation of an atoroidal $`3`$–manifold $`M`$. Then there are co–ordinates on the universal cover $`M=^2\times `$ such that leaves of $`\stackrel{~}{}`$ are horizontal planes $`^2\times \text{const}`$ and $`\pi _1(M)`$ acts by elements of $`\text{Homeo}(^2)\times \text{Homeo}()`$.
Using this structure, one has at least the following partial result:
###### Theorem 2.3.4.
Let $``$ be a co–oriented $``$–covered foliation of a negatively curved closed $`3`$–manifold $`M`$. Then there are continuous projections from $`s_i:\mathrm{\Sigma }_i\mathrm{\Sigma }_i^m`$, which are compatible with $``$ in the sense that $`s_{i1}=s_i`$.
Proof: To avoid cumbersome notation, we define instead straightenings of simplicial maps to $`\stackrel{~}{M}`$ which are continuous in the compact–open topology, and which are equivariant under the action of $`\pi _1(M)`$.
For each $`\sigma :\mathrm{\Delta }^i\stackrel{~}{M}`$ we can define two functions $`\rho ,\tau `$ in terms of the co–ordinates on $`\stackrel{~}{M}=^2\times `$ by writing
$$\sigma :t(\rho (t),\tau (t))^2\times $$
We begin by defining $`s_1`$. Let $`\tau (0)=l,\tau (1)=r`$ and let $`J`$ be the set of numbers bounded by $`r`$ and $`l`$. There is an obvious retract $`\varphi :J`$ which sends everything above $`J`$ to the maximum of $`J`$, and everything below $`J`$ to the minimum of $`J`$. Set $`\tau _1(t)=\varphi \tau (t)`$. Now set
$$\tau ^{}(t)=inf(s:\text{there exists }t_1tt_2\text{ with }\tau _1(t_1)=\tau _1(t_2)=s)$$
Then $`\tau ^{}:I`$ is a monotone map, and we can replace the map $`\sigma :t(\rho (t),\tau (t))`$ with $`s_1(\sigma ):t(\rho (t),\tau ^{}(t))`$.
Now for $`\sigma :\mathrm{\Delta }^2\stackrel{~}{M}`$, we first straighten $`\sigma `$ using $`s_1`$. The maps $`\tau :\mathrm{\Delta }^2`$ are already monotone, so for each $`p\mathrm{\Delta }^2`$, there is a unique $`t`$ such that $`p`$ is in the convex hull of the points in $`\mathrm{\Delta }`$ which are mapped by $`\tau `$ to $`t`$. Define $`\tau ^{}(p)`$ to be this value $`t`$, and set $`s_2(\sigma ):p(\rho (p),\tau ^{}(p))`$.
Finally, for $`\sigma :\mathrm{\Delta }^3\stackrel{~}{M}`$, straighten $`\sigma `$ using $`s_2`$. For each $`t`$ in the interior of the image of $`\tau |_{\mathrm{\Delta }^3}`$, the level set $`\tau ^1(t)\mathrm{\Delta }^3`$ is a cellular subset homotopy equivalent to a circle. It has two frontiers in $`\mathrm{\Delta }^3`$, on the side where $`\tau `$ is greater than $`t`$ and the side where $`\tau `$ is less than $`t`$. Define $`D_t\mathrm{\Delta }^3`$ to be the minimal surface spanned by the upper frontier of $`\tau ^1(t)`$ (see e.g. for basic facts about minimal surfaces in $`3`$–manifolds). For distinct $`s,t`$ the disks $`D_t,D_s`$ are disjoint, so we can define $`\tau _{}^{}{}_{}{}^{1}(t)`$ on $`\mathrm{\Delta }^3`$ to be the subset of points above $`D_s`$ for all $`s<t`$ and contained in or below $`D_t`$. Then set $`s_3(\sigma ):p(\rho (p),\tau ^{}(p))`$.
These constructions use only the order structure of $``$, and are therefore equivariant under the action of $`\pi _1(M)`$.
### 2.4. The norm is non–trivial on $`[M]`$
For $`M^n`$ hyperbolic, we know $`[M]_G=\text{vol}(M)/v_n`$ where $`v_n`$ is the volume of the regular ideal $`n`$–simplex. For a hyperbolic manifold, any cycle can be chain homotoped to a geodesic cycle by replacing each singular map of a simplex $`\sigma _i:\mathrm{\Delta }^nM`$ with the geodesic simplex $`\sigma _i^g:\mathrm{\Delta }^nM`$ having the same endpoints.
Let $`C_j`$ be a sequence of geodesic chains whose norms converge to the Gromov norm of $`M`$. Let $`\stackrel{~}{C_j}`$ denote the $`\pi _1(M)`$–equivariant infinite chain obtained by lifting $`C_j`$ to $`\stackrel{~}{M}`$. Let $`X`$ be the infinite $`(n+1)`$–valent tree with basepoint, and let $`T`$ be the abstract complex obtained by gluing together infinitely many ideal $`n`$–simplices along their faces in the pattern described by $`X`$. Fix some regular ideal simplex $`\mathrm{\Delta }^n`$. Then choosing an identification of $`\mathrm{\Delta }`$ with some simplex of $`T`$, there is a natural developing map $`\text{dev}:T^n`$ taking each simplex of $`T`$ to a regular ideal simplex. If $`n=3`$ this has as its image the standard regular tessellation of $`^3`$ by ideal simplices. Otherwise, the representation $`Aut(T)Isom^+(^n)`$ is indiscrete.
###### Lemma 2.4.1.
With notation as above, for any $`t,ϵ>0`$ there is a $`j`$ such that for any $`k>j`$, there is a collection $`S_k`$ of singular maps in the support of $`\stackrel{~}{C_k}`$ and an element $`\alpha _kIsom^+(^n)`$ such that $`S_k`$ and $`\alpha _k\text{dev}(T)`$ agree on the ball of radius $`t`$ about $`0`$ to within $`ϵ`$.
Proof: Let $`C_k=_ir_i\sigma _i`$ be a geodesic chain which very nearly realizes the Gromov norm of $`M`$. Then by definition
$$\underset{i}{}|r_i|<\text{vol}(M)/v_n+\delta $$
for some small $`\delta `$. On the other hand,
$$\underset{i}{}r_i\text{vol}(\sigma _i(\mathrm{\Delta }^n))=\text{vol}(M)$$
so the weighted average
$$\frac{_ir_i\text{vol}(\sigma _i(\mathrm{\Delta }^n))}{_i|r_i|}>v_n\delta ^{}$$
for some small $`\delta ^{}`$; that is, “most” of the $`\sigma _i(\mathrm{\Delta }^n)`$, as weighted by $`r_i`$, have volume very close to $`v_n`$. This implies that they are geometrically very close to regular ideal simplices, on a big compact set containing most of their mass. (see e.g. )
Fix a fundamental domain $`D`$ in $`\stackrel{~}{M}`$. For each $`\sigma \text{supp}(C_k)`$, choose a lift $`\stackrel{~}{\sigma }`$ of $`\sigma `$ whose center of gravity is in $`D`$. Since the bundle of frames over $`D`$ is compact, there is some ideal tetrahedron $`\mathrm{\Delta }^{}`$ in $`^n`$ with center of mass in $`D`$ for which some definite mass of $`\stackrel{~}{\sigma }`$ is geometrically close to $`\mathrm{\Delta }^{}`$ on a big set. Call $`S`$ the set of lifts sufficiently close to $`\mathrm{\Delta }^{}`$. It follows that $`S`$ is geometrically close to $`\mathrm{\Delta }`$. Since $`\stackrel{~}{C_k}=0`$, most of the mass of this boundary must be absorbed by simplices which are geometrically close to being regular and ideal.
It follows that for $`\delta `$ sufficiently small, we can find lifts $`\stackrel{~}{\sigma }`$ whose images are close on a big set to the ideal simplices obtained by reflecting $`\mathrm{\Delta }^{}`$ in each of its boundary faces. Let $`\alpha _k(\mathrm{\Delta })=\mathrm{\Delta }^{}`$. Continuing inductively, if we propagate outwards in $`X`$ until we cover a big ball, we can find corresponding simplices in $`\text{supp}(\stackrel{~}{C_k})`$ which agree with the simplices in $`\alpha _k\text{dev}(T)`$ to within a suitable tolerance, by taking $`\delta `$ sufficiently small.
###### Corollary 2.4.2 (Jungreis).
For a hyperbolic $`M^n`$ with $`n3`$, a regular ideal simplex $`\mathrm{\Delta }^n`$ and any sequence $`C_j`$ of geodesic chains whose norms converge to $`M_G`$, for sufficiently large $`j`$ there is a $`\sigma _j`$ in the support of $`C_j`$ which lifts to a geodesic simplex arbitrarily close to $`\mathrm{\Delta }`$.
Proof: By the previous lemma, there are simplices in the support of $`\stackrel{~}{C_j}`$ which stay close to $`\alpha _k\text{dev}(T)`$ on a big ball about some fixed point. We can identify the set of framed ideal regular simplices in $`^n`$ with $`Isom^+(^n)`$. In the limit, the set of ideal regular simplices in the support of $`\stackrel{~}{C_j}`$ is invariant under the action of $`\pi _1(M)`$ on the left and $`\text{dev}(Aut(T))`$ on the right. If $`n>3`$, $`\text{dev}(Aut(T))`$ is already dense in $`Isom^+(^n)`$. In dimension $`3`$, following Jungreis and Ratner, using both the left and right actions we can find simplices in the support of $`C_j`$ arbitrarily close to $`\mathrm{\Delta }`$ for $`j`$ sufficiently large. See for details.
###### Definition 2.4.3.
Say that a foliation $``$ of a hyperbolic $`n`$–manifold is asymptotically separated if for some leaf $`\lambda `$ of $`\stackrel{~}{}`$, there are a pair of open hemispheres $`H^+,H^{}^n`$ in the complement of $`\lambda `$ which are separated by $`\lambda `$.
###### Example 2.4.4.
For $``$ a finite depth foliation which is not a perturbation of a surface bundle over a circle, the compact leaves lift to quasigeodesically embedded planes in $`^3`$. Hence every leaf has the separation property, and $``$ is asymptotically separated.
###### Theorem 2.4.5.
Suppose $`M^n`$ is hyperbolic and $``$ is asymptotically separated. Then
$$[M]_G<([M],)_{FG}$$
Proof: By passing to a finite cover if necessary, we can assume that $``$ is co–oriented. If we can show that every chain whose norm is sufficiently close to the Gromov norm contains an edge whose endpoints are not joined by an arc transverse to $``$, then we will be done.
Let $`\lambda `$ be a leaf of $`\stackrel{~}{}`$ and $`H^+,H^{}`$ a pair of hemispheres in the complement of $`\lambda `$ as provided by the definition of asymptotically separated. These determine a pair of disks $`D^+,D^{}S_{\mathrm{}}^{n1}`$ above and below $`\lambda `$ respectively. Let $`\alpha \pi _1(M)`$ be an element taking the complement of $`D^+`$ inside $`D^+`$. Then any infinite line from $`D^{}`$ to $`\alpha (D^{})`$ must fail to be transverse to $`\stackrel{~}{}`$ somewhere, since when it crosses $`\lambda `$ it is going in the positive direction, and when it crosses $`\alpha (\lambda )`$ it is going in the negative direction, with respect to the co–orientation on $`\stackrel{~}{}`$ which is preserved by $`\alpha `$.
It is easy to find an ideal regular simplex $`\mathrm{\Delta }`$ which has a pair of endpoints in $`D^{}`$ and $`\alpha (D^{})`$ respectively. For any chain $`C_j`$ with norm sufficiently close to $`[M]_G`$, there is a $`\sigma `$ in the support of $`C_j`$ whose geodesic representative stays very close to $`\mathrm{\Delta }`$ on an arbitrarily large compact piece. Such a $`\sigma `$ has endpoints on incomparable leaves, and therefore cannot be straightened (keeping its endpoints fixed) to a transverse simplex. It follows that no transverse chain can have norm too close to $`[M]_G`$, and the strict inequality is proved.
### 2.5. Limit sets of leaves of taut foliations
To investigate the asymptotic separation property, we must investigate the limit sets of leaves of taut foliations.
###### Lemma 2.5.1.
Let $``$ be taut, and let $`\lambda `$ be a leaf of $`\stackrel{~}{}`$ on which we have chosen a co–orientation. Denote by $`\lambda _{\mathrm{}}`$ the limit set of $`\lambda `$. Then each region $`D`$ in the complement of $`\lambda _{\mathrm{}}`$ is either above or below $`\lambda `$, in the sense that for any two sequences $`\{p_i\}^3`$ and $`\{q_i\}^3`$ with $`p_ipD`$ and $`q_iqD`$, the points $`p_i,q_i`$ are eventually on the same side of $`\lambda `$.
Proof: The points $`p`$ and $`q`$ can be joined by an arc $`\alpha `$ in $`S_{\mathrm{}}^2`$ which avoids $`\lambda _{\mathrm{}}`$. This arc $`\alpha `$ is the Hausdorff limit in $`\overline{^3}`$ of a sequence of arcs $`\alpha _i`$ joining $`p_i`$ to $`q_i`$. If each of the $`\alpha _i`$ intersected $`\lambda `$, this would give rise to a sequence of points in $`\lambda `$ converging to some point in $`\alpha `$, contrary to the hypothesis that $`\alpha `$ avoids $`\lambda _{\mathrm{}}`$. It follows that $`p_i`$ and $`q_i`$ are eventually on the same side of $`\lambda `$, and therefore the “side” of $`D`$ is unambiguously defined.
Notice that “above” and “below” as defined in the previous lemma are not the same as $`<`$ and $`>`$ in the partial order on $`L`$. Each leaf in the universal cover of a taut foliation has two sides; a co–orientation on the leaf defines one of the sides to be above and one below, and every other leaf falls into one of these two possibilities. This does not define a partial ordering on leaves.
###### Definition 2.5.2.
Say that a foliation $``$ has two–sided branching if in the partial order on the leaf space $`L`$ of $`\stackrel{~}{}`$, there are triples of leaves $`\lambda ,\lambda _l^+,\lambda _r^+`$ and $`\mu ,\mu _l^{},\mu _r^{}`$ such that
$$\lambda <\lambda _l^+,\lambda <\lambda _r^+$$
$$\lambda _l^+\text{ and }\lambda _r^+\text{ are incomparable}$$
$$\mu _l^{}<\mu ,\mu _r^{}<\mu $$
$$\mu _l^{}\text{ and }\mu _r^{}\text{ are incomparable}$$
Observe that if $``$ is taut and has two–sided branching, then we may choose any leaf as $`\mu =\lambda `$. Moreover, if $``$ is not co–orientable, or covers some foliation which is not co–orientable, then either $``$ is $``$–covered or it has two–sided branching.
###### Example 2.5.3.
Let $``$ be a foliation of $`T^3`$ with one horizontal torus leaf, and the complementary $`T^2\times I`$ foliated as a Reeb foliation of the annulus $`\times S^1`$. Then a pair of transversals whose initial segments agree and cross the horizontal torus leaf must thereafter be leafwise homotopic; that is, there is no branching in the positive direction from that point on. Thus, we cannot choose $`\mu =\lambda `$ in this example.
###### Example 2.5.4.
In G. Meigniez constructs examples of taut foliations which branch on only one side, say the negative side. Furthermore, some of these examples are obtained as perturbations of surface bundles over circles, and therefore have pseudo–Anosov flows transverse to them. In we construct new examples of such foliations, and show that this situation holds in general: taut foliations of atoroidal $`3`$–manifolds with one–sided branching have transverse pseudo–Anosov flows which are regulating: that is, flow lines are properly embedded in the leaf space of the universal cover.
Easy examples of foliations with one–sided branching are obtained by starting with $``$–covered foliations with (approximately) projectively invariant transverse measures, and then taking branched covers over a curve which lifts to a line in $`\stackrel{~}{M}`$, one end of which is properly embedded in the leaf space and one end of which is not.
###### Theorem 2.5.5.
Let $``$ be a taut foliation. If $``$ has two–sided branching and is not asymptotically separated, then every leaf $`\lambda `$ of $`\stackrel{~}{}`$ has $`m(\lambda _{\mathrm{}})>0`$, where $`m`$ is some normalized Lebesgue measure on the sphere at infinity $`S_{\mathrm{}}^2`$ of $`^3`$. In fact, $`\lambda _{\mathrm{}}`$ must have non–empty interior.
Proof: Assume without loss of generality that $``$ is co–oriented.
Suppose that $`m(\lambda _{\mathrm{}})=0`$ for some $`\lambda `$. Then certainly there is some complementary domain to $`\lambda _{\mathrm{}}`$ in $`S_{\mathrm{}}^2`$. If there are domains $`D^\pm `$ both above and below $`\lambda `$, in the sense of lemma 2.5.1 then there are half–spaces $`H^\pm `$ bounded by circles in $`D^\pm `$ which avoid $`\lambda `$, and $``$ is asymptotically separated. Otherwise without loss of generality, all the complementary regions to $`\lambda _{\mathrm{}}`$ are contained above $`\lambda `$. It follows that the subset of $`\stackrel{~}{M}`$ below $`\lambda `$ has limit set exactly equal to $`\lambda _{\mathrm{}}`$.
If $``$ is $``$–covered, one knows that $`\lambda _{\mathrm{}}=S_{\mathrm{}}^2`$ for every $`\lambda `$, so we may assume $``$ is not $``$–covered. (see e.g. )
If $``$ has two–sided branching, then there are a pair of positive transversals to $`\stackrel{~}{}`$ emanating from $`\lambda `$ and ending on two incomparable translates $`\alpha (\lambda ),\beta (\lambda )`$ both $`>\lambda `$ in the partial order on $`L`$. Now, the subset of $`\stackrel{~}{M}`$ below $`\alpha (\lambda )`$ has limit set equal to $`\alpha (\lambda _{\mathrm{}})`$, which has measure $`0`$. However, the subset of $`\stackrel{~}{M}`$ above $`\beta (\lambda )`$ is itself a subset of the subset of $`\stackrel{~}{M}`$ below $`\alpha (\lambda )`$. It follows that we can write $`^3`$ as the union of two sets (the sets above and below $`\beta (\lambda )`$), each of which has a limit set of measure $`0`$, which is absurd. More generally, if $`\lambda _{\mathrm{}}`$ has no interior, we could write $`S_{\mathrm{}}^2`$ as the union of two closed sets without interior, which is absurd.
The following theorem is proved in :
###### Theorem 2.5.6 (Fenley).
Let $``$ be a Reebless foliation in $`M^3`$ closed, hyperbolic. Suppose that $`\lambda _{\mathrm{}}S_{\mathrm{}}^2`$ for some $`\lambda `$, and assume that there is branching in the positive and negative directions of $`\stackrel{~}{}`$. Then there is a $`k<2`$ such that the limit set of every leaf has Hausdorff dimension less than $`k`$. In particular, every such limit set has zero Lebesgue measure.
###### Corollary 2.5.7.
No leaf in the universal cover of a taut foliation of a hyperbolic $`3`$–manifold with two–sided branching is quasi–isometric (as a subset of $`^3`$) to a totally degenerate surface group.
Proof: The limit set of a totally degenerate surface group is a dendrite: a closed set of measure $`0`$ whose complement is connected (see e.g. ). If the limit set of a leaf of a taut foliation has measure $`0`$, it has at least $`2`$ complementary regions: one above and one below.
###### Corollary 2.5.8.
Let $``$ have two–sided branching. Then either $`\lambda _{\mathrm{}}=S_{\mathrm{}}^2`$ for every leaf $`\lambda `$ of $`\stackrel{~}{}`$ or every leaf is asymptotically separated, and the foliated norm of $`[M]`$ is strictly greater than the usual norm.
Fenley has conjectured that $`\lambda _{\mathrm{}}=S_{\mathrm{}}^2`$ iff $``$ is $``$–covered.
By contrast, we have the following theorem:
###### Theorem 2.5.9.
Suppose that $``$ is a taut foliation with one–sided branching. Then there is an equality of norms
$$[M]_G=([M],)_{FG}$$
Proof: We declare that the branching takes place in the positive direction, with respect to some choice of co–orientation on $`\stackrel{~}{}`$.
Lift the singular maps in the support of a chain $`C=_ir_i\sigma _i`$ to maps $`\stackrel{~}{\sigma _i}:\mathrm{\Delta }^3\stackrel{~}{M}`$. It is possible that some vertices of $`\mathrm{\Delta }^3`$ are mapped by some $`\stackrel{~}{\sigma _i}`$ to incomparable leaves of $`\stackrel{~}{}`$. However, any pair of points in $`\stackrel{~}{M}`$ can be made comparable after a finite isotopy in the negative direction. Moreover, if two points are already on comparable leaves, then they are still on comparable leaves after such an isotopy. Since there are only finitely many $`\stackrel{~}{\sigma _i}`$, we can push the images of the vertices under $`\sigma _i`$ in the negative direction to get a new chain, homotopic to $`C`$, for which each $`\stackrel{~}{\sigma _i}`$ sends the vertices of $`\mathrm{\Delta }^3`$ to comparable leaves of $`\stackrel{~}{}`$. By lemma 2.2.8, we can straighten this new chain relative to its vertices to be transverse to $``$.
###### Corollary 2.5.10.
If $``$ has one–sided branching, the leaves of $`\stackrel{~}{}`$ are not asymptotically separated.
### 2.6. Foliations with Reeb components
One might suppose at least that the existence of Reeb components should be detected by the foliated Gromov norm. However, this is not the case, as the following example shows.
###### Example 2.6.1.
Let $`M=S^2\times S^1`$ and $``$ the standard foliation by two Reeb components. Let $`C`$ be any chain representing $`M`$. Let $`\pi :MM`$ be the unique connected double cover of $`M`$. Then $`\pi ^{}=`$ up to isotopy. However, $`\pi _{}C=2[M]`$, so the sequence $`2^n\pi _{}^nC`$ of chains can be made transverse after isotopy and have norm $`0`$. Hence
$$[M]_G=([M],)_{FG}=0$$
in this case.
Despite this example, it is easy to arrange a sequence of foliations of a given manifold $`M`$ with more and more Reeb components where the foliated Gromov norm grows without bound, as the following theorem shows.
###### Theorem 2.6.2.
There is a function $`k:`$ with $`lim_n\mathrm{}k(n)=\mathrm{}`$ such that if $``$ is any foliation of a $`3`$–manifold $`M`$ with $`n`$ generalized Reeb components whose complement is atoroidal, then
$$([M],)_{FG}k(n)$$
Proof: A generalized Reeb component, also known as a dead end component, is a region of the foliated manifold bounded by torus or Klein bottle compact leaves, such that no path transverse to the foliation which enters the component can leave again.
Let $`C=_ir_i\sigma _i`$ be a transverse chain representing $`[M]`$ and let $`NM`$ be a dead end component. Suppose that $`\sigma _i`$ is a singular simplex whose image intersects $`N`$. Let $`p\mathrm{\Delta }^3`$ be such that $`\sigma _i(p)N`$ and let $`\alpha `$ be a path in $`\mathrm{\Delta }^3`$ transverse to the foliation induced by $`\sigma _i^1()`$ running from the top to the bottom vertices which passes through $`p`$. Then $`\sigma _i(\alpha )`$ is transverse to $``$ and intersects $`N`$; it follows that the image of at least one of the vertices of $`\mathrm{\Delta }^3`$ must be contained in $`N`$.
If we “truncate” $`M`$ by removing the Reeb components, we get a $`3`$–manifold $`M^{}`$ with at least $`n`$ torus or Klein bottle cusps. By the fact above, each truncated simplex can be collapsed to an edge or a face, or else is a normal simplex possibly with some ideal points. The resulting truncated chain $`C`$ represents $`[M^{}]H_3(M^{},M^{};)`$, so the foliated Gromov norm of $`M`$ can be estimated by the usual Gromov norm of $`M^{}`$. Then set $`k(n)`$ equal to the minimum Gromov norm of a hyperbolic $`3`$–manifold with $`n`$ cusps.
After the work of Thurston (see e.g. ), one knows that $`lim_n\mathrm{}k(n)\mathrm{}`$.
## 3. Extending the norm to $`H_{}(M;)`$
### 3.1. Semicontinuity of the norm
It is clear that the definition of the foliated Gromov norm can be extended to a norm on $`H_i(M;)`$ for a manifold $`M`$ foliated by $``$. As before, for each homology class $`\mu `$ we can consider transverse singular chains representing $`\mu `$, and take the $`L_1`$ norm of such representatives. Denote the value of this norm on a class $`\mu `$ by $`(\mu ,)_{FG}`$.
Notice that unlike the usual Gromov norm, this norm may be non–trivial even on $`H_1(M;)`$, as the following example shows:
###### Example 3.1.1.
Let $``$ be the foliation of $`T^2\times I`$ obtained by multiplying a Reeb foliation of the cylinder $`S^1\times I`$ by $`S^1`$. Glue the top and bottom of $`T^2\times I`$ together to get a foliation of $`T^3`$ also denoted by $``$. Let $`\alpha H_1(M;)`$ be the generator obtained from the $`I`$ factor by the gluing. Let $`\beta H_1(M;)=r\alpha `$. Then $`(\beta ,)_{FG}r/2`$. For, each $`\stackrel{~}{\sigma }:\mathrm{\Delta }^1^3`$ obtained by lifting a map in the support of a chain representing $`\beta `$ must have length $`2`$ in its projection to the vertical factor, since such a chain cannot cross the torus leaf twice.
###### Theorem 3.1.2.
Let $`_i`$ be a sequence of taut foliations of $`M^n`$ which converge geometrically (as $`(n1)`$–plane fields) to $``$. Then
$$(A,)_{FG}lim\; sup(A,_i)_{FG}$$
for any $`AH_{}(M;)`$.
Proof: Let $`C_j=_ir_{ij}\sigma _{ij}`$ be a sequence of cycles transverse to $``$ representing $`A`$ whose norms converge to $`(A,)_{FG}`$. Then for any $`j`$, every $`\sigma _{ij}`$ is transverse to $``$ and therefore by compactness there is an $`ϵ_j`$ such that the $`1`$–skeleta of the images of $`\mathrm{\Delta }^n`$ under the $`\sigma _{ij}`$ make an angle of at least $`ϵ`$ with $``$ everywhere. It follows that for sufficiently large $`k`$, the $`1`$–skeleta of $`\text{supp}(C_j)`$ is transverse to $`_k`$. By lemma 2.2.8 we can straighten $`C_j`$ to be transverse to $`_k`$.
This implies that the foliated Gromov norm is lower semi–continuous: the norm can jump up at a limit, but never down. The following example shows, however, that the norm is not actually continuous.
###### Example 3.1.3.
Suppose $`M^3`$ is hyperbolic and fibers over the circle, and has $`b_22`$. Let $`_i`$ be a sequence of fiberings contained in some top dimensional face of the Thurston norm converging to some foliation $``$ which is at a vertex. Then $`([M],_i)_{FG}=[M]_G`$ by theorem 2.2.10. On the other hand, $``$ contains a quasigeodesically embedded compact leaf, so $`([M],)_{FG}>[M]_G`$ by theorem 2.4.5.
An interesting phenomenon in the theory of foliations occurs when a sequence of isotopies of a fixed foliation $``$ converges geometrically to a topologically distinct foliation $`𝒢`$. We give a simple example of this phenomenon.
###### Example 3.1.4.
Let $`S`$ be the cylinder $`I\times S^1`$ foliated by horizontal circles $`\text{point}\times S^1`$. For an end–preserving homeomorphism $`f:II`$ we can produce a foliation $`_f`$ of $`T^2\times I`$ which is the suspension of the foliation of $`S`$ by the map $`f\times \text{id}:SS`$. Any two topologically conjugate maps $`f,g:II`$ give isotopic foliations. Now, it is well–known that any two strictly increasing homeomorphisms of the open interval to itself are topologically conjugate. One can easily find a sequence $`f_i`$ of these which converge (as maps $`II`$) to the identity. The foliations $`_{f_i}`$ are all isotopic, but distinct from $`_{\text{id}}`$.
In any case, upper–semicontinuity of the norm implies the following
###### Corollary 3.1.5.
Let $`,𝒢`$ be taut foliations of $`M`$ and suppose
$$(A,)_{FG}>(A,𝒢)_{FG}$$
for some $`A`$. Then no sequence of isotopies of $``$ can converge geometrically to $`𝒢`$.
Dual to the $`L_1`$ norm on $`C_{}(M;)`$ defined by a foliation, there is an $`L_{\mathrm{}}`$ norm on $`C^{}(M;)`$ defined as the supremum of the value of the cochain on transverse singular maps. There is an associated foliated bounded cohomology, denoted by $`H_{}^{}(M;)`$. This may contain nontrivial elements even in dimension $`1`$, in contrast with the usual bounded cohomology.
### 3.2. Length of a free homotopy class
There is a homotopy–theoretic refinement of the norm on $`H_1`$. Say that a free homotopy class $`[\alpha ]`$ of loops is transverse to $``$ if $`\alpha `$ is freely homotopic to a transverse circle. More generally, define the length of $`[\alpha ]`$, denoted $`\mathrm{}([\alpha ])`$, to be the minimum number of subdivisions of $`S^1`$ needed to make a representative transverse on each subdivision. Say that this length is $`0`$ if no subdivision is necessary: that is, if some loop representing $`\alpha `$ is either transverse as a circle to $``$ or can be homotoped into a leaf of $``$.
Notice that for a co–oriented foliation, $`\mathrm{}`$ takes on only even values.
###### Lemma 3.2.1.
The length of a free homotopy class of loops $`[\alpha ]`$ is upper semi–continuous in the geometric topology.
Proof: The only non–obvious point to check is that if $`[\alpha ]`$ has a representative which is contained in a leaf of $``$, then $`\mathrm{}([\alpha ])=0`$ for any $`𝒢_i`$ sufficiently close to $``$ in the geometric topology. Lift $`\alpha `$ to an arc $`\stackrel{~}{\alpha }`$ in $`\stackrel{~}{M}`$ contained in a leaf of $`\stackrel{~}{}`$. Then there is a big ball $`B`$ containing $`\stackrel{~}{\alpha }`$ which is foliated in a standard way by $`\stackrel{~}{}`$. Therefore, for any $`𝒢_i`$ sufficiently close to $``$ in the geometric topology, $`\stackrel{~}{𝒢_i}`$ foliates some slightly smaller ball, also containing $`\alpha `$, in a standard way. The arc $`\alpha `$ may be made transverse or tangential to $`\stackrel{~}{𝒢_i}`$ in this ball, implying that $`\mathrm{}([\alpha ])=0`$ for $`𝒢_i`$.
###### Lemma 3.2.2.
For any co–oriented taut foliation $``$ which has two–sided branching, there is an $`[\alpha ]\pi _1(M)`$ with $`\mathrm{}([\alpha ])2`$.
Proof: Let $`\tau _1,\tau _2`$ be two positive transversals to $`\stackrel{~}{}`$ emanating from the same point whose upper endpoints are on incomparable leaves. Let $`\sigma _1,\sigma _2`$ be two negative transversals to $`\stackrel{~}{}`$ emanating from the same point whose lower endpoints are on incomparable leaves.
Map $`\tau _1\tau _2`$ and $`\sigma _1\sigma _2`$ downstairs to $`M`$. Since $``$ is taut, one can extend the image of $`\tau _1`$ in the positive direction by an arc $`\rho _1`$ until it joins up with $`\sigma _1`$, and do the same with $`\tau _2`$. The union makes up a loop $`\alpha `$, consisting of two transverse arcs $`\alpha _1=\tau _1\rho _1\sigma _1`$ and $`\alpha _2=\tau _2\rho _2\sigma _2`$. We orient $`\alpha `$ so that $`\alpha _1`$ is positive and $`\alpha _2`$ is negative. Let $`\stackrel{~}{\alpha }`$ be a lift to $`\stackrel{~}{M}`$, and consider its projection to the leaf space $`L`$ of $`\stackrel{~}{}`$: this consists of an alternating sequence of positive and negative arcs, which pass over a branch of $`L`$ at each stage. If $`t_i`$ and $`b_i`$ denote the alternating sequence of top and bottom leaves of the projection, then for all $`i`$,
$$t_i>b_i,t_i>b_{i1}$$
$$b_i\text{ and }b_{i+1}\text{ are incomparable}$$
$$t_i\text{ and }t_{i+1}\text{ are incomparable}$$
We can find a sequence of points $`m_i`$ with $`b_i<m_i<t_i`$ and each $`m_i,b_{i1}`$ and $`m_i,t_{i+1}`$ pairwise incomparable. It is clear that for any $`\alpha ^{}`$ homotopic to $`\alpha `$, the projection of the corresponding lift $`\stackrel{~}{\alpha ^{}}`$ to $`L`$ must intersect each $`m_i`$. In particular, such a lift intersects incomparable leaves of $`\stackrel{~}{}`$, so $`\alpha ^{}`$ cannot be transverse. See figure 3.
The following corollary answers a question posed by W. Thurston:
###### Corollary 3.2.3.
Let $``$ with branching on at most one side and $`𝒢`$ with two–sided branching, be taut foliations of $`M`$. Then there is no sequence of isotopies $`𝒢_i`$ of $`𝒢`$ which converges geometrically to $``$.
Proof: Lift to a finite cover where the foliations are co–oriented. Then observe that for a foliation with branching on at most one side, the length of any free homotopy class is $`0`$. For, suppose $``$ does not branch in the positive direction, and let $`\alpha `$ be any loop in $`M`$, and suppose $`\alpha `$ is in general position with respect to $``$ so that it has a finite number of isolated minima and maxima. Let $`p`$ be such a minimum, lying between maxima $`q,r`$. Lift the segment $`\tau `$ between $`q,r`$ to $`\stackrel{~}{\tau }`$ in $`\stackrel{~}{M}`$. Then $`\stackrel{~}{q},\stackrel{~}{r}`$ are comparable, since they are both $`>\stackrel{~}{p}`$ and by hypothesis there is no branching in the positive direction, so without loss of generality we can assume $`\stackrel{~}{q}\stackrel{~}{r}`$. So we can push the local minima corresponding to $`p`$ in the positive direction until it cancels the local maxima corresponding to $`q`$ without introducing any new critical points; in particular, the number of critical points can be reduced. Continuing inductively, they can all be eliminated.
### 3.3. Virtually fine triangulations
###### Definition 3.3.1.
Let $`M`$ be a $`3`$–manifold. A triangulation $`\tau `$ of $`M`$ is fine if for every taut foliation $``$ of $`M`$ $`\tau `$ can be isotoped to be transverse to $``$. A triangulation $`\tau `$ is virtually fine if for every taut foliation $``$ of $`M`$ there is a finite cover $`\widehat{M}`$ of $`M`$ such that the pulled–back triangulation $`\widehat{\tau }`$ can be isotoped to be transverse to the pulled–back foliation $`\widehat{}`$.
One of the main theorems of states that for any $`M`$ there is a fine triangulation $`\tau `$. This leads naturally to the question of what conditions are necessary and sufficient on a triangulation to be fine. An obvious condition is that the triangulation admit a transverse foliation locally. For a geodesic triangulation of a hyperbolic manifold, this is obvious, since in the projective model of hyperbolic space, a hyperbolically geodesic triangulation looks like a Euclidean triangulation of the ball, and a foliation by horizontal planes will be transverse. It has been an open question whether every geodesic triangulation of a hyperbolic $`3`$–manifold $`M`$ is virtually fine.
It turns out this guess is incorrect: there are geodesic triangulations whose simplices have diameters arbitrarily small compared to the injectivity radius of the ambient manifold, which cannot be made transverse to certain taut foliations.
###### Theorem 3.3.2.
Let $`M`$ be a hyperbolic $`3`$–manifold, and $``$ any taut foliation with $`2`$–sided branching. Then there is a geodesic triangulation $`\tau `$ of $`M`$ which cannot be made transverse to $``$. Furthermore, $`\tau `$ cannot be made transverse to $``$ in any finite cover (i.e. $`\tau `$ is not virtually fine).
Proof: Let $`\gamma `$ be a closed loop with $`\mathrm{}(d)=l2`$. Then we can choose a geodesic representative of $`\gamma `$ and make it an edge of a geodesic triangulation. Such a triangulation can obviously not be made transverse to $``$. Now, for any finite cover $`\widehat{M}`$ of $`M`$, we can lift $`\gamma `$ to some $`\widehat{\gamma }`$ which covers $`\gamma `$ with degree $`d`$, and has $`d`$ segments in the lifted triangulation. With respect to $`\widehat{}`$, the new length of $`\widehat{\gamma }`$ is $`ld`$, so the lifted triangulation cannot be made transverse to the lifted foliation.
## 4. Laminations and order trees
### 4.1. Genuine laminations
Laminations of $`3`$–manifolds are defined in .
###### Definition 4.1.1.
A lamination in a $`3`$–manifold is a foliation of a closed subset of $`M`$ by $`2`$ dimensional leaves. The complement of this closed subset falls into connected components, called complementary regions. A lamination is essential if it contains no spherical leaf or torus leaf bounding a solid torus, and furthermore if $`C`$ is the closure (with respect to the path metric in $`M`$) of a complementary region, then $`C`$ is irreducible and $`C`$ is both incompressible and end incompressible in $`C`$. Here an end compressing disk is a properly embedded $`(D^2(\text{closed arc in }D^2))`$ in $`C`$ which is not properly isotopic relative to the $``$ in $`C`$ to an embedding into a leaf. Finally, an essential lamination is genuine if it has some complementary region which is not an $`I`$-bundle.
An essential lamination simultaneously generalizes both Reebless foliations and incompressible surfaces. It is not true that an essential lamination lifts in a finite cover to a co–orientable lamination. Consequently the leaf space of an essential lamination in the universal cover does not carry a natural partial order. The leaf space of a foliation is like a train–track: there is a natural combing near any branch point. The leaf space of a lamination is more like a tree: there is no natural way to say whether branches approach a branch point from the same or from opposite directions.
Nevertheless, we can still talk about transversality of a simplicial map in a laminated context.
###### Definition 4.1.2.
Let $`\mathrm{\Lambda }`$ be an essential lamination of $`M`$. A map $`\sigma :\mathrm{\Delta }^1M`$ is transverse if there is no back–tracking; i.e. there is no subinterval of $`\mathrm{\Delta }^1`$ whose image can be homotoped relative to its endpoints into a leaf of $`\mathrm{\Lambda }`$. A map $`\sigma :\mathrm{\Delta }^iM`$ with $`i2`$ is transverse if the induced lamination $`\sigma ^1(\mathrm{\Lambda })`$ of $`\mathrm{\Delta }^i`$ is non–singular and can be completed to an affine foliation of $`\mathrm{\Delta }^i`$
Denote by $`(A,\mathrm{\Lambda })_{LG}`$ the norm of a homology class $`A`$ with respect to $`\mathrm{\Lambda }`$.
For $`\sigma :\mathrm{\Delta }^3M`$ a transverse map and $`\mathrm{\Lambda }`$ nowhere dense, we can perturb $`\sigma `$ to be nondegenerate; that is, so that $`\sigma ^1(\mathrm{\Lambda })`$ is a collection of normal triangles and quadrilaterals compatible with a total ordering of the vertices of $`\mathrm{\Delta }^3`$. This can be done by wiggling $`\sigma `$ slightly so that no vertex is taken into a leaf of $`\mathrm{\Lambda }`$.
Let $`𝒯`$ be a triangulation of $`M`$, and let $`n(𝒯)`$ denote the number of tetrahedra in $`𝒯`$. Then any minimal genuine lamination $`\mathrm{\Lambda }`$ (i.e. one with every leaf dense in $`\mathrm{\Lambda }`$) can be put in normal form with respect to $`𝒯`$, by a theorem of M. Brittenham . For instance, an incompressible surface is an example of a minimal lamination. This establishes the following estimate
###### Theorem 4.1.3.
Let $`\mathrm{\Lambda }`$ be a minimal genuine lamination of $`M`$. Then
$$([M],\mathrm{\Lambda })_{LG}\underset{𝒯}{\mathrm{min}}4n(𝒯)$$
Proof: Let $`𝒯`$ be a triangulation of $`M`$. Then we can isotope $`\mathrm{\Lambda }`$ to be in normal form with respect to $`𝒯`$. Now we can subdivide $`𝒯`$, replacing each tetrahedron by $`4`$ tetrahedra, each of which only contains normal disks compatible with a total ordering on its vertices.
On the other hand, corollary 2.4.2 implies
###### Theorem 4.1.4.
For $`S`$ an incompressible surface in a hyperbolic $`3`$–manifold $`M`$, either $`S`$ is a fiber of a fibration of $`M`$ over $`S^1`$ or
$$([M],S)_{LG}>[M]_G$$
Proof: Either $`S`$ is a fiber of a fibration over $`S^1`$, or $`S`$ is quasigeodesic. In the second case, we can find three lifts $`\stackrel{~}{S}_1,\stackrel{~}{S}_2,\stackrel{~}{S}_3`$ of $`S`$ to $`\stackrel{~}{M}=^3`$ which bound disjoint half–spaces. There is an ideal triangle with one vertex in each of these half–spaces, and there is a regular ideal tetrahedron, one of whose faces is this triangle. It follows that any chain representing $`[M]`$ whose norm is sufficiently close to $`[M]_G`$ cannot be transverse to $`S`$.
### 4.2. Order trees
The following definition is found in .
###### Definition 4.2.1.
An order tree is a set $`T`$ together with a collection $`𝒮`$ of linearly ordered subsets called segments, each with distinct least and greatest elements called the initial and final ends. If $`\sigma `$ is a segment, $`\sigma `$ denotes the same subset with the reverse order, and is called the inverse of $`\sigma `$. The following conditions should be satisfied:
* $`\sigma 𝒮\sigma 𝒮`$
* Any closed subinterval of a segment is a segment (if it has more than one element).
* Any two elements of $`T`$ can be joined by a finite sequence of segments $`\sigma _i`$ with the final end of $`\sigma _i`$ equal to the initial end of $`\sigma _{i+1}`$.
* Given a cyclic word $`\sigma _0\sigma _1\mathrm{}\sigma _{k1}`$ (subscripts mod $`k`$) with the final end of $`\sigma _i`$ equal to the initial end of $`\sigma _{i+1}`$, there is a subdivision of the $`\sigma _i`$ yielding a cyclic word $`\rho _0\rho _1\mathrm{}\rho _{n1}`$ which becomes the trivial word when adjacent inverse segments are canceled.
* If $`\sigma _1`$ and $`\sigma _2`$ are segments whose intersection is a single element which is the final element of $`\sigma _1`$ and the initial element of $`\sigma _2`$ then $`\sigma _1\sigma _2`$ is a segment containing $`\sigma _1`$ and $`\sigma _2`$.
An order tree is topologized by the order topology on segments. We assume in the sequel that our order trees are $``$–order trees — that is, $`2`$nd countable order trees whose segments are order isomorphic to compact intervals of $``$.
Let $`\mathrm{\Gamma }=\pi _1(M)`$ act by automorphisms on an order tree $`T`$ and suppose that we have a $`\mathrm{\Gamma }`$–equivariant surjective map
$$\varphi :\stackrel{~}{M}T$$
A singular map $`\sigma :\mathrm{\Delta }^iM`$ is transverse if for any lift $`\stackrel{~}{\sigma }:\mathrm{\Delta }^i\stackrel{~}{M}`$ the composition $`\varphi \stackrel{~}{\sigma }:\mathrm{\Delta }^iT`$ maps $`\mathrm{\Delta }^i`$ to a totally ordered segment of $`T`$.
Say that a sequence of such maps $`\varphi _i:\stackrel{~}{M}T_i`$ converges to $`\varphi :\stackrel{~}{M}T`$ if every map $`\sigma :\mathrm{\Delta }M`$ transverse with respect to $`\varphi `$ is eventually transverse with respect to $`\varphi _i`$, for sufficiently large $`i`$.
For a representation $`\rho :\mathrm{\Gamma }\text{Aut}(T)`$, an equivariant map $`\varphi :\stackrel{~}{M}T`$ and a homology class $`\mu H_i(M;)`$ we can define a norm $`(\mu ,\varphi )_{FG}`$ as before as the $`L_1`$ norm on the singular chains representing $`\mu `$ which are transverse with respect to $`\varphi `$. Observe that this norm does not really depend on the map $`\varphi `$, since it is determined up to equivariant homotopy by $`\rho `$, and thus this is really a norm on $`H_i(\mathrm{\Gamma };)`$ depending only on $`\rho `$.
###### Theorem 4.2.2.
Let $`\mathrm{\Gamma }=\pi _1(M)`$ and let $`\varphi _i:\stackrel{~}{M}T_i`$ be a sequence of equivariant maps for actions $`\rho _i:\mathrm{\Gamma }\text{Aut}(T)`$. Suppose this sequence converges to $`\varphi :\stackrel{~}{M}T`$ equivariant for some action $`\rho :\mathrm{\Gamma }\text{Aut}(T)`$.
Then given $`\mu H_i(M;)`$ we have the inequality
$$(\mu ,\varphi )_{FG}limsup(\mu ,\varphi _j)_{FG}$$
Proof: Any geometric chain in $`M`$ representing $`\mu `$ which is transverse for $`\varphi `$ will be transverse for $`\varphi _i`$ for sufficiently large $`i`$.
Examples of group actions on order trees arise in the study of essential laminations, where the lamination in the universal cover is dual to an order tree which can be taken to be an $``$–order tree by replacing isolated leaves by foliated $`I`$–bundles over those leaves. We have already seen from the example of foliations that this norm is not trivial and can vary for different representations of a fixed group.
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# Hard synchrotron BL Lacs: the case of 1ES 1101-232
## 1. Introduction
Overall spectral energy distributions (SED) of BL Lacs and blazars in general show two broad peaks: the synchrotron one at low energies and the inverse Compton scattering peak at high energies. The position of the synchrotron peak defines different classes of BL Lacs: the HBL (High-peaked BL Lacs) and the LBL (Low-peaked BL Lacs). Ghisellini et al. (1998) and Fossati et al. (1998) propose a sequence for blazars in which the energy of the peak is anti-correlated with the bolometric luminosity, and fainter objects, as HBL, should have a peak in the UV–X-ray band.
1ES 1101–232 (z=0.186) is an extreme case of HBL in which the synchrotron component peaks in the X-ray band ($``$ 1 keV), as shown by our previous observation (Wolter et al. 1998). Even if not as extreme as that of the flaring states of Mkn 501 (Pian et al. 1998) and 1ES 2344+514 (Catanese et al. 1998), the SED of 1ES 1101–232 makes it a good candidate for TeV emission.
## 2. X-ray data
BeppoSAX has observed 1ES1101–232 on two occasions, on 4 Jan 97 and 19 Jun 98. A single power law fit with Galactic absorption at low energy is rejected for both observations, while a broken power law yields an acceptable $`\chi ^2`$. In Wolter et al. (1999) all the details of the fits are reported. The broken power law model is preferred, from a statistical point of view besides for physical reasons, even over a single power law with intrinsic absorption. The PDS observations, being so short, are not of sufficient statistical significance to put a real constraint on the spectrum.
The position of the break energy ($`E_0`$) and the slope of the low energy part of spectrum ($`\alpha _1`$) are the same in the two observations within the errors. On the contrary, the portion of the spectrum at higher energies (i.e. above $`E_0`$) has changed between the two observations. We therefore fit the two datasets together, by using an appropriate model; the best fit of a broken power law model, in which only the high energy index $`\alpha _2`$ is untied between the two observations, is acceptable (see Table 1).
The fluxes are consistent with those obtained by the separate fits. Only the intensity above 2 keV changed (of $`32\%`$) between the two observations. Even if the flux variation is small, this result might bear an impact on spectral variability models in BL Lacs.
## 3. Spectral Energy Distribution
Figure 1: SED, points from literature and BeppoSAX observations. See text for an explanation of the model. Light gray line and dots refer to the Jan 1997 observation, while dark line and dots to the June 1998 BeppoSAX observation.
By using the same data as reported in Wolter et al. (1998) and adding the second BeppoSAX observation we construct the SED of Figure 1 and 2. Furthermore, 1ES1101-232 has been observed on the nights of 19-27 May 1998 with the Durham University Mark 6 atmospheric $`\stackrel{ˇ}{\mathrm{C}}`$erenkov telescope (Chadwick et al. 1999). The source was not detected and an upper limit of $`f_{TeV}`$ ($`>300`$ GeV) = 3.7 $`\times 10^{11}`$ photons cm<sup>-2</sup> s<sup>-1</sup> has been derived from the observation. This value also has been plotted in Figure 1.
We can reproduce the observed SED by using the homogeneous Synchrotron–Self Compton model described in detail in Ghisellini et al. (1998). A power-law distribution of electrons with slope $`n`$ and minimum Lorentz factor $`\gamma _{min}`$ is continuously injected in a spherical region with radius $`R`$. The source is in relativistic motion toward the observer and relativistic effects are expressed by the Doppler factor $`\delta `$. Electrons are free to cool and form the low energy flat spectrum with spectral index $`\alpha =0.5`$.
Figure 2: SED, enlargement of the X-ray band with the two observations and model. Light gray line and dots refer to the Jan 1997 observation, while dark line and dots to the June 1998 BeppoSAX observation. The agreement between the model and the X-ray points in the two observations is evident.
The model over-imposed on the SED is derived assuming a radius of $`R=1\times 10^{16}`$ cm, $`\delta `$=15, $`L_{\mathrm{inj}}^{\mathrm{intr}.}=9.3\times 10^{41}`$ erg/s; $`\gamma _{max}=4\times 10^6`$, with no external photons. The slope of the injected electrons is $`s`$=2.7 (1998) or $`s`$=1.95 (1997). $`B`$ = 0.6 Gauss (and $`\gamma _{min}^{inj}=5.\times 10^4`$) for the continuous line; $`B`$ = 1.2 Gauss (and $`\gamma _{min}^{inj}=3.\times 10^4`$) for the dashed line.
## 4. Magnetic field
A small change in the magnetic field, while still consistent with the X-ray (BeppoSAX) observations (see Figure 1), produces a very different TeV emission. The TeV band data can therefore put stringent constraints on the magnetic field.
The TeV upper limit indicates that the Compton peak cannot be higher than the synchrotron peak ($`L_C/L_S1`$); using the analytical relations discussed in Tavecchio et al. (1998) we can calculate the minimum $`B`$ allowed by the observed TeV upper limit for different values of $`\nu _c`$ and $`\delta `$. The values that produce a SED in agreement with both the X-ray spectra and the TeV upper limit are very similar to those found for Mkn 501 (e.g. $`\delta =15`$ and $`B`$=0.2 G; Kataoka et al. 1999) implying that the physical conditions of the two sources are also quite similar.
## 5. Conclusions
The X-ray spectrum of 1ES 1101-232 is fitted by a broken power law (a single or an absorbed power law are not statistically acceptable) with a break at 1.3 - 1.9 keV. From the first to the second observation, the spectrum varied at high energies, becoming softer (steeper). The flux has therefore decreased by about 32%, in the 2-10 keV band.
The TeV observation has not yielded a detection. However, since the TeV emission is largely sensitive to parameters like the magnetic field that produces the Synchrotron emission, interesting limits can be put on this quantity. Of course, more sensitive TeV instruments will produce more stringent constraints on the higher energy part of the spectrum and therefore on the emission mechanisms.
Multifrequency, simultaneous observations (e.g. optical, X-ray, TeV) will thus allow us to explain the variability of the sources, both from the energetic and the spectral distribution point of view.
## ACKNOWLEDGEMENTS
This work has received partial financial support from the Italian Space Agency and from the European Commission, TMR Programme, Research Network Contract ERBFMRXCT96-0034 “CERES”
## REFERENCES
Catanese M. et al., 1998, ApJ, 501, 616.
Chadwick P.M. et al., 1999, ApJ, 513, 161.
Fossati G. et al., 1998, MNRAS, 299, 433.
Ghisellini G. et al., 1998, MNRAS, 301, 451.
Kataoka, J., et al. 1999, ApJ, 514, 138.
Pian E. et al., 1998, ApJL, 492, L17.
Tavecchio F. et al. 1998, ApJ, 509, 608.
Wolter A. et al. 1998, A&A, 335, 899.
Wolter A. et al. 1999, A&A, submitted.
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# The Superfield Formalism Applied to the Noncommutative Wess-Zumino Model
## 1 Introduction
There are no doubts that the concept of space-time as a differentiable manifold cannot be extrapolated to extremely short distances . Simple heuristic arguments combining the principles of both general relativity and quantum theory imply that it is impossible to locate a particle with an arbitrarily small uncertainty . This means that standard differential geometry is certainly not an adapted method for physics at short distances. On the other hand, our standard description of fundamental interactions is exclusively based upon standard differential geometry – fibre bundles for the standard model and Riemannian geometry for gravity.
If standard differential geometry is not appropriate – what else should replace it? A promising candidate is *noncommutative geometry* pioneered by Alain Connes , see for a recent review. An excellent book on this subject will soon appear. Noncommutative geometry is the attempt to extend the principles of quantum mechanics to geometry itself: The use of operator algebras, Hilbert spaces, functional analysis. Within this framework the analogue of gauge theory has been developed which reduces to Yang-Mills theory if the geometry is commutative. The first example of such a theory on a noncommutative space appeared almost ten years ago, when Connes and Rieffel developed classical two dimensional Yang-Mills theory on the noncommutative torus . There is also an interesting class of “almost commutative” geometries which allow to treat Yang-Mills and Higgs fields on an equal footing and lead to a new understanding of spontaneous symmetry breaking, see e.g. . The analogue of external field quantization on noncommutative spaces was proposed in .
Thus, the strategy of noncommutative geometry is to generalize the mathematical structures encountered in experimentally confirmed physics. Another approach to short-distance physics is string theory/ M theory which tries to guess physics from new first principles. At first sight it seems unlikely that noncommutative geometry and string theory could be related. However, it has been shown that certain noncommutative geometries arise as limiting cases of string theory. The first hint came from where compactifications of M theory on the noncommutative torus were introduced, leading to the interpretation that the step from the commutative to the noncommutative torus corresponds to turning on a constant background 3-form $`C`$. Then, in it was shown that this situation is obtained by starting with a type IIa superstring theory with non-zero Neveu-Schwarz $`B`$ field and taking a scaling limit according to . That idea was thoroughly investigated in . Using the results of about instantons on noncommutative $`^4`$, Seiberg and Witten argued that there is an equivalence between the Yang-Mills theories on standard $`^4`$ and $`_{nc}^4`$.
It should be mentioned that matrix theories were studied long before M theory was proposed, and that these matrix theories did contain certain noncommutative features. These models live on a lattice, and the number of degrees of freedom is reduced when the size $`N`$ of the matrix goes to infinity . Putting them on a torus instead of a lattice, twisted boundary conditions are possible. Then the action can be rewritten in terms of noncommuting matrix derivatives $`[\mathrm{\Gamma }^{(j)},.]`$, with $`[\mathrm{\Gamma }^{(2j)},\mathrm{\Gamma }^{(2j+1)}]=2\pi \mathrm{i}/N`$, see .
The Seiberg-Witten paper inspired numerous attempts to formulate quantum field theories on noncommutative geometries. Nevertheless, quantum field theory on noncommutative spaces is also interesting in its own right. As standard quantum field theory is the art to deal with problems of interactions at short distances, see e.g. the proceedings , one should expect interesting features when doing quantum field theory on spaces with different short-distance structure. Singularities in standard quantum field theories are a consequence of the point-like interactions. There has been some hope that smearing out the points it is possible to avoid these ultraviolet divergences. Example of geometries where points are replaced by some sort of cells are the fuzzy spaces, see e.g. . Such fuzzy spaces also arise as limits of brane dynamics .
That divergences are not avoided on $`_{nc}^4`$ was first noticed by Filk . He showed that the noncommutative model contains Feynman graphs which are identical with their commutative counterparts. The $`_{nc}^4`$ is defined by the following commutator of the coordinate operators $`\{q^\mu \}`$:
$$[q^\mu ,q^\nu ]=\mathrm{i}\mathrm{\Sigma }^{\mu \nu },\mathrm{\Sigma }^{\mu \nu }=\mathrm{\Sigma }^{\nu \mu }.$$
Integrals corresponding to Feynman graphs in noncommutative QFTs differ from their commutative counterparts by phase factors $`\mathrm{e}^{\frac{\mathrm{i}}{2}\mathrm{\Sigma }^{\mu \nu }p_\mu k_\nu }`$, where $`p,k`$ are internal or external momenta. The case $`p=k`$ is possible, and in this situation the integrals of the commutative and the noncommutative theory coincide.
This raised the question whether the noncommutative QFT is renormalizable. On the one-loop level this was affirmed for Yang-Mills theory on $`_{nc}^4`$ and the noncommutative 4-torus as well as for supersymmetric Yang-Mills theory in $`(2+1)`$ dimensions, with space being the noncommutative 2-torus . Quantum electrodynamics on $`_{nc}^4`$ was treated in and the BF-Yang-Mills theory in .
These results lead to the hope that Yang-Mills theory on $`_{nc}^4`$ is renormalizable to all orders in perturbation theory. It was shown however by Minwalla, Van Raamsdonk and Seiberg that at least for scalar theories ($`\varphi ^4`$ on $`_{nc}^4`$ and $`\varphi ^3`$ on $`_{nc}^6`$) there is a new type of infrared divergences which ruins the perturbative renormalization beyond one loop. This follows immediately from the work of Filk , it was nevertheless completely unexpected: The oscillatory factors $`\mathrm{e}^{\frac{\mathrm{i}}{2}\mathrm{\Sigma }^{\mu \nu }p_\mu k_\nu }`$ render in four dimensions an otherwise (superficially) divergent integral convergent if both $`p,k`$ are internal momenta. If e.g. $`p`$ is external and $`k`$ is internal, the integral is convergent as well, but of course only as long as $`p0`$, where the ultraviolet divergence of the commutative theory reappears. This manifests as an infrared divergence coming from a ultraviolet-dangerous integration (UV/IR mixing). It turns out that the power counting degree of the new superficial IR divergence of the noncommutative theory coincides with the degree of the superficial UV divergence of the commutative theory .
Inspite of some rumour in the literature that Yang-Mills theory has only logarithmic divergences, it turned out that Yang-Mills on $`_{nc}^4`$ has quadratic IR divergences for the gluon propagator and linear IR divergences for the 3-gluon vertex which prevent the perturbative renormalization. This result was also derived from the scaling limit of string theory . Graphs made exclusively of nested ghost propagator corrections have only logarithmic divergences and are renormalizable at any loop order . In it was also shown that adding fermions in the adjoint representation cancels the quadratic and linear IR divergences at one loop. This was a hint that supersymmetric Yang Mills could be renormalizable, which is also strongly supported by the divergence analysis in .
The UV/IR mixing was also observed in noncommutative complex scalar $`\varphi ^4`$ theory where the interaction potential $`a\varphi ^{}\varphi \varphi ^{}\varphi +b\varphi ^{}\varphi ^{}\varphi \varphi `$ is only one-loop renormalizable for $`b=0`$ or $`a=b`$. UV/IR mixing does also occur in non-relativistic models . Anomalies were investigated in and the operator product expansion in . Discrete symmetries (CPT) were investigated in .
A completely finite model is Chern-Simons theory on $`_{nc}^3`$ . Examples of models where the limit $`\mathrm{\Sigma }^{\mu \nu }0`$ is smooth can be found in . Yang-Mills with fermions on $`_{nc}^3`$ was studied in with respect to the induced Chern-Simons action. Noncommutative topological massive Yang-Mills was treated in . There are also interesting two dimensional noncommutative models such as the nonlinear $`\sigma `$ model and the noncommutative Wess-Zumino-Witten model .
Even at the tree-level a field theory on $`_{nc}^4`$ shows unusual features such as violation of causality and S-matrix unitarity if time is noncommutative. Experimental limits on $`\mathrm{\Sigma }^{\mu \nu }`$ are discussed in . Beside the renormalizability problem, a quadratic IR divergence in the one-loop gluon propagator leads to a confinement of the model to a region which size is of the order $`\sqrt{|\mathrm{\Sigma }^{\mu \nu }|}`$. Thus, such a model cannot have the standard model as its low energy limit. The standard model can only be extended to noncommutative space-time if the quadratic IR divergences cancel, which is probably the case of supersymmetric versions. This motivates the interest in supersymmetry on noncommutative spaces, apart from the scaling limit of superstring theory.
Concerning supersymmetry, in a deformation also of the anticommutator of the fermionic superspace coordinates $`\theta `$ was considered, it was shown however that such a deformation is not compatible with supertranslations and chiral superfields. This result was also derived in via the scaling limit of string theory. A superspace formulation (at the classical level) of the Wess-Zumino model and super Yang-Mills was given in . Employing the component formulation it was eventually proved in that the Wess-Zumino (WZ) model on $`_{nc}^4`$ is renormalizable to all orders in perturbation theory. One-loop renormalizability of $`N=2`$ super Yang-Mills was obtained in .
In this paper we extend the work of Filk to the superfield formalism and apply the techniques to the noncommutative Wess-Zumino model. The paper is organized as follows: In section 2, following closer the work of Filk than , we introduce the notion of a superoperator on noncommutative space-time. In section 3 the Wess-Zumino model in the superfield formalism is introduced at the classical level, and in section 4 we postulate a noncommutative version of the Gell-Mann Low formula, and apply it to compute the one loop self-energy contributions. To be complete, we deduce the superfield Feynman rules in section 5, and apply them to the one loop vertex corrections. In section 6 we discuss the renormalizability to all loop orders.
## 2 Superfields with noncommutative coordinates
Following we consider the space-time coordinates of a flat space as self-adjoint operators in a Hilbert space with the following algebra
$$[q^\mu ,q^\nu ]=\mathrm{i}\mathrm{\Sigma }^{\mu \nu },[\mathrm{\Sigma }^{\mu \nu },q^\rho ]=0,$$
(1)
where $`\mathrm{\Sigma }^{\mu \nu }`$ is real and antisymmetric. In order to describe the superfields consistently within the framework of Filk one defines the operator
$$T(k)=e^{\mathrm{i}k_\mu q^\mu },$$
(2)
with the properties
$`T^+(k)`$ $`=`$ $`T(k),`$ (3)
$`T(k)T(k^{})`$ $`=`$ $`T(k+k^{})e^{\mathrm{i}k\times k^{}},k\times k^{}={\displaystyle \frac{1}{2}}\mathrm{\Sigma }^{\mu \nu }k_\mu k_\nu ^{},`$ (4)
and formally
$`\text{tr}\left[T(k)\right]`$ $`=`$ $`{\displaystyle d^4qe^{\mathrm{i}k_\mu q^\mu }}={\displaystyle \underset{\mu }{}}\delta (k_\mu )(2\pi )^4.`$ (5)
Additionally, we are dealing with classical chiral and anti-chiral superfields $`\mathrm{\Phi }(x,\theta ,\overline{\theta })`$ and $`\overline{\mathrm{\Phi }}(x,\theta ,\overline{\theta })`$ in the real representation , defined on an ordinary manifold by
$`\mathrm{\Phi }(x,\theta ,\overline{\theta })`$ $`=`$ $`\mathrm{\Phi }_1(x\mathrm{i}\theta \sigma \overline{\theta },\theta )=e^{\mathrm{i}\theta \sigma \overline{\theta }}\mathrm{\Phi }_1(x,\theta ),`$ (6)
$`\overline{\mathrm{\Phi }}(x,\theta ,\overline{\theta })`$ $`=`$ $`\overline{\mathrm{\Phi }}_2(x+\mathrm{i}\theta \sigma \overline{\theta },\overline{\theta })=e^{\mathrm{i}\theta \sigma \overline{\theta }}\overline{\mathrm{\Phi }}_2(x,\overline{\theta }),`$ (7)
where
$`\mathrm{\Phi }_1(x,\theta )`$ $`=`$ $`A(x)+\theta ^\alpha \psi _\alpha (x)+\theta ^\alpha \theta _\alpha F(x),`$ (8)
$`\overline{\mathrm{\Phi }}_2(x,\overline{\theta })`$ $`=`$ $`\overline{A}(x)+\overline{\theta }_{\dot{\alpha }}\overline{\psi }^{\dot{\alpha }}(x)+\overline{\theta }_{\dot{\alpha }}\overline{\theta }^{\dot{\alpha }}\overline{F}(x).`$ (9)
The corresponding covariant derivatives and further technical material on the superfield formalism is collected in appendix A. To a chiral classical superfield $`\mathrm{\Phi }_1(x_1,\theta _1)\mathrm{\Phi }_1(1)`$ one defines the Fourier-transform as<sup>1</sup><sup>1</sup>1We use the notation $`dx=d^4x`$, $`dk=\frac{d^4k}{(2\pi )^4}`$
$$\mathrm{\Phi }_1(1)=𝑑pe^{\mathrm{i}px_1}\stackrel{~}{\mathrm{\Phi }}_1(p,\theta _1),$$
(10)
where the coefficients of $`\stackrel{~}{\mathrm{\Phi }}_1(p,\theta _1)`$ belong to $`𝒮\left(^4\right)`$. The “inverse” reads
$`\stackrel{~}{\mathrm{\Phi }}_1(p,\theta _1)`$ $`=`$ $`{\displaystyle 𝑑S_2e^{\mathrm{i}px_2}\delta (\theta _{12})\mathrm{\Phi }_1(x_2,\theta _2)}`$ (11)
$`=`$ $`{\displaystyle 𝑑x_2e^{\mathrm{i}px_2}\mathrm{\Phi }_1(x_2,\theta _1)},`$
where $`dS_2`$ and $`\delta (\theta _{12})`$ are given in appendix A. Corresponding to Filk , one associates to the classical chiral superfield $`\mathrm{\Phi }_1(x_1,\theta _1)`$ the following superoperator
$`\mathrm{\Phi }_1(q_1,\theta _1)`$ $`=`$ $`{\displaystyle 𝑑S_2𝑑kT(k)e^{\mathrm{i}kx_2}\delta (\theta _{12})\mathrm{\Phi }_1(x_2,\theta _2)}`$ (12)
$`=`$ $`{\displaystyle 𝑑x_2𝑑kT(k)e^{\mathrm{i}kx_2}\mathrm{\Phi }_1(x_2,\theta _1)}`$
$`=`$ $`{\displaystyle 𝑑kT(k)\stackrel{~}{\mathrm{\Phi }}_1(k,\theta _1)}.`$
The trace operation allows to recover the classical chiral superfield $`\mathrm{\Phi }_1(x_1,\theta _1)`$
$$\mathrm{\Phi }_1(1)=𝑑ke^{\mathrm{i}kx_1}\text{tr}\left[\mathrm{\Phi }_1(q_1,\theta _1)T^+(k)\right].$$
(13)
Following Filk’s idea we are now able to construct a $``$-product, the so-called Moyal product, of two classical superfields
$$(\mathrm{\Phi }_1\mathrm{\Phi }_1)(1)=𝑑ke^{\mathrm{i}kx_1}\text{tr}\left[\mathrm{\Phi }_1(q_1,\theta _1)\mathrm{\Phi }_1(q_1,\theta _1)T^+(k)\right].$$
(14)
With the relations (3)-(5) eq.(14) becomes
$$(\mathrm{\Phi }_1\mathrm{\Phi }_1)(1)=𝑑k_1𝑑k_2e^{\mathrm{i}\left(k_1+k_2\right)x_1}e^{\mathrm{i}k_1\times k_2}\stackrel{~}{\mathrm{\Phi }}_1(k_1,\theta _1)\stackrel{~}{\mathrm{\Phi }}_1(k_2,\theta _1),$$
(15)
and with (11) there follows
$`(\mathrm{\Phi }_1\mathrm{\Phi }_1)(1)`$ $`=`$ $`{\displaystyle 𝑑k_1𝑑k_2e^{\mathrm{i}(k_1+k_2)x_1}e^{\mathrm{i}k_1\times k_2}}`$ (16)
$`\times {\displaystyle }dS_1^{}{\displaystyle }dS_2^{}e^{\mathrm{i}k_1x_1^{}\mathrm{i}k_2x_2^{}}\delta (\theta _{11^{}})\delta (\theta _{12^{}})\mathrm{\Phi }_1(x_1^{},\theta _1^{})\mathrm{\Phi }_1(x_2^{},\theta _2^{}).`$
Additionally, one has also
$$(\mathrm{\Phi }_1\mathrm{\Phi }_1\mathrm{\Phi }_1)(1)=𝑑ke^{\mathrm{i}kx_1}\text{tr}\left[\mathrm{\Phi }_1(q_1,\theta _1)\mathrm{\Phi }_1(q_1,\theta _1)\mathrm{\Phi }_1(q_1,\theta _1)T^+(k)\right].$$
(17)
Repeating the same steps as before, eq.(17) may be rewritten as
$`(\mathrm{\Phi }_1\mathrm{\Phi }_1\mathrm{\Phi }_1)(1)`$ $`=`$ $`{\displaystyle 𝑑k_1𝑑k_2𝑑k_3e^{\mathrm{i}(k_1+k_2+k_3)x_1}e^{\mathrm{i}_{i<j}^3k_i\times k_j}\stackrel{~}{\mathrm{\Phi }}_1(k_1,\theta _1)\stackrel{~}{\mathrm{\Phi }}_1(k_2,\theta _1)\stackrel{~}{\mathrm{\Phi }}_1(k_3,\theta _1)}`$ (18)
$`=`$ $`{\displaystyle 𝑑k_1𝑑k_2𝑑k_3e^{\mathrm{i}(k_1+k_2+k_3)x_1}e^{\mathrm{i}_{i<j}^3k_i\times k_j}}`$
$`\times {\displaystyle }dS_1^{}{\displaystyle }dS_2^{}{\displaystyle }dS_3^{}e^{\mathrm{i}k_1x_1^{}\mathrm{i}k_2x_2^{}\mathrm{i}k_3x_3^{}}`$
$`\times \delta (\theta _{11^{}})\delta (\theta _{12^{}})\delta (\theta _{13^{}})\mathrm{\Phi }_1(x_1^{},\theta _1^{})\mathrm{\Phi }_1(x_2^{},\theta _2^{})\mathrm{\Phi }_1(x_3^{},\theta _3^{}).`$
The $`\mathrm{\Phi }_1^3`$ is required to describe the corresponding interactions of the WZ-model. With the functional derivative for chiral superfields
$$\frac{\delta \mathrm{\Phi }_1(1)}{\delta \mathrm{\Phi }_1(2)}=\delta _S(1,2)=\delta (\theta _{12})\delta (x_1x_2)=\frac{1}{4}\theta _{12}^2\delta (x_1x_2),$$
(19)
one checks that the above definitions imply also
$$\frac{\delta }{\delta \mathrm{\Phi }_1(2)}𝑑S_1\left(\mathrm{\Phi }_1\mathrm{\Phi }_1\mathrm{\Phi }_1\right)(1)=3(\mathrm{\Phi }_1\mathrm{\Phi }_1)(2).$$
(20)
Finally, from (15) we get the useful relation
$$𝑑S_1\left(\mathrm{\Phi }_1\mathrm{\Phi }_1\right)(1)=𝑑S_1\mathrm{\Phi }_1(1)\mathrm{\Phi }_1(1).$$
(21)
## 3 Superfield formulation of the noncommutative Wess-Zumino model at the tree level
In four dimensional Minkowskian space-time the Wess-Zumino model in terms of superfields is defined at the tree level by the following action
$`\mathrm{\Gamma }^{(0)}`$ $`=`$ $`\mathrm{\Gamma }_{kin}^{(0)}+\mathrm{\Gamma }_m^{(0)}+\mathrm{\Gamma }_I^{(0)}`$ (22)
$`=`$ $`{\displaystyle \frac{1}{16}}{\displaystyle 𝑑V\overline{\mathrm{\Phi }}\mathrm{\Phi }}+{\displaystyle \frac{m}{8}}\left[{\displaystyle 𝑑S\mathrm{\Phi }^2}+{\displaystyle 𝑑\overline{S}\overline{\mathrm{\Phi }}^2}\right]+{\displaystyle \frac{g}{48}}\left[{\displaystyle 𝑑S\mathrm{\Phi }^3}+{\displaystyle 𝑑\overline{S}\overline{\mathrm{\Phi }}^3}\right].`$
Since $`\mathrm{\Phi }`$ and $`\overline{\mathrm{\Phi }}`$ are chiral and anti-chiral superfields the kinetic part of $`\mathrm{\Gamma }^{(0)}`$ can be rewritten as
$$\mathrm{\Gamma }_{kin}^{(0)}=\frac{1}{16}𝑑S\overline{D}^2\overline{\mathrm{\Phi }}(x,\theta ,\overline{\theta })\mathrm{\Phi }(x,\theta ,\overline{\theta })=\frac{1}{16}𝑑\overline{S}\overline{\mathrm{\Phi }}(x,\theta ,\overline{\theta })D^2\mathrm{\Phi }(x,\theta ,\overline{\theta }).$$
(23)
Carrying out the $`\theta `$ and $`\overline{\theta }`$ integration, which is in fact a differentiation, one gets always the last component of a superfield (products of superfields are again superfields) which furnishes the component formulation of supersymmetric field models. In order to derive the corresponding superfield propagators one uses the Legendre transformation between the functional for connected Green functions $`Z_c`$ and the vertex functional $`\mathrm{\Gamma }^{(0)}`$. Introducing external chiral and anti-chiral sources $`J`$ and $`\overline{J}`$ we have
$$Z_c[J,\overline{J}]=\mathrm{\Gamma }^{(0)}[\mathrm{\Phi },\overline{\mathrm{\Phi }}]+𝑑SJ\mathrm{\Phi }+𝑑\overline{S}\overline{J}\overline{\mathrm{\Phi }},$$
(24)
with
$$\frac{\delta \mathrm{\Gamma }^{(0)}}{\delta \mathrm{\Phi }}=J,\frac{\delta \mathrm{\Gamma }^{(0)}}{\delta \overline{\mathrm{\Phi }}}=\overline{J},$$
(25)
and
$$\frac{\delta Z_c}{\delta J}=\mathrm{\Phi },\frac{\delta Z_c}{\delta \overline{J}}=\overline{\mathrm{\Phi }}.$$
(26)
Solving eq.(25) for $`\mathrm{\Phi }=\mathrm{\Phi }[J,\overline{J}]`$ and $`\overline{\mathrm{\Phi }}=\overline{\mathrm{\Phi }}[J,\overline{J}]`$ one gets the desired superfield propagators
$`T\mathrm{\Phi }(1)\mathrm{\Phi }(2)_{(0)}`$ $`=`$ $`{\displaystyle \frac{\delta }{\mathrm{i}\delta J(1)}}\mathrm{\Phi }[J,\overline{J}](2)`$ (27)
$`=`$ $`{\displaystyle \frac{4\mathrm{i}m\delta _S(1,2)}{\mathrm{}+m^2}},`$
$`T\mathrm{\Phi }(1)\overline{\mathrm{\Phi }}(2)_{(0)}`$ $`=`$ $`{\displaystyle \frac{\delta }{\mathrm{i}\delta J(1)}}\overline{\mathrm{\Phi }}[J,\overline{J}](2)`$ (28)
$`=`$ $`{\displaystyle \frac{\mathrm{i}D_2^2\delta _S(1,2)}{\mathrm{}+m^2}}`$
$`=`$ $`{\displaystyle \frac{\mathrm{i}\overline{D}_1^2\delta _{\overline{S}}(1,2)}{\mathrm{}+m^2}},`$
and
$$T\overline{\mathrm{\Phi }}(1)\overline{\mathrm{\Phi }}(2)_{(0)}=\frac{4\mathrm{i}m\delta _{\overline{S}}(1,2)}{\mathrm{}+m^2}.$$
(29)
Having defined the WZ-model at the tree level with its superfield propagators we are now able to discuss radiative corrections at the one loop level with the help of the Gell-Mann Low formula in terms of superfields.
## 4 One loop self-energy corrections
The one loop calculations for the self-energy are governed by the Gell-Mann Low formula
$$G(1,\mathrm{},n)=T\mathrm{\Phi }(1)\mathrm{}\mathrm{\Phi }(n)=R\frac{T\mathrm{\Phi }(1)\mathrm{}\mathrm{\Phi }(n)e^{\mathrm{i}\mathrm{\Gamma }_I}_{(0)}}{Te^{\mathrm{i}\mathrm{\Gamma }_I}_{(0)}},$$
(30)
where we use for $`\mathrm{\Gamma }_I`$ the “deformed” interaction of the form
$$\mathrm{\Gamma }_I=\frac{g}{48}\left[𝑑S\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }+𝑑\overline{S}\overline{\mathrm{\Phi }}\overline{\mathrm{\Phi }}\overline{\mathrm{\Phi }}\right].$$
(31)
No attempt is made to prove formula (30). In our approach we define the model intuitively by (30) with a deformed interaction. In performing “modified” Wick contractions we will see that our procedure gives a meaningful result which is in agreement with a recent analysis in components .
The main advantage of our superfield procedure is the fact that one gets a very compact result in form of supergraphs. In order to demonstrate the power of the superfield formulation it is sufficient to discuss one representative. For this reason we calculate the one loop graph with one chiral and one anti-chiral external leg.
Up to some numerical factors the corresponding contribution is given by
$$G(1,\overline{2})T\mathrm{\Phi }(1)\overline{\mathrm{\Phi }}(2)\left[𝑑S_3\mathrm{\Phi }^3(3)+𝑑\overline{S}_3\overline{\mathrm{\Phi }}^3(3)\right]\left[𝑑S_4\mathrm{\Phi }^3(4)+𝑑\overline{S}_4\overline{\mathrm{\Phi }}^3(4)\right]_{(0)}.$$
(32)
By Wick contractions in the presence of deformed interactions we calculate the graph shown in fig.1, which is one of the four contributions corresponding to<sup>2</sup><sup>2</sup>2For further details in the commutative case see . (32). We find
$`G(1,\overline{2})_{3\overline{4}}`$ $``$ $`{\displaystyle 𝑑S_3\underset{i=1}{\overset{3}{}}dk_ie^{\mathrm{i}\left(k_1+k_2+k_3\right)x_3}e^{\mathrm{i}_{i<j}k_i\times k_j}𝑑x_{31}𝑑x_{32}𝑑x_{33}e^{\mathrm{i}k_1x_{31}\mathrm{i}k_2x_{32}\mathrm{i}k_3x_{33}}}`$ (33)
$`\times `$ $`{\displaystyle 𝑑\overline{S}_4\underset{i=1}{\overset{3}{}}dk_i^{}e^{\mathrm{i}\left(k_1^{}+k_2^{}+k_3^{}\right)x_4}e^{\mathrm{i}_{i<j}k_i^{}\times k_j^{}}𝑑x_{41}𝑑x_{42}𝑑x_{43}e^{\mathrm{i}k_1^{}x_{41}\mathrm{i}k_2^{}x_{42}\mathrm{i}k_3^{}x_{43}}}`$
$`\times `$ $`T\mathrm{\Phi }(1)\mathrm{\Phi }(3)_{(0)}\left(T\mathrm{\Phi }(3)\overline{\mathrm{\Phi }}(4)_{(0)}\right)^2T\overline{\mathrm{\Phi }}(4)\overline{\mathrm{\Phi }}(2)_{(0)}.`$
The required free superfield propagators are defined in eqs.(27)-(29). A straightforward but lengthy calculation leads to<sup>3</sup><sup>3</sup>3We omit the $`\mathrm{i}ϵ`$-prescription.
$$G(1,\overline{2})_{3\overline{4}}𝑑p_1e^{\mathrm{i}p_1(x_1x_2)}\frac{1}{\left(p_1^2m^2\right)^2}\mathrm{\Gamma }_{\mathrm{\Phi }\overline{\mathrm{\Phi }}}(p_1),$$
(34)
where $`\mathrm{\Gamma }_{\mathrm{\Phi }\overline{\mathrm{\Phi }}}`$ is the self-energy 1PI-vertex, in the one loop approximation (up to some numerical factors) given by
$`\mathrm{\Gamma }_{\mathrm{\Phi }\overline{\mathrm{\Phi }}}^{(1)}(p_1)`$ $``$ $`{\displaystyle 𝑑k\frac{1}{2}\left[1+\text{cos}(2p_1\times k)\right]\frac{D_2^2(p_1k)\stackrel{~}{\delta }_S(1,2)}{(p_1k)^2m^2}\frac{D_2^2(k)\stackrel{~}{\delta }_S(1,2)}{k^2m^2}}`$ (35)
$`=`$ $`{\displaystyle 𝑑k\text{cos}^2\left(p_1\times k\right)\frac{D_2^2(p_1k)\stackrel{~}{\delta }_S(1,2)}{(p_1k)^2m^2}\frac{D_2^2(k)\stackrel{~}{\delta }_S(1,2)}{k^2m^2}}.`$
Using
$$D_2^2(p)\stackrel{~}{\delta }_S(1,2)=\overline{D}_1^2(p)\stackrel{~}{\delta }_{\overline{S}}(1,2)=e^{E_{12}p},$$
(36)
one gets finally
$$\mathrm{\Gamma }_{\mathrm{\Phi }\overline{\mathrm{\Phi }}}^{(1)}(p_1)e^{E_{12}p_1}𝑑k\frac{1}{(p_1k)^2m^2}\frac{1}{k^2m^2}\text{cos}^2\left(p_1\times k\right),$$
(37)
where $`E_{12}`$ is defined by
$$E_{12}=\theta _1\sigma \overline{\theta }_1+\theta _2\sigma \overline{\theta }_22\theta _1\sigma \overline{\theta }_2.$$
(38)
The result in the form of eq.(37) shows in a very elegant manner that the total “undeformed” $`\theta `$-dependence is encoded in the exponent $`E_{12}`$, whereas the remaining Feynman integral represents the “component” result.
From (35) it is seen that the one loop self-energy corrections are just the sum of a usual planar contribution (the ’$`1`$’ term in $`[\mathrm{}]`$ of (35)) and a non-planar contribution (the cos(.) term in $`[\mathrm{}]`$). The planar contribution contains the expected logarithmically divergent wave function renormalization . Finally we must show that the non-planar integral
$$I(p,\stackrel{~}{p})=𝑑k\frac{e^{\mathrm{i}k\stackrel{~}{p}}}{\left(\left(pk\right)^2m^2+\mathrm{i}ϵ\right)\left(k^2m^2+\mathrm{i}ϵ\right)},$$
(39)
leads to a finite result for non-vanishing $`\stackrel{~}{p}^\mu =\mathrm{\Sigma }^{\mu \nu }p_\nu `$. The calculations are given in appendix B. We find
$`I(p,\stackrel{~}{p})`$ $`={\displaystyle \frac{2\mathrm{i}}{(4\pi )^2}}{\displaystyle _0^1}𝑑xK_0\left(\sqrt{(m^2p^2x(1x))(\stackrel{~}{p}^2)}\right).`$
## 5 Feynman rules in momentum space and one loop vertex corrections
In order to be complete this section is devoted to discuss the one loop vertex corrections directly with the help of the Feynman rules in momentum space. With the conventions of one has the following Feynman rules in momentum space
$`\text{}T\mathrm{\Phi }(1)\mathrm{\Phi }(2)_{(0)}^\stackrel{~}{}`$ $`=`$ $`{\displaystyle \frac{4\mathrm{i}m\stackrel{~}{\delta }_S(1,2)}{p^2m^2+\mathrm{i}ϵ}},`$ (40)
$`T\overline{\mathrm{\Phi }}(1)\overline{\mathrm{\Phi }}(2)_{(0)}^\stackrel{~}{}`$ $`=`$ $`{\displaystyle \frac{4\mathrm{i}m\stackrel{~}{\delta }_{\overline{S}}(1,2)}{p^2m^2+\mathrm{i}ϵ}},`$ (41)
$`T\mathrm{\Phi }(1)\overline{\mathrm{\Phi }}(2)_{(0)}^\stackrel{~}{}`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}D_2^2(p)\stackrel{~}{\delta }_S(1,2)}{p^2m^2+\mathrm{i}ϵ}}={\displaystyle \frac{\mathrm{i}\overline{D}_1^2(p)\stackrel{~}{\delta }_{\overline{S}}(1,2)}{p^2m^2+\mathrm{i}ϵ}},`$ (42)
$`T\overline{\mathrm{\Phi }}(1)\mathrm{\Phi }(2)_{(0)}^\stackrel{~}{}`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}\overline{D}_2^2(p)\stackrel{~}{\delta }_{\overline{S}}(1,2)}{p^2m^2+\mathrm{i}ϵ}}={\displaystyle \frac{\mathrm{i}D_1^2(p)\stackrel{~}{\delta }_S(1,2)}{p^2m^2+\mathrm{i}ϵ}},`$ (43)
$`{\displaystyle \frac{1}{8}}g(2\pi )^4\delta \left(p_1+p_2+p_3\right)\delta \left(\theta _{12}\right)\delta \left(\theta _{13}\right)\text{cos}\left(p_2\times p_3\right),`$ (45)
$`{\displaystyle \frac{1}{8}}g(2\pi )^4\delta \left(p_1+p_2+p_3\right)\delta \left(\overline{\theta }_{12}\right)\delta \left(\overline{\theta }_{13}\right)\text{cos}\left(p_2\times p_3\right).`$
With the Feynman rules (40)-(45) one easily confirms the result of eq.(35). Additionally, one can state the non-renormalization theorem for $`\mathrm{\Gamma }_{\mathrm{\Phi }\mathrm{\Phi }}`$ and $`\mathrm{\Gamma }_{\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }}`$ with only chiral (or anti-chiral) internal lines. A possible self-energy correction of this type is
$$\mathrm{\Gamma }_{\mathrm{\Phi }\mathrm{\Phi }}g^2𝑑k\frac{4\mathrm{i}m\stackrel{~}{\delta }_S(1,2)}{(p+k)^2m^2+\mathrm{i}ϵ}\frac{4\mathrm{i}m\stackrel{~}{\delta }_S(2,1)}{k^2m^2+\mathrm{i}ϵ}\text{cos}^2\left(p\times k\right).$$
(46)
Due to $`\stackrel{~}{\delta }_S(1,2)\stackrel{~}{\delta }_S(2,1)\frac{1}{16}\left(\theta _{12}^2\right)^2=0`$ one has: $`\mathrm{\Gamma }_{\mathrm{\Phi }\mathrm{\Phi }}=0`$. In a similar manner one can show that also the one loop vertex correction shown in fig.2a vanishes
$`\mathrm{\Gamma }_{\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }}`$ $``$ $`g^3{\displaystyle 𝑑k(4\mathrm{i}m)^3\frac{\stackrel{~}{\delta }_S(1,2)}{(p_1+k)^2m^2+\mathrm{i}ϵ}\frac{\stackrel{~}{\delta }_S(2,3)}{(p_1+p_2+k)^2m^2+\mathrm{i}ϵ}\frac{\stackrel{~}{\delta }_S(3,1)}{k^2m^2+\mathrm{i}ϵ}}`$ (47)
$`\times \text{cos}\left(p_1\times k\right)\text{cos}\left(p_2\times p_1+p_2\times k\right)\text{cos}\left(p_1\times k+p_2\times k\right).`$
Using the appendix A we find that the product of the three $`\stackrel{~}{\delta }`$-functions vanishes. However, there could exist a non-vanishing one loop correction if one allows mixed propagators, see fig.2b. Applying the above Feynman rules one gets
$`\mathrm{\Gamma }_{\mathrm{\Phi }\mathrm{\Phi }\overline{\mathrm{\Phi }}}`$ $``$ $`g^3{\displaystyle 𝑑k\frac{4\mathrm{i}m\stackrel{~}{\delta }_S(1,2)}{(p_1+k)^2m^2+\mathrm{i}ϵ}\frac{\mathrm{i}D_3^2(p_1+p_2+k)\stackrel{~}{\delta }_S(2,3)}{(p_1+p_2+k)^2m^2+\mathrm{i}ϵ}\frac{\mathrm{i}D_3^2(k)\stackrel{~}{\delta }_S(3,1)}{k^2m^2+\mathrm{i}ϵ}}`$ (48)
$`\times {\displaystyle \frac{1}{4}}[\text{cos}(p_1\times p_2)+\text{cos}(p_2\times p_1+2(p_1+p_2)\times k)+\text{cos}(p_2\times p_12p_1\times k)`$
$`+\text{cos}(p_2\times p_1+2p_2\times k)],`$
where we have separated the vertex correction in a planar (the $`\text{cos}\left(p_1\times p_2\right)`$ term, which does not depend on the internal momentum) and a non-planar contribution. Using (75) we conclude that the integral (48) is finite, since it goes asymptotically like $`\frac{1}{k^6}`$. This is required in order to secure stability of the classical action.
## 6 Renormalization at all orders and conclusion
Using the same integration techniques as in appendix B one can prove that the superficial integration encoded in any non-planar 1PI Feynman (sub)graph $`\gamma `$ produces a Bessel function $`K_.(.)`$ depending on the external momenta of $`\gamma `$. The Bessel function $`K_.(.)`$ tends exponentially to zero if the external momenta of $`\gamma `$ become large. This suffices to render the integration of a graph $`\gamma ^{}`$ containing $`\gamma `$ as a subgraph UV finite. The only problem could be an IR singularity of the Bessel function. However, since in the commutative Wess-Zumino model there are only logarithmic divergences (in the 2-point function), and because the only difference on the noncommutative space are phase factors, there can only be a logarithmic IR singularity coming from the Bessel function. Nested logarithmic singularities are IR-integrable, as it was explicitly demonstrated in . In conclusion, a graph in the WZ model which contains non-planar sectors leads always to a convergent integral. In particular, a non-planar graph in the standard sense of the commutative world is always convergent. Divergences come only from completely planar graphs, and they are subtracted e.g. by the BPHZ procedure as in . We therefore conclude that the Wess-Zumino model on noncommutative Minkowski space is renormalizable to all loop orders, a result which was already obtained in and conjectured in . Note that the $`\beta `$ functions of the noncommutative and the commutative theory differ because the standard non-planar graphs become finite on the noncommutative space.
In this paper we have demonstrated the strength of the superfield formalism. Especially we would like to emphasize that the superfield formalism enables us to apply eq.(36), which lowers the degree of divergence (both IR and UV) by two. Furthermore is the number of graphs to be computed considerable lower than in the work of . We believe that this formalism will prove useful for further investigations, in particular for super Yang-Mills theories on noncommutative $`^4`$.
## Appendix A Conventions and useful formulae in superspace
Let us briefly summarize some of the conventions and rules concerning supersymmetry and superspace (most of the rules are taken from ).
### Metrics, index transport and scalar products
The metric tensor of Minkowski space is given by $`g_{\mu \nu }=\mathrm{diag}(1,1,1,1)`$ and we use the following spinor metric:
$$ϵ^{\alpha \beta }=\mathrm{i}\sigma ^2=ϵ^{\dot{\alpha }\dot{\beta }},$$
(49)
$$ϵ_{\alpha \beta }=\mathrm{i}\sigma ^2=ϵ_{\dot{\alpha }\dot{\beta }},$$
(50)
$$ϵ_{\alpha \beta }ϵ^{\beta \gamma }=\delta _\alpha ^\gamma ,$$
(51)
$$\theta \eta =\theta ^\alpha \eta _\alpha ,\overline{\theta }\overline{\eta }=\overline{\theta }_{\dot{\alpha }}\overline{\eta }^{\dot{\alpha }},$$
(52)
$$\theta _\alpha =ϵ_{\alpha \beta }\theta ^\beta ,\overline{\theta }^{\dot{\alpha }}=ϵ^{\dot{\alpha }\dot{\beta }}\overline{\theta }_{\dot{\beta }}.$$
(53)
### Pauli matrices
$$\sigma ^0=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\sigma ^1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma ^2=\left(\begin{array}{cc}0& \mathrm{i}\\ \mathrm{i}& 0\end{array}\right),\sigma ^3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),$$
(62)
$$\stackrel{}{\sigma }=(\sigma ^1,\sigma ^2,\sigma ^3),$$
$$\sigma ^\mu =(1_2,\stackrel{}{\sigma })=(\sigma ^\mu )_{\alpha \dot{\beta }},$$
$$\overline{\sigma }^\mu =(1_2,\stackrel{}{\sigma })=(\overline{\sigma }^\mu )_{\dot{\alpha }\beta }.$$
(63)
### Covariant derivatives
$$D_\alpha =_\alpha \mathrm{i}(\sigma ^\mu )_{\alpha \dot{\beta }}\overline{\theta }^{\dot{\beta }}_\mu ,\overline{D}_{\dot{\alpha }}=\overline{}_{\dot{\alpha }}+\mathrm{i}\theta ^\beta (\overline{\sigma }^\mu )_{\dot{\alpha }\beta }_\mu ,$$
(64)
$$\{D_\alpha ,\overline{D}_{\dot{\beta }}\}=2\mathrm{i}(\sigma ^\mu )_{\alpha \dot{\beta }}_\mu ,\{D_\alpha ,D_\beta \}=\{\overline{D}_{\dot{\alpha }},\overline{D}_{\dot{\beta }}\}=0.$$
(65)
### Integration
$`{\displaystyle 𝑑V\mathrm{\Phi }}`$ $`={\displaystyle d^4xD^2\overline{D}^2\mathrm{\Phi }}\text{for any superfield }\mathrm{\Phi },`$ (66)
$`{\displaystyle 𝑑S\mathrm{\Phi }}`$ $`={\displaystyle d^4xD^2\mathrm{\Phi }}\text{for a chiral superfield }\mathrm{\Phi }\text{ (i.e. }\overline{D}\mathrm{\Phi }=0\text{)},`$ (67)
$`{\displaystyle 𝑑\overline{S}\overline{\mathrm{\Phi }}}`$ $`={\displaystyle d^4x\overline{D}^2\overline{\mathrm{\Phi }}}\text{for an anti-chiral superfield }\overline{\mathrm{\Phi }}\text{ (i.e. }D\overline{\mathrm{\Phi }}=0\text{)}.`$ (68)
### Functional differentiation, delta-functions and representations
$`{\displaystyle \frac{\delta \mathrm{\Phi }(1)}{\delta \mathrm{\Phi }(2)}}`$ $`=\delta _V(1,2)\text{for any superfield},`$ (69)
$`{\displaystyle \frac{\delta \mathrm{\Phi }_1(1)}{\delta \mathrm{\Phi }_1(2)}}`$ $`=\delta _S(1,2)\text{for a chiral superfield},`$ (70)
$`{\displaystyle \frac{\delta \overline{\mathrm{\Phi }}_2(1)}{\delta \overline{\mathrm{\Phi }}_2(2)}}`$ $`=\delta _{\overline{S}}(1,2)\text{for an anti-chiral superfield},`$ (71)
the numbers denoting points in superspace (for instance, $`(1)`$ is a shorthand-notation for $`((x^\mu )_1,(\theta _\alpha )_1,(\overline{\theta }^{\dot{\alpha }})_1)`$). The above delta-functions are in position space given by:
$`\delta _V(1,2)`$ $`={\displaystyle \frac{1}{16}}\theta _{12}^2\overline{\theta }_{12}^2\delta ^4(x_1x_2),`$ $`{\displaystyle 𝑑V_1\mathrm{\Phi }(1)\delta _V(1,2)}`$ $`=\mathrm{\Phi }(2),`$ (72)
$`\delta _S(1,2)`$ $`={\displaystyle \frac{1}{4}}\theta _{12}^2\delta ^4(x_1x_2),`$ $`{\displaystyle 𝑑S_1\mathrm{\Phi }(1)\delta _S(1,2)}`$ $`=\mathrm{\Phi }(2),`$ (73)
$`\delta _{\overline{S}}(1,2)`$ $`={\displaystyle \frac{1}{4}}\overline{\theta }_{12}^2\delta ^4(x_1x_2),`$ $`{\displaystyle 𝑑\overline{S}_1\overline{\mathrm{\Phi }}(1)\delta _{\overline{S}}(1,2)}`$ $`=\overline{\mathrm{\Phi }}(2),`$ (74)
Note that the first equalities of eq.(73) and eq.(74) are only valid in the chiral and anti-chiral representation, respectively. The “Fourier-transforms” of the $`\delta `$-functions in the real representation are
$`\stackrel{~}{\delta }_S(1,2)`$ $`={\displaystyle \frac{1}{4}}\theta _{12}^2e^{(\theta _1\sigma \overline{\theta }_2\theta _2\sigma \overline{\theta }_1)p},`$ $`D_1^2(p)\stackrel{~}{\delta }_S(1,2)`$ $`=e^{E_{12}p},`$ (75)
$`\stackrel{~}{\delta }_{\overline{S}}(1,2)`$ $`={\displaystyle \frac{1}{4}}\overline{\theta }_{12}^2e^{(\theta _1\sigma \overline{\theta }_2\theta _2\sigma \overline{\theta }_1)p},`$ $`\overline{D}_1^2(p)\stackrel{~}{\delta }_{\overline{S}}(1,2)`$ $`=e^{E_{12}p},`$ (76)
where
$$E_{12}=\theta _1\sigma \overline{\theta }_1+\theta _2\sigma \overline{\theta }_22\theta _1\sigma \overline{\theta }_2,$$
(77)
$$\theta _{12}^2=(\theta _1\theta _2)^2,\overline{\theta }_{12}^2=(\overline{\theta }_1\overline{\theta }_2)^2.$$
(78)
Different representations (real and (anti)chiral) of (anti)chiral superfields are connected by the following relations:
$`\mathrm{\Phi }(x,\theta ,\overline{\theta })`$ $`=\mathrm{\Phi }_1(x\mathrm{i}\theta \sigma \overline{\theta },\theta )=e^{\mathrm{i}\theta \sigma \overline{\theta }}\mathrm{\Phi }_1(x,\theta ),`$ (79)
$`\overline{\mathrm{\Phi }}(x,\theta ,\overline{\theta })`$ $`=\overline{\mathrm{\Phi }}_2(x+\mathrm{i}\theta \sigma \overline{\theta },\overline{\theta })=e^{\mathrm{i}\theta \sigma \overline{\theta }}\overline{\mathrm{\Phi }}_2(x,\overline{\theta }),`$ (80)
$`1`$ denoting the chiral and $`2`$ the anti-chiral representation; also, (anti)chiral fields have a simplified $`\theta `$-expansion:
$`\mathrm{\Phi }_1(x,\theta )`$ $`=A(x)+\theta ^\alpha \psi _\alpha (x)+\theta ^\alpha \theta _\alpha F(x),`$ (81)
$`\overline{\mathrm{\Phi }}_2(x,\overline{\theta })`$ $`=\overline{A}(x)+\overline{\theta }_{\dot{\alpha }}\overline{\psi }^{\dot{\alpha }}(x)+\overline{\theta }_{\dot{\alpha }}\overline{\theta }^{\dot{\alpha }}\overline{F}(x),`$ (82)
and the covariant derivatives are given by:
$`(D_\alpha \mathrm{\Phi })_1`$ $`=\left(_\alpha 2\mathrm{i}(\sigma ^\mu \overline{\theta })_\alpha _\mu \right)\mathrm{\Phi }_1,`$ $`(\overline{D}_{\dot{\alpha }}\mathrm{\Phi })_1`$ $`=\overline{}_{\dot{\alpha }}\mathrm{\Phi }_1,`$ (83)
$`(D_\alpha \mathrm{\Phi })_2`$ $`=_\alpha \mathrm{\Phi }_2,`$ $`(\overline{D}_{\dot{\alpha }}\mathrm{\Phi })_2`$ $`=\left(\overline{}_{\dot{\alpha }}+2\mathrm{i}(\theta \sigma ^\mu )_{\dot{\alpha }}_\mu \right)\mathrm{\Phi }_2.`$ (84)
## Appendix B Calculation of the non-planar self-energy graph integral
We are going to compute
$$I(p,\stackrel{~}{p})=\underset{ϵ0}{lim}\frac{d^4k}{(2\pi )^4}\frac{e^{\mathrm{i}(\stackrel{~}{p}_0k_0\stackrel{}{\stackrel{~}{p}}\stackrel{}{k})}}{((k_0p_0)^2|\stackrel{}{k}\stackrel{}{p}|^2m^2+\mathrm{i}ϵ)(k_0^2|\stackrel{}{k}|^2m^2+\mathrm{i}ϵ)}.$$
We apply Zimmermann’s $`ϵ`$ trick :
$`I(p,\stackrel{~}{p})`$ $`=\underset{ϵ0}{lim}{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{e^{\mathrm{i}(\stackrel{~}{p}_0k_0\stackrel{}{\stackrel{~}{p}}\stackrel{}{k})}}{((k_0p_0)^2(|\stackrel{}{k}\stackrel{}{p}|^2+m^2)(1\mathrm{i}ϵ))(k_0^2(|\stackrel{}{k}|^2+m^2)(1\mathrm{i}ϵ))}}`$
$`=\underset{ϵ0}{lim}{\displaystyle \frac{d^4k}{(2\pi )^4}_0^1𝑑x\frac{e^{\mathrm{i}(\stackrel{~}{p}_0k_0\stackrel{}{\stackrel{~}{p}}\stackrel{}{k})}}{((k_0^22p_0k_0x+p_0^2x)+(|\stackrel{}{k}|^22\stackrel{}{p}\stackrel{}{k}x+|\stackrel{}{p}|^2x+m^2)(\mathrm{i}ϵ1))^2}}`$
$`=\underset{ϵ0}{lim}{\displaystyle \frac{d^4k}{(2\pi )^4}_0^1𝑑x_0^{\mathrm{}}𝑑\alpha \alpha (ϵ^{}\mathrm{i})^2}`$
$`\times e^{\alpha (ϵ^{}\mathrm{i})((k_0^22p_0k_0x+p_0^2x)+(|\stackrel{}{k}|^22\stackrel{}{p}\stackrel{}{k}x+|\stackrel{}{p}|^2x+m^2)(\mathrm{i}ϵ1))+\mathrm{i}(\stackrel{~}{p}_0k_0\stackrel{}{\stackrel{~}{p}}\stackrel{}{k})}`$
$`=\underset{ϵ0}{lim}{\displaystyle \frac{d^4k}{(2\pi )^4}_0^1𝑑x_0^{\mathrm{}}𝑑\alpha \alpha (ϵ^{}\mathrm{i})^2e^{\mathrm{i}(\stackrel{~}{p}_0p_0\stackrel{}{\stackrel{~}{p}}\stackrel{}{p})x}}`$
$`\times e^{\alpha (ϵ^{}\mathrm{i})(k_0p_0x\frac{\mathrm{i}\stackrel{~}{p}_0}{2\alpha (ϵ^{}\mathrm{i})})^2\alpha (ϵϵ^{}+\mathrm{i}+ϵϵ^{}\mathrm{i})|\stackrel{}{k}\stackrel{}{p}x+\frac{\mathrm{i}\stackrel{}{\stackrel{~}{p}}}{2\alpha (ϵϵ^{}+\mathrm{i}+ϵϵ^{}\mathrm{i})}|^2}`$
$`\times e^{\alpha (ϵ^{}\mathrm{i})p_0^2x(1x)\alpha (ϵϵ^{}+\mathrm{i}+ϵϵ^{}\mathrm{i})(|\stackrel{}{p}|^2x(1x)+m^2)\frac{\stackrel{~}{p}_0^2}{4\alpha (ϵ^{}\mathrm{i})}\frac{|\stackrel{}{\stackrel{~}{p}}|^2}{4\alpha (ϵϵ^{}+\mathrm{i}+ϵϵ^{}\mathrm{i})}}.`$
For $`ϵ^{}<ϵ`$ we perform the Gaussian $`k`$ integration:
$`I(p,\stackrel{~}{p})`$ $`=\underset{ϵ0}{lim}{\displaystyle \frac{1}{(4\pi )^2}}{\displaystyle _0^1}𝑑x{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\alpha }{\alpha }}\left({\displaystyle \frac{ϵ^{}\mathrm{i}}{ϵϵ^{}+\mathrm{i}+ϵϵ^{}\mathrm{i}}}\right)^{3/2}\mathrm{e}^{\mathrm{i}(\stackrel{~}{p}_0p_0\stackrel{}{\stackrel{~}{p}}\stackrel{}{p})x}`$
$`\times e^{\alpha (ϵ^{}\mathrm{i})p_0^2x(1x)\alpha (ϵϵ^{}+\mathrm{i}+ϵϵ^{}\mathrm{i})(|\stackrel{}{p}|^2x(1x)+m^2)\frac{\stackrel{~}{p}_0^2}{4\alpha (ϵ^{}\mathrm{i})}\frac{|\stackrel{}{\stackrel{~}{p}}|^2}{4\alpha (ϵϵ^{}+\mathrm{i}+ϵϵ^{}\mathrm{i})}}`$
The $`\alpha `$ integration yields a result independent of $`ϵ^{}<ϵ`$:
$`I(p,\stackrel{~}{p})`$ $`=\underset{ϵ0}{lim}{\displaystyle \frac{2}{(4\pi )^2}}{\displaystyle _0^1}𝑑x\left({\displaystyle \frac{1\mathrm{i}ϵ}{1+ϵ^2}}\right)^{3/2}e^{\mathrm{i}(\stackrel{~}{p}_0p_0\stackrel{}{\stackrel{~}{p}}\stackrel{}{p})x}`$
$`\times K_0\left(\sqrt{\left(p_0^2x(1x)(1\mathrm{i}ϵ)(|\stackrel{}{p}|^2x(1x)+m^2)\right)\left(\stackrel{~}{p}_0^2\frac{1+\mathrm{i}ϵ}{1+ϵ^2}|\stackrel{}{\stackrel{~}{p}}|^2\right)}\right).`$
Now we can put $`ϵ0`$ and switch to Minkowskian scalar products:
$`I(p,\stackrel{~}{p})`$ $`={\displaystyle \frac{2\mathrm{i}}{(4\pi )^2}}{\displaystyle _0^1}𝑑xK_0\left(\sqrt{(m^2p^2x(1x))(\stackrel{~}{p}^2)}\right).`$ (85)
On the mass shell we have $`p^2=m^2`$ and $`\stackrel{~}{p}^20`$: If the particle moves for instance in 3-direction then $`\stackrel{~}{p}^2=\mathrm{\Sigma }_{03}^2m^2_{i=1,2}(\mathrm{\Sigma }_{i3}|\stackrel{}{p}|+\mathrm{\Sigma }_{i0}\sqrt{|\stackrel{}{p}|^2+m^2})^2`$.
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# Unbraiding the braided tensor product
## 1 Introduction and main theorem
As is well known, given two associative unital algebras $`𝒜_1,𝒜_2`$ (over the field $``$, say), one can build a new module algebra $`𝒜`$ which is as a vector space the tensor product $`𝒜=𝒜_1𝒜_2`$ of the two vector spaces (over the same field) by postulating the product law
$$(a_1a_2)(b_1b_2)=a_1b_1a_2b_2.$$
(1.1)
The resulting algebra is the ordinary tensor product algebra. (1.1) is equivalent to the set of relations
$`(a_1\mathrm{𝟏}_2)(b_1\mathrm{𝟏}_2)`$ $`=`$ $`a_1b_1\mathrm{𝟏}_2,`$ (1.2)
$`(a_1\mathrm{𝟏}_2)(\mathrm{𝟏}_1a_2)`$ $`=`$ $`a_1a_2,`$ (1.3)
$`(\mathrm{𝟏}_1a_2)(\mathrm{𝟏}_1b_2)`$ $`=`$ $`\mathrm{𝟏}_1a_2b_2,`$ (1.4)
$`(\mathrm{𝟏}_1a_2)(a_1\mathrm{𝟏}_2)`$ $`=`$ $`(a_1\mathrm{𝟏}_2)(\mathrm{𝟏}_1a_2).`$ (1.5)
However, in many cases the same goal can be reached also by replacing (1.5) by some suitable nontrivial commutation relations. With a standard abuse of notation we shall denote in the sequel $`a_1a_2`$ by $`a_1a_2`$ for any $`a_1𝒜_1`$, $`a_2𝒜_2`$ and omit all units $`\mathrm{𝟏}_i`$ when multiplied by non-unit elements; consequently (1.2-1.4) take trivial forms, whereas (1.5) becomes the commutation relation
$$a_2a_1=a_1a_2.$$
(1.6)
If $`𝒜_1,𝒜_2`$ are module algebras of a Lie algebra g , and we require $`𝒜`$ to be too, then (1.6) has no alternative, because any $`g\text{}`$ acts as a derivation on the (algebra as well as tensor) product of any two elements or, in Hopf algebra language, because the coproduct $`\mathrm{\Delta }(g)=g_{(1)}g_{(2)}`$ (at the rhs we have used Sweedler notation) of the Hopf algebra $`HU\text{}`$ is cocommutative. In the main part of this paper we shall work with right-module algebras (instead of left ones), and denote by $`:(a_i,g)𝒜_i\times Ha_ig𝒜_i`$ the right action; the reason is that they are equivalent to left comodule algebras, which are used in much of the literature. In section 5 we shall give the formulae for left module algebras. We recall that a right action $`:(a,g)𝒜\times Hag𝒜`$ by definition fulfills
$`a(gg^{})=(ag)g^{},`$ (1.7)
$`(aa^{})g=(ag_{(1)})(a^{}g_{(2)}).`$ (1.8)
If we “$`q`$-deform” this setting by taking as Hopf algebra $`H`$ the quantum group $`U_q\text{g}`$ , and as $`𝒜_i`$ the corresponding $`q`$-deformed module algebras, then it is also known that although $`\mathrm{\Delta }(g)`$ is no longer cocommutative, it is still possible to build the deformed counterpart of $`𝒜`$ if one replaces (1.6) with nontrivial commutation relations of the form
$$a_2a_1=(a_1^{(1)})(a_2^{(2)}).$$
(1.9)
Here $`^{(1)}^{(2)}H^+H^{}`$ denotes the socalled universal $`R`$-matrix of $`HU_q\text{g}\text{ }`$ , and $`H^\pm `$ denote the Hopf positive and negative Borel subalgebras of $`H`$. This yields instead of $`𝒜`$ a braided tensor product algebra $`𝒜^+=𝒜_1\underset{¯}{}^+𝒜_2`$ . An alternative one $`𝒜^{}=𝒜_1\underset{¯}{}^{}𝒜_2`$ is obtained by replacing in the previous formula $``$ by $`_{21}^1`$:
$$a_2a_1=(a_1^{1(2)})(a_2^{1(1)}).$$
(1.10)
Both $`𝒜^+`$ and $`𝒜^{}`$ go to the ordinary tensor product algebra $`𝒜`$ in the limit $`q1`$.
This is a particular example of a more general notion, that of a crossed (or twisted) tensor product of two unital associative algebras.
In view of (1.9) or (1.10) studying representations of $`𝒜^\pm `$ is a more difficult task than just studying the representations of $`𝒜_1,𝒜_2`$, taking their tensor products and studying the irreducible ones there contained. The degrees of freedom of $`𝒜_1,𝒜_2`$ are so to say “coupled”. One might ask whether one can “decouple” them by a transformation of generators.
In this work we present a sufficient condition for the construction of a transformation making $`𝒜^+`$ equal to an ordinary tensor product $`𝒜_1\stackrel{~}{𝒜}_2^+`$, with $`\stackrel{~}{𝒜}_2^+`$ a subalgebra of $`𝒜^+`$ isomorphic to $`𝒜_2`$ and commuting \[in the sense (1.6)\] with $`𝒜_1`$, although - of course - no longer a $`H`$-submodule; and similarly for $`𝒜^{}`$. In a quantum theory framework one could thus interpret the generators of $`𝒜_1,\stackrel{~}{𝒜}_2^\pm `$ as pertaining to decoupled degrees of freedom, describing e.g. some composite or “quasiparticle” excitations. Reducing $`𝒜^\pm `$ to a form $`𝒜_1\stackrel{~}{𝒜}_2^\pm `$ will be called an unbraiding of the braided tensor product algebra $`𝒜^\pm =𝒜_1\underset{¯}{}^\pm 𝒜_2`$. The sufficient condition is that there respectively exists an algebra homomorphism $`\phi _1^+`$ or an algebra homomorphism $`\phi _1^{}`$
$$\phi _1^\pm :𝒜_1>H^\pm 𝒜_1$$
(1.11)
acting as the identity on $`𝒜_1`$, namely for any $`a_1𝒜_1`$
$$\phi _1^\pm (a_1)=a_1.$$
(1.12)
(Note that, as a consequence of (1.12), $`\phi _1^\pm `$ is idempotent, $`(\phi _1^\pm )^2=\phi _1^\pm `$). Here $`𝒜_1>H^\pm `$ denotes the cross product between $`𝒜_1`$ and $`H^\pm `$. In other words, this amounts to assuming that $`\phi _1^+(H^+)`$ \[resp. $`\phi _1^{}(H^{})`$\] provides a realization of $`H^+`$ (resp. $`H^{}`$) within $`𝒜_1`$. In fact, $`\stackrel{~}{𝒜}_2`$ is found using the main result of this work:
###### Theorem 1
Let $`\{H,\}`$ be a quasitriangular Hopf algebra and $`H^+,H^{}`$ be Hopf subalgebras of $`H`$ such that $`H^+H^{}`$. Let $`𝒜_1,𝒜_2`$ be respectively a $`H^+`$\- and a $`H^{}`$-module algebra, so that we can define $`𝒜^+`$ as in (1.9), and $`\phi _1^+`$ be a homomorphism of the type (1.11), (1.12), so that we can define the “unbraiding” map $`\chi ^+:𝒜_2𝒜^+`$ by
$$\chi ^+(a_2):=\phi _1^+(^{(1)})(a_2^{(2)}).$$
(1.13)
Alternatively, let $`𝒜_1,𝒜_2`$ be respectively a $`H^{}`$\- and a $`H^+`$-module algebra, so that we can define $`𝒜^{}`$ as in (1.10), and $`\phi _1^{}`$ be a homomorphism of the type (1.11), (1.12), so that we can define the “unbraiding” map $`\chi ^{}:𝒜_2𝒜^{}`$ by
$$\chi ^{}(a_2):=\phi _1^{}(^{1(2)})(a_2^{1(1)}).$$
(1.14)
In either case $`\chi ^\pm `$ are then injective algebra homomorphisms and
$$[\chi ^\pm (a_2),𝒜_1]=0,$$
(1.15)
namely the subalgebras $`\stackrel{~}{𝒜}_2^\pm :=\chi ^\pm (𝒜_2)𝒜_2`$ commute with $`𝒜_1`$. Moreover $`𝒜^\pm =𝒜_1\stackrel{~}{𝒜}_2^\pm `$.
Proof . We start by recalling the content of the hypotheses stated in the theorem. The algebra $`𝒜_1>H^\pm `$ as a vector space is the tensor product of $`𝒜_1`$ and $`H^\pm `$, whereas its product law is obtained combining the product laws of these two tensor factors with the cross-product law,
$$a_1g=g_{(1)}(a_1g_{(2)}),$$
(1.16)
for any $`a_1𝒜_1`$ and $`gH^\pm `$. $`\phi _1^\pm `$ being an algebra homomorphism means that for any $`\xi ,\xi ^{}𝒜_1>H^\pm `$
$$\phi _1^\pm (\xi \xi ^{})=\phi _1^\pm (\xi )\phi _1^\pm (\xi ^{}).$$
(1.17)
For $`\xi a𝒜_1𝒜_1>H^\pm `$, $`\xi ^{}gH^\pm 𝒜_1>H^\pm `$ this implies
$$a\phi ^\pm (g)=\phi ^\pm (g_{(1)})(ag_{(2)})$$
(1.18)
Hereby we have also used (1.12) and (1.16). After these preliminaries, note that under the assumption (1.9), for any $`a_1𝒜_1`$ and $`a_2𝒜_2`$
$`a_1\chi ^+(a_2)`$ $`\stackrel{(\text{1.13})}{=}`$ $`a_1\phi _1^+(^{(1)})(a_2^{(2)})`$
$`\stackrel{(\text{1.18})}{=}`$ $`\phi ^+(_{(1)}^{(1)})(a_1_{(2)}^{(1)})(a_2^{(2)})`$
$`\stackrel{(\text{A.1.3})}{=}`$ $`\phi ^+(^{(1)})(a_1^{(1^{})})(a_2^{(2)}^{(2^{})})`$
$`\stackrel{(\text{1.9})}{=}`$ $`\phi ^+(^{(1)})(a_2^{(2)})a_1`$
$`\stackrel{(\text{1.13})}{=}`$ $`\chi ^+(a_2)a_1,`$
which proves (1.15) in this case. Moreover
$`\chi ^+(a_2a_2^{})`$ $`\stackrel{(\text{1.13})}{=}`$ $`\phi _1^+(^{(1)})(a_2a_2^{}^{(2)})`$
$`\stackrel{(\text{1.8})}{=}`$ $`\phi ^+(^{(1)})(a_2_{(1)}^{(2)})(a_2^{}_{(2)}^{(2)})`$
$`\stackrel{(\text{A.1.4})}{=}`$ $`\phi ^+(^{(1^{})}^{(1)})(a_2^{(2)})(a_2^{}^{(2^{})})`$
$`\stackrel{(\text{1.17})}{=}`$ $`\phi ^+(^{(1^{})})\phi (^{(1)})(a_2^{(2)})(a_2^{}^{(2^{})})`$
$`\stackrel{(\text{1.13})}{=}`$ $`\phi ^+(^{(1^{})})\chi ^+(a_2)(a_2^{}^{(2^{})})`$
$`\stackrel{(\text{1.15})}{=}`$ $`\chi ^+(a_2)\phi (^{(1^{})})(a_2^{}^{(2^{})})`$
$`\stackrel{(\text{1.13})}{=}`$ $`\chi ^+(a_2)\chi ^+(a_2^{}),`$
proving that $`\chi ^+`$ is a homomorphism. To prove injectivity we show that $`\chi ^+`$ can be inverted on $`\chi ^+(𝒜_2)`$, and the inverse is given by
$$(\chi ^+)^1(\stackrel{~}{a}_2)=V^1\left([\phi ^+(S^1^{(1)})\stackrel{~}{a}_2]^{(2)}\right)$$
(1.19)
where $`V𝒜_1`$ is the invertible element defined by $`V:=\phi _1^+(S^1^{(1)})^{(2)}`$ ($`V`$ is invertible because $``$ is). In fact,
$`V^1[\phi _1^+(S^1^{(1)})\chi ^+(a_2)]^{(2)}`$
$`\stackrel{(\text{1.13})}{=}`$ $`V^1[\phi _1^+(S^1^{(1)})\phi _1^+(^{(1^{})})(a_2^{(2^{})})]^{(2)}`$
$`\stackrel{(\text{A.1.5}),(\text{1.17})}{=}`$ $`V^1\left\{\phi _1^+[S^1(^{1(1^{})}^{(1)})](a_2^{1(2^{})})\right\}^{(2)}`$
$`\stackrel{(\text{1.8})}{=}`$ $`V^1\phi _1^+[S^1(^{1(1^{})}^{(1)})]_{(1)}^{(2)}(a_2^{1(2^{})})_{(2)}^{(2)}`$
$`\stackrel{(\text{A.1.4})}{=}`$ $`V^1\phi _1^+[S^1(^{1(1^{})}^{(1)}^{(1\mathrm{"})})]^{(2\mathrm{"})}(a_2^{1(2^{})})^{(2)})`$
$`\stackrel{(\text{1.7})}{=}`$ $`V^1\phi _1^+[S^1(^{(1\mathrm{"})})]^{(2\mathrm{"})}a_2`$
$`=`$ $`V^1Va_2=a_2.`$
In fact, if $`\phi _1^+`$ can be extended to an algebra homomorphism $`\phi _1:𝒜_1>H𝒜_1`$ a little calculation with the help of eq.’s (A.2.1), (A.1.7) shows that $`V=\phi _1(v)`$, where $`vH`$ is the invertible central element defined by (A.1.8). We know that $`𝒜_1\stackrel{~}{𝒜}_2^+𝒜^+`$. To prove that $`𝒜^+=𝒜_1\stackrel{~}{𝒜}_2^+`$ note first that by (1.9) any element in $`𝒜^+`$ can be written as a sum of products $`a_1a_2`$, with $`a_1𝒜_1`$ and $`a_2𝒜_2`$. So we need to show that
$$a_1a_2=b^{(1)}\chi ^+(b^{(2)})$$
(1.20)
for some $`b^{(1)}𝒜_1`$, $`b^{(2)}𝒜_2`$ (at the rhs a sum of many terms is implicitly understood). Now this can be proved as follows:
$`a_1a_2`$ $`=`$ $`a_1\phi _1^+(\mathrm{𝟏}_H)(a_2\mathrm{𝟏}_H)=a_1\phi _1^+(^{1(1)}^{(1^{})})[a_2(^{1(2)}^{(2^{})})]`$
$`\stackrel{(\text{1.7})}{=}`$ $`a_1\phi _1^+(^{1(1)})\phi _1^+(^{(1^{})})(a_2^{1(2)})^{(2^{})}`$
$`\stackrel{(\text{1.13})}{=}`$ $`a_1\phi _1^+(^{1(1)})\chi ^+(a_2^{1(2)}),`$
which is of the form (1.20).
The proof for $`\chi ^{}`$ under the corresponding assumptions is completely analogous. $``$$``$
In the next Section we shall need an alternative expression for $`\chi ^\pm `$, which we prove in the appendix:
###### Proposition 1
$`\chi ^+(a_2)=(a_2^{1(2)})\phi _1^+(S^{1(1)}),`$ (1.21)
$`\chi ^{}(a_2)=(a_2^{(1)})\phi _1^{}(S^{(2)}).`$ (1.22)
The rest of the paper is essentially devoted to illustrate the application of Theorem 1 to some algebras $`𝒜_i`$ for which homomorphisms $`\phi _1^\pm `$ are known. In Ref. algebra homomorphisms $`\phi _1^\pm `$ have been found for (a slightly enlarged version $`𝒜_1`$ of) the algebra of functions on the $`N`$-dimensional quantum Euclidean space $`_q^N`$, corresponding to $`H=U_qso(N)`$. The explicit forms of $`\phi _1^\pm `$ on the Faddeev-Reshetikhin-Takhtadjan (FRT) generators $`^\pm _j^i`$ of $`U_qso(N)`$ are recalled in the appendix A.3. The maps $`\phi _1^\pm `$ for $`N=3`$ are given also in Ref. . The same maps do the job also on the quotient spaces obtained by setting $`x^ix_i=1`$ \[quantum $`(N1)`$-dimensional spheres $`S_q^{N1}`$\], and the appropriate maps for the $`q`$-deformed fuzzy sphere $`S_{q,M}^2`$ have been found in . Therefore $`U_q\text{so}(N)`$and the quantum Euclidean spaces/spheres provide nontrivial $`H`$ and $`𝒜_1`$ for the application of the above theorem. In fact, the constructions of the frame given in Ref. can be interpreted as an application of the theorem with $`𝒜_1_q^N`$ and $`𝒜_2`$ the $`N!`$-dim exterior algebra generated by the differentials $`dx^i`$ of the $`U_q\text{so}(N)`$-covariant differential calculus (although with a universal $`R`$-matrix $``$ slightly modified by multiplication by the coproduct $`\mathrm{\Delta }(\mathrm{\Lambda })=\mathrm{\Lambda }\mathrm{\Lambda }`$ of a new element $`\mathrm{\Lambda }`$ generating dilatations); consequently, in agreement with the philosophy of Ref. , the algebra of differential forms on $`_q^N`$ can be written as $`_q^N\stackrel{~}{𝒜}_2`$, where $`\stackrel{~}{𝒜}_2`$ is the $`N!`$-dim exterior algebra generated by the frame elements. On the other hand, the existence of algebra homomorphisms $`\phi :𝒜_1>H𝒜_1`$, for $`H=U_q\text{so}(N),U_qsl(N)`$ and $`𝒜_1`$ respectively equal to (a suitable completion of) the $`U_q\text{so}(N)`$-covariant Heisenberg algebra or the $`U_qsl(N)`$-covariant Heisenberg or Clifford algebras, has been known for even a longer time , so the theorem also applies if we choose as $`(H,𝒜_1)`$ one of these pairs of algebras.
Of course the above theorem can be used iteratively to completely unbraid an algebra $`𝒜`$ obtained by repeated braided tensor product \[through prescription (1.9), or prescription (1.10)\] of an arbitrary number of $`H`$-module algebras $`𝒜_1,𝒜_2,\mathrm{},𝒜_M`$. We shall explicitly consider the particular case that the latter be $`M`$ identical copies of the $`U_q\text{so}(N)`$-covariant quantum space/sphere (Section 3), of the $`U_qso(3)`$-covariant $`q`$-fuzzy sphere (Section 6), or of the $`U_q\text{so}(N)`$\- or $`U_qsl(N)`$-covariant Heisenberg algebra (Section 4). There we shall explicitly write down the generators of $`\stackrel{~}{𝒜}_2^\pm `$ for the lowest $`N`$ examples.
In appendix A.3 we analyze the properties of $`\phi ^\pm `$ under the main real sections of $`U_q\text{so}(N)`$, what was left aside in Ref. . In section 2 we investigate in the context of general position the properties of $`\chi ^\pm `$ under the $``$-structures.
## 2 The unbraiding under the $``$-structures
Assume $`H`$ is a Hopf $``$-algebra, namely the coproduct $`\mathrm{\Delta }`$ and counit $`\epsilon `$ are $``$-homomorphisms,
$$\mathrm{\Delta }(g^{})(g^{})_{(1)}(g^{})_{(2)}=(g_{(1)})^{}(g_{(2)})^{},$$
(2.1)
and $`𝒜_1`$, $`𝒜_2`$ are $`H`$-module $``$-algebras, namely for any $`a_i𝒜_i`$
$$(a_ig)^{}=a_i^{}S^1g^{}$$
(2.2)
(here $`S`$ denotes the antipode of $`H`$); we have used and shall use the same symbol $``$ for the $``$-structure on all algebras $`H,𝒜_1`$, etc. Then $``$ is a $``$-structure also for $`𝒜_1>H`$. The same statement is not automatically true for the braided tensor product algebra $`𝒜^\pm =𝒜_1\underset{¯}{}^\pm 𝒜_2`$, because the basic requirement that the latter be antimultiplicative
$$(a_2a_1)^{}=a_1^{}a_2^{}$$
(2.3)
(note that this would make $`𝒜^\pm `$ also a $`H`$-module $``$-algebra) is not automatically guaranteed. In fact, applying this would-be $``$ to rhs of (1.9) one finds
$`(a_2a_1)^{}`$ $`\stackrel{(\text{1.9})}{=}`$ $`\left[(a_1^{(1)})(a_2^{(2)})\right]^{}`$ (2.4)
$`\stackrel{(\text{2.3})}{=}`$ $`(a_2^{(2)})^{}(a_1^{(1)})^{}`$
$`\stackrel{(\text{2.2})}{=}`$ $`(a_2^{}S^1^{(2)}{}_{}{}^{})(a_1^{}S^1^{(1)}{}_{}{}^{})`$
$`\stackrel{(\text{1.9})}{=}`$ $`(a_1^{}S^1^{(1)}{}_{}{}^{}_{}^{(1^{})})(a_2^{}S^1^{(2)}{}_{}{}^{}_{}^{(2^{})});`$
in order that this be equal to the rhs of (2.3) it is necessary that $`(S^1S^1)^{}=^1`$, which upon use of (A.1.5) is equivalent to
$$^{}=^1$$
(2.5)
(here $`^{}`$ means $`^{(1)}{}_{}{}^{}^{(2)}^{}`$). This condition is fulfilled only for the standard noncompact sections (A.1.28) of $`U_q\text{}`$, for $`|q|=1`$; as a consequence, $`𝒜^+=𝒜_1\underset{¯}{}^+𝒜_2`$ becomes a $`H`$-module $``$-algebra if one extends the $``$-structures of the tensor factors to $`𝒜`$ using (2.3). The same holds for $`𝒜^{}`$.
On the contrary, the compact section, which requires $`q`$, is characterized by
$$^{}=_{21}.$$
(2.6)
In the latter case the map $``$ introduced through (2.3) makes sense only as an involutive antimultiplicative antilinear map $`𝒜^+𝒜^{}`$, if both $`𝒜^+`$ and $`𝒜^{}`$ exist. In fact, in this case the last line in (2.4) will be replaced by
$$\stackrel{(\text{1.10}),(\text{2.6})}{=}(a_1^{}S^1^{(2)}^{1(2^{})})(a_2^{}S^1^{(1)}^{1(1^{})})\stackrel{(\text{A.1.5})}{=}a_1^{}a_2^{},$$
as required. Alternatively, if $`𝒜_1,𝒜_2`$ are two copies of the same algebra and we denote by $`\psi :𝒜_1𝒜_2`$ the map associating to each $`a_1𝒜_1`$ the equivalent element in $`𝒜_2`$, one can define an alternative $``$-structure $``$ in $`𝒜^\pm `$ by setting
$$a_1^{}=\psi (a_1^{})a_2^{}=\psi ^1(a_2^{}),$$
(2.7)
since this is instead compatible with (1.9). In fact, (2.4) will become
$`(a_2a_1)^{}`$ $`\stackrel{(\text{1.9})}{=}`$ $`\left[(a_1^{(1)})(a_2^{(2)})\right]^{}`$ (2.8)
$`\stackrel{(\text{2.3})}{=}`$ $`(a_2^{(2)})^{}(a_1^{(1)})^{}`$
$`\stackrel{(\text{2.2})}{=}`$ $`\psi ^1(a_2^{}S^1^{(2)}{}_{}{}^{})\psi (a_1^{}S^1^{(1)}{}_{}{}^{})`$
$`\stackrel{(\text{1.9})}{=}`$ $`\psi (a_1^{}S^1^{(2)}^{1(2^{})})\psi ^1(a_2^{}S^1^{(1)}^{1(1^{})})`$
$`\stackrel{(\text{A.1.5})}{=}`$ $`\psi (a_1^{})\psi ^1(a_2^{})a_1`$
$`\stackrel{(\text{2.7})}{=}`$ $`a_1^{}a_2^{}`$
A similar trick can be used also if one considers an iterated braided tensor product of $`M>2`$ copies of the same algebra, see next section. However, such $``$’s have not the standard commutative limit, because of the presence of the map $`\psi `$.
Inspired by the applications of the next two Sections, we now assume that $`\phi _1^\pm `$ fulfill some specific conditions relating its action before and after the application of the involution $``$, and analyze the identities relating the action of $`\chi ^\pm `$ before and after the application of $``$ which follow herefrom.
###### Proposition 2
Assume that the conditions of Theorem 1 for defining $`\chi ^+`$ or $`\chi ^{}`$ are fulfilled. If $`^{}=^1`$ and for any $`g^\pm H^\pm `$
$$[\phi _1^\pm (g^\pm )]^{}=\phi _1^\pm (g^\pm {}_{}{}^{}),$$
(2.9)
in other words $`\phi _1^\pm `$ are $``$-homomorphisms, then
$$[\chi ^\pm (a_2)]^{}=\chi ^\pm (a_2^{}).$$
(2.10)
If $`^{}=_{21}`$ and $`:H^\pm H^{}`$ fulfills
$$[\phi _1^\pm (g)]^{}=\phi _1^{}(g^{}),$$
(2.11)
then
$$[\chi ^\pm (a_2)]^{}=\chi ^{}(a_2^{}).$$
(2.12)
Proof . Under the first assumptions, for any $`a_2𝒜_2`$,
$`[\chi ^+(a_2)]^{}`$ $`\stackrel{(\text{1.21})}{=}`$ $`[(a_2^{1(2)})\phi _1^+(S^{1(1)})]^{}`$
$`=`$ $`[\phi _1^+(S^{1(1)})]^{}(a_2^{1(2)})^{}`$
$`\stackrel{(\text{2.9}),(\text{2.2})}{=}`$ $`\phi _1^+(S^1^{1(1)}{}_{}{}^{})(a_2^{}S^1^{1(2)}{}_{}{}^{})`$
$`\stackrel{(\text{2.5})}{=}`$ $`\phi _1^+(S^1^{(1)})(a_2^{}S^1^{(2)})`$
$`\stackrel{(\text{A.1.5})}{=}`$ $`\phi _1^+(^{(1)})(a_2^{}^{(2)})`$
$`\stackrel{(\text{1.13})}{=}`$ $`\chi ^+(a_2^{}).`$
Similarly one proves (2.10) for $`\chi ^{}`$. Under the second assumptions, for any $`a_2𝒜_2`$,
$`[\chi ^+(a_2)]^{}`$ $`\stackrel{(\text{1.21})}{=}`$ $`[(a_2^{1(2)})\phi _1^+(S^{1(1)})]^{}`$
$`=`$ $`[\phi _1^+(S^{1(1)})]^{}(a_2^{1(2)})^{}`$
$`\stackrel{(\text{2.11}),(\text{2.2})}{=}`$ $`\phi _1^{}(S^1^{1(1)}{}_{}{}^{})(a_2^{}S^1^{1(2)}{}_{}{}^{})`$
$`\stackrel{(\text{2.6})}{=}`$ $`\phi _1^{}(S^1^{1(2)})(a_2^{}S^1^{1(1)})`$
$`\stackrel{(\text{A.1.5})}{=}`$ $`\phi _1^{}(^{1(2)})(a_2^{}^{1(1)})`$
$`\stackrel{(\text{1.14})}{=}`$ $`\chi ^{}(a_2^{}).`$
By similar arguments one proves the claim for $`\chi ^{}`$. $``$$``$
It should be noted that there also exist non–standard star structures on $`U_q\text{}`$ for $`|q|=1`$, in particular the compact form $`X_{i}^{\pm }{}_{}{}^{}=X_i^{}`$, $`K_i^{}=K_i^1`$ in terms of the Cartan–Weyl generators. Then
$$^{}=_{21}^1,$$
(2.13)
while the coproduct does not fulfill (2.1) as in a standard Hopf $``$-algebra but becomes flipped under the star. This nevertheless has the correct classical limit, because the coproduct is cocommutative for $`q=1`$. In certain cases (in particular on the fuzzy quantum sphere discussed in section 6, but see also ), it is then possible to define a star structure on each $`𝒜_i`$, which takes the form $`a_{i;k}^{}=\pm \mathrm{\Omega }_ia_{i;k}\mathrm{\Omega }_i^1`$ on the generators $`a_{i;k}`$ of $`𝒜_i`$. Here $`\mathrm{\Omega }_i={}_{}{}^{4}\sqrt{v_i}_{}^{1}\omega _i`$, where $`v_i`$ and $`\omega _i`$ are the realizations in $`𝒜_i`$ (using an algebra map from $`U_q\text{}`$ to $`𝒜_i`$ as above) of the central element $`vU_q\text{}`$ (A.1.8) and the “universal Weyl element” $`\omega `$ in an extension of $`U_q\text{}`$ . All this must be defined in some representation of $`𝒜_i`$; for more details see<sup>*</sup><sup>*</sup>*The $`v`$ in is the square root of our $`v`$ here. . If moreover there exists an element $`\mathrm{\Omega }`$ which realizes $`{}_{}{}^{4}\sqrt{v}_{}^{1}\omega `$ in $`𝒜^+=𝒜_1\underset{¯}{}^+𝒜_2`$ or a “physical” subspace thereof, then it follows easily from (2.13) that the star structure $`a_{i;k}^{}=\pm \mathrm{\Omega }a_{i;k}\mathrm{\Omega }^1`$ on $`𝒜^+`$ is consistent with the commutation relations of the braided tensor product algebra $`𝒜^+`$. This star then has the correct classical limit, and the same construction also works for $`𝒜^{}`$.
## 3 Unbraiding ‘chains’ of braided quantum Euclidean spaces or spheres
In this section we consider the braided tensor product of $`M2`$ copies $`𝒜_1,𝒜_2`$, $`\mathrm{},𝒜_M`$ of the quantum Euclidean space $`_q^N`$ (the $`U_q\text{so}(N)`$-covariant quantum space), i.e. of the unital associative algebra generated by $`x^i`$ fulfilling the relations
$$𝒫_a{}_{hk}{}^{ij}x_{}^{h}x^k=0,$$
(3.1)
where $`𝒫_a`$ denotes the $`q`$-deformed antisymmetric projector appearing in the decomposition of the braid matrix $`\widehat{R}`$ of $`U_q\text{so}(N)`$ \[given in formula (A.1)\], or of the quotient space of the latter obtained by setting $`r^2:=x^ix_i=1`$ \[the quantum $`(N1)`$-dimensional sphere $`S_q^{N1}`$\]. The multiplet $`(x^i)`$ carries the fundamental vector representation $`\rho `$ of $`U_q\text{so}(N)`$: for any $`gU_q\text{so}(N)`$
$$x^ig=\rho _j^i(g)x^j.$$
(3.2)
We shall enumerate the different copies of the quantum Euclidean space or sphere by attaching an additional greek index to them, e.g. $`\alpha =1,2,\mathrm{},M`$. The prescription (1.10) to glue $`𝒜_1,\mathrm{},𝒜_M`$ into a $`U_q\text{so}(N)`$-module associative algebra $`𝒜^{}`$ gives the following cross commutation relations between their respective generators:
$$x^{\alpha ,i}x^{\beta ,j}=\widehat{R}_{hk}^{ij}x^{\alpha ,h}x^{\beta ,k}$$
(3.3)
whenever $`\alpha <\beta `$. Note that prescriptions (1.10), (1.9) go into each other under the inverse reordering $`1,2,\mathrm{},MM,\mathrm{},2,1`$. Applying iteratively Theorem 1 we shall be able to completely unbraid this iterated tensor product.
To define $`\phi _1^\pm `$ one actually needs a slightly enlarged version of $`_q^N`$ (or $`S_q^{N1}`$). One has to introduce some new generators $`\sqrt{r_a}`$, with $`1a\frac{N}{2}`$, together with their inverses $`(\sqrt{r_a})^1`$, requiring that
$$r_a^2=\underset{h=a}{\overset{a}{}}x^hx_h=\underset{h=a}{\overset{a}{}}g_{hk}x^hx^k$$
(3.4)
(note that, having set $`n:=\left[\frac{N}{2}\right]`$, $`r_n^2`$ coincides with $`r^2`$). Moreover for odd $`N`$ we add also also $`\sqrt{x^0}`$ and its inverse as new generators). In fact, the commutation relations involving these new generators can be fixed consistently, and turn out to be simply $`q`$-commutation relations. $`r`$ plays the role of ‘deformed Euclidean distance’ of the generic ‘point of coordinates’ $`(x^i)`$ of $`_q^N`$ from the ‘origin’; $`r_a`$ is the ‘projection’ of $`r`$ on the ‘subspace’ $`x^i=0`$, $`|i|>a`$. In the previous equation $`g_{hk}`$ denotes the ‘metric matrix’ of $`SO_q(N)`$:
$$g_{ij}=g^{ij}=q^{\rho _i}\delta _{i,j}.$$
(3.5)
It is a $`SO_q(N)`$-isotropic tensor and is a deformation of the ordinary Euclidean metric. Here and in the sequel $`n:=\left[\frac{N}{2}\right]`$ is the rank of $`so(N)`$, the indices take the values $`i=n,\mathrm{},1,0,1,\mathrm{}n`$ for $`N`$ odd, and $`i=n,\mathrm{},1,1,\mathrm{}n`$ for $`N`$ even. Moreover, we have introduced the notation $`(\rho _i)=(n\frac{1}{2},\mathrm{},\frac{1}{2},0,\frac{1}{2},\mathrm{},\frac{1}{2}n)`$ for $`N`$ odd, $`(n1,\mathrm{},0,0,\mathrm{},1n)`$ for $`N`$ even. In the case of even $`N`$ one needs to include also the FRT generator $`^{}_1^1`$ and its inverse $`^+_1^1`$ (which are generators of $`U_qso(N)`$ belonging to the Cartan subalgebra) among the generators of $`𝒜_1`$. They satisfy the commutation relations
$$^{}{}_{1}{}^{1}x_{}^{\pm 1}=q^{\pm 1}x^{\pm 1}^{}{}_{1}{}^{1},^{}{}_{1}{}^{1}x_{}^{\pm i}=x^{\pm i}^{}{}_{1}{}^{1}\text{ for }i>1.$$
(3.6)
with $`𝒜_1`$, and the standard FRT relations with the rest of $`U_qso(N)`$. One can easily show that the extension of the action of $`U_qso(N)`$ to $`\sqrt{r_a},(\sqrt{r_a})^1`$ is uniquely determined by the constraints the latter fulfil; it is a bit complicated and therefore will be omitted, since we will not need its explicit expression. The action of $`H`$ on $`^{}_1^1`$ is the standard (right) adjoint action. Note that the maps $`\phi _1^\pm `$ have no analog in the “undeformed” case ($`q=1`$), because $`𝒜_1^N`$ is abelian, whereas $`HU_qso(N)`$ is not.
The unbraiding procedure is recursive. We use the homomorphism $`\phi _1`$ found in Ref. and start by unbraiding the first copy from the others. Following Theorem 1, we perform the following change of generators in $`𝒜^{}`$
$$\begin{array}{c}y^{1,i}:=x^{1,i}\hfill \\ y^{\alpha ,i}:=\chi ^{}(x^{\alpha ,i})\stackrel{(\text{1.14})}{=}\phi _1(^{1(2)})\rho _j^i(^{1(1)})x^{\alpha ,j}=\phi _1(^{}{}_{j}{}^{i})x^{\alpha ,j},\alpha >1.\hfill \end{array}$$
In the last equality we have used the definition (A.1.15) of the FRT generators of $`U_q\text{so}(N)`$. In appendix A.3 we recall the $`\phi ^\pm `$ images of the latter. In view of formula (A.3.2) we thus find
$$\begin{array}{c}y^{1,i}:=x^{1,i}\hfill \\ y^{\alpha ,i}:=g^{ih}[\mu _h^1,x^{1,k}]_qg_{kj}x^{\alpha ,j},\alpha >1.\hfill \end{array}$$
(3.7)
The suffix 1 in $`\mu _a^1`$ means that the special elements $`\mu _a`$ defined in (A.3.3) must be taken as elements of the first copy of $`_q^N`$. In view of (A.3.3) we see that $`g^{ih}[\mu _h^1,x^{1,k}]_qg_{kj}`$ are rather simple polynomials in $`x^i`$ and $`r_a^1`$, homogeneous of total degree 1 in the coordinates $`x^i`$ and $`r_a`$. Hence (3.7) is a transformation of polynomial type and therefore likely to be implemented as a well-defined operator transformation also when representing $`𝒜^{}`$ as an algebra of operators on some linear space. Using the results (A.3.6) given in the appendix we give now the explicit expression of (3.7) for $`N=3`$:
$`y^{\alpha ,}=qh\gamma _1{\displaystyle \frac{r}{x^0}}x^{\alpha ,}`$
$`y^{\alpha ,0}=\sqrt{q}(q+1){\displaystyle \frac{1}{x^0}}x^+x^{\alpha ,}+x^{\alpha ,0}`$ (3.8)
$`y^{\alpha ,+}={\displaystyle \frac{\sqrt{q}(q+1)}{h\gamma _1rx^0}}(x^+)^2x^{\alpha ,}+{\displaystyle \frac{q^1+1}{h\gamma _1r}}x^+x^{\alpha ,0}{\displaystyle \frac{1}{qh\gamma _1r}}x^0x^{\alpha ,+}`$
for any $`\alpha =2,\mathrm{},M`$. Here we have set $`x^ix^{1,i}`$, $`h\sqrt{q}1/\sqrt{q}`$, replaced for simplicity the values $`1,0,1`$ of the indices by the ones $`,0,+`$ and denoted by $`\gamma _1`$ a free parameter.
As a consequence of the theorem we find
###### Corollary 1
$`[y^{1,i},y^{\alpha ,j}]=0\alpha >1`$ (3.9)
$`y^{\alpha ,i}y^{\beta ,j}=\widehat{R}_{hk}^{ij}y^{\alpha ,h}y^{\beta ,k}1<\alpha <\beta `$ (3.10)
$`𝒫_a{}_{hk}{}^{ij}y_{}^{\alpha ,h}y^{\alpha ,k}=0`$ (3.11)
By (3.9) the subalgebra $`\stackrel{~}{𝒜}_1^{}𝒜_1`$ of $`𝒜^{}`$ generated by $`y^{1,i}x^{1,i}`$ commutes with the subalgebra generated by $`y^{2,i},\mathrm{},y^{M,i}`$, which we shall call $`\stackrel{~}{𝒜}^{}`$. This was the first step of the unbraiding procedure. Now we can reiterate the latter for $`\stackrel{~}{𝒜}^{}`$, with $`y^{2,i}`$ playing the role of $`x^{1,i}`$. After $`M1`$ steps, we shall have determined $`M`$ independent commuting subalgebrae of $`𝒜^{}`$ which we shall call $`\stackrel{~}{𝒜}_\alpha ^{}`$, $`\alpha =1,\mathrm{},M`$.
Theunbraiding procedure for the alternative braided tensor product stemming from prescription (1.9) arises by iterating the change of generators
$$\begin{array}{c}y^{}{}_{}{}^{M,i}:=x^{M,i}\hfill \\ y^{}{}_{}{}^{\alpha ,i}:=\phi _M(^+{}_{j}{}^{i})x^{\alpha ,j}=g^{ih}[\overline{\mu }_h^M,x^{M,k}]_{q^1}g_{kj}x^{\alpha ,j},\alpha <M.\hfill \end{array}$$
(3.12)
$`\overline{\mu }_a^M`$ are the special elements defined in (A.3.5) belonging to the $`M`$-th copy of $`_q^N`$. Using the results (A.3.7) given in the appendix we give the explicit expression of (3.12) for $`N=3`$: for any $`\alpha =1,\mathrm{},M1`$,
$`y^{}{}_{}{}^{\alpha ,}=h\overline{\gamma }_1{\displaystyle \frac{z^0}{r_z}}x^{\alpha ,}+{\displaystyle \frac{k\overline{\gamma }_1}{\sqrt{q}r_z}}z^{}x^{\alpha ,0}+{\displaystyle \frac{q^2k\overline{\gamma }_1}{r_zz^0}}(z^{})^2x^{\alpha ,+}`$
$`y^{}{}_{}{}^{\alpha ,0}=x^{\alpha ,0}+q^{\frac{1}{2}}(q^1+1){\displaystyle \frac{1}{z^0}}z^{}x^{\alpha ,+}`$ (3.13)
$`y^{}{}_{}{}^{\alpha ,+}={\displaystyle \frac{r_z}{h\overline{\gamma }_1z^0}}x^{\alpha ,+}`$
Here we have set $`z^ix^{M,i}`$, $`r_z^2x^{M,i}x_i^M`$, $`kqq^1`$ and $`\overline{\gamma }_1`$ is a free parameter.
Again, the subalgebra $`\stackrel{~}{𝒜}_M^+_q^N`$ of $`𝒜^+`$ generated by $`y^{M,i}x^{M,i}`$ commutes with the subalgebra generated by $`y^{1,i},\mathrm{},y^{M1,i}`$, which we shall call $`\stackrel{~}{𝒜}^+`$. This was the first step of the unbraiding procedure. Now we can reiterate the latter for $`\stackrel{~}{𝒜}^+`$, with $`y^{M1,i}`$ playing the role of $`x^{M,i}`$. After $`M1`$ steps, we shall have determined $`M`$ independent commuting subalgebrae of $`𝒜^+`$ which we shall call $`\stackrel{~}{𝒜}_\alpha ^+`$.
We summarize the results of this section:
###### Proposition 3
Let $`𝒜_1,𝒜_2,\mathrm{},𝒜_M`$ be $`M`$ copies of the $`U_q\text{so}(N)`$-covariant quantum Euclidean space (or sphere). Then $`𝒜_1\underset{¯}{}^\pm 𝒜_2\underset{¯}{}^\pm \mathrm{}\underset{¯}{}^\pm 𝒜_M=𝒜_1\stackrel{~}{𝒜}_2^\pm \mathrm{}\stackrel{~}{𝒜}_M^\pm `$, where $`\stackrel{~}{𝒜}_2^\pm ,\mathrm{},\stackrel{~}{𝒜}_M^\pm `$ are subalgebras of the lhs isomorphic to $`𝒜_1`$.
By a suitable choice of $`\gamma _1,\overline{\gamma }_1`$, as well as of the other free parameters appearing in the definitions of $`\phi ^\pm `$ for $`N>3`$ (see appendix A.3), one can make $`\phi ^\pm `$ into $``$-homomorphisms when $`|q|=1`$, and make them satisfy the relation
$$[\phi ^\pm (g)]^{}=\phi ^{}(g^{})$$
(3.14)
when $`q^+`$. Since these relations are of the type considered in proposition 2.12, the claims of the latter for $`\chi ^\pm `$ and their consequences hold. In particular, when $`|q|=1`$ one has a well-defined $``$ on the braided tensor product of $`𝒜_1,\mathrm{},𝒜_M`$ mapping each of the independent, commuting subalgebras $`\stackrel{~}{𝒜}_i^\pm `$ into itself. On the contrary for real $`q`$ one can consider the map $`:𝒜^+𝒜^{}`$ defined by (2.3) or a $``$-structure on $`𝒜^\pm `$ defined in a way similar to what we have done in (2.7),
$$(x^{\alpha ,i})^{}=x^{M\alpha +1,j}g_{ji}.$$
(3.15)
The latter has not the standard classical limit. A short calculation shows that the latter implies
$$(y^{\alpha ,i})^{}=y^{}{}_{}{}^{M\alpha +1,j}g_{ji}^{}.$$
(3.16)
## 4 Unbraiding ‘chains’ of braided Heisenberg algebras
In this section we consider the braided tensor product of $`M2`$ copies $`𝒜_1,𝒜_2`$, $`\mathrm{},𝒜_M`$ of the $`U_q\text{g}`$ -covariant deformed Heisenberg algebra $`𝒟_{ϵ,\text{}}`$, $`\text{}=sl(N)`$, $`so(N)`$ , i.e. the unital associative algebra generated by $`x^i,_j`$ fulfilling the relations
$`𝒫_a{}_{hk}{}^{ij}x_{}^{h}x^k=0`$
$`𝒫_a{}_{hk}{}^{ij}_{j}^{}_i=0`$ (4.1)
$`_ix^j=\delta _j^i+(q\gamma \widehat{R})^ϵ{}_{ih}{}^{jk}x_{}^{h}_k,`$
where $`\gamma =q^{\frac{1}{N}},1`$ respectively for $`\text{}=sl(N),so(N)`$, and the exponent $`ϵ`$ can take either value $`ϵ=1,1`$. $`\widehat{R}`$ denotes the braid matrix of $`U_q\text{g}`$ \[given in formulae (A.1.22) and (A.1)\], and again $`𝒫_a`$ the antisymmetric projector appearing in the decomposition of the latter. The coordinates $`x^i`$ transform according to the fundamental vector representation of $`U_q\text{g}`$ , as in (3.2), whereas the ‘partial derivatives’ trasform according the contragredient representation,
$$_ig=_h\rho _i^h(S^1g).$$
(4.2)
The indices will take the values $`i=1,\mathrm{},N`$ if $`\text{}=sl(N)`$, the same values considered in the previous section if $`\text{}=so(N)`$. Clearly in the latter case $`𝒟_{ϵ,\text{}}`$ has the quantum Euclidean space generated by $`x^i`$ as a module subalgebra.
Again, we shall enumerate the different copies by attaching to them an additional greek index, e.g. $`\alpha =1,2,\mathrm{},M`$. The prescription (1.9) to glue $`𝒜_1,\mathrm{},𝒜_M`$ into a $`U_q\text{g}`$ -module associative algebra $`𝒜^+`$ (see also Ref. ) gives the following cross commutation relations between their respective generators
$$\begin{array}{cc}x^{\alpha ,i}x^{\beta ,j}=\widehat{R}_{hk}^{ij}x^{\beta ,h}x^{\alpha ,k}\hfill & _{\alpha ,i}_{\beta ,j}=\widehat{R}_{ji}^{kh}_{\beta ,h}_{\alpha ,k}\hfill \\ _{\alpha ,i}x^{\beta ,j}=\widehat{R}^1{}_{ik}{}^{jh}x_{}^{\beta ,k}_{\alpha ,h}\hfill & _{\beta ,i}x^{\alpha ,j}=\widehat{R}_{ik}^{jh}x^{\alpha ,k}_{\beta ,h}\hfill \end{array}$$
(4.3)
when $`\alpha >\beta `$. With respect to Ref. we have called the generators $`x^i,_j`$ instead of $`A^i,A_j^+`$, inverted the order of the product due to covariance w.r.t. the right (instead of the left) $`U_q\text{g}`$ -action, and for the sake of simplicity we have put equal to one possible factors at the rhs of (4.3).
In Ref. algebra homomorphisms $`\phi :𝒟_{ϵ,\text{}}>H𝒟_{ϵ,\text{}}`$ have been determined for $`\text{}=so(n)`$ and $`\text{}=sl(N),so(N)`$ respectively. This is the $`q`$-analog of vector field realization of g on the corresponding g -covariant (undeformed) space, e.g. $`\phi _1(E_j^i)=x^i_j\frac{1}{N}\delta _j^i`$ in the $`\text{}=sl(N)`$ case. The searched maps $`\phi ^\pm `$ will be simply the restrictions of $`\phi `$ to $`𝒟_{ϵ,\text{}}>H^\pm `$. In Ref. there are among others the $`\phi `$-images of the Chevalley generators of $`U_q\text{so}(N)`$One should take care of the fact that in Ref. we considered $`U_q\text{so}(N)`$acting by a left action, instead of a right one, what manifests itself in a replacement $`qq^1`$, or equivalently in an opposite coproduct. The rules for passing from right to left are described in Sect. 5., in Ref. there are the $`\phi `$-images of the generators of $`U_q\text{g}`$ playing the role of “vector fields” on $`G_q`$. By the change of generators described in Ref. one can easily pass from the Chevalley to the FRT generators $`^\pm _j^i`$ (A.1.15), whereas the relation between the latter and the vector fields is recalled in (A.4.2). The FRT generators are the ones explicitly needed in writing down $`\chi ^\pm (x^i)`$ and $`\chi ^\pm (_i)`$. For example, for $`\text{}=sl(2)`$ and $`ϵ=1`$ one finds
$`\phi (^+{}_{1}{}^{1})=\phi (^{}{}_{2}{}^{2})=[\phi (^{}{}_{1}{}^{1})]^1=[\phi (^+{}_{2}{}^{2})]^1=\alpha \mathrm{\Lambda }^{\frac{1}{2}}[1+(q^21)x^2_2]^{\frac{1}{2}}`$
$`\phi (^+{}_{2}{}^{1})=\alpha kq^1\mathrm{\Lambda }^{\frac{1}{2}}[1+(q^21)x^2_2]^{\frac{1}{2}}x^1_2`$ (4.4)
$`\phi (^{}{}_{1}{}^{2})=\alpha kq^3\mathrm{\Lambda }^{\frac{1}{2}}[1+(q^21)x^2_2]^{\frac{1}{2}}x^2_1,`$
where $`\alpha `$ is fixed by (A.1.14) to be $`\alpha =\pm 1,\pm i`$ and we have set
$$\mathrm{\Lambda }^2:=1+(q^21)x^i_i.$$
(4.5)
Whereas for $`\text{}=so(3)`$ and $`ϵ=1`$ one finds on the positive Borel subalgebra
$$\begin{array}{c}\phi (^+{}_{}{}^{})=\alpha \mathrm{\Lambda }[1+(q1)x^0_0+(q^21)x^+_+]\hfill \\ \phi (^+{}_{0}{}^{})=\alpha k\mathrm{\Lambda }(x^{}_0\sqrt{q}x^0_+)\hfill \\ \phi (^+{}_{+}{}^{})=\frac{1}{1+q^1}\phi (^+{}_{0}{}^{})\phi (^+{}_{+}{}^{0})\hfill \\ \phi (^+{}_{0}{}^{0})=1\hfill \\ \phi (^+{}_{+}{}^{0})=q^{\frac{1}{2}}[\phi (^+{}_{}{}^{})]^1\phi (^+{}_{0}{}^{})\hfill \\ \phi (^+{}_{+}{}^{+})=[\phi (^+{}_{}{}^{})]^1\hfill \end{array}$$
(4.6)
and on the negative Borel subalgebra
$$\begin{array}{c}\phi (^{}{}_{}{}^{})=\left(\alpha \mathrm{\Lambda }[1+(q1)x^0_0+(q^21)x^+_+]\right)^1\hfill \\ \phi (^{}{}_{}{}^{0})=\alpha q^2k\phi (^{}{}_{}{}^{})\mathrm{\Lambda }(x^0_{}\sqrt{q}x^+_0)\hfill \\ \phi (^{}{}_{}{}^{+})=\frac{1}{1+q}\phi (^{}{}_{0}{}^{+})\phi (^{}{}_{}{}^{0})\hfill \\ \phi (^{}{}_{0}{}^{0})=1\hfill \\ \phi (^{}{}_{0}{}^{+})=\alpha q^{\frac{3}{2}}k\mathrm{\Lambda }(x^0_{}\sqrt{q}x^+_0)\hfill \\ \phi (^{}{}_{+}{}^{+})=[\phi (^{}{}_{}{}^{})]^1.\hfill \end{array}$$
(4.7)
Here we have set
$$\mathrm{\Lambda }^2:=[1+(q^21)x^i_i+\frac{(q^21)^2}{\omega _1^2}(g_{ij}x^ix^j)(g^{hk}_k_h)],$$
(4.8)
where
$$\omega _a:=(q^{\rho _a}+q^{\rho _a}),$$
and replaced for simplicity the values $`1,0,1`$ of the indices by the ones $`,0,+`$. In either case the $`\phi `$-images of $`^+_j^i`$ and $`^{}_i^j`$ for $`i>j`$ vanish, because the latter do.
We see that strictly speaking $`\phi `$ takes values in some appropriate completion of $`𝒟_{ϵ,\text{}}`$, containing at least the square root and inverse square root of the polynomial $`\mathrm{\Lambda }^2`$ respectively defined in (4.5), (4.8), as well as the square root of $`[1+(q^21)x^2_2]`$ and its inverse, when $`\text{}=sl(2)`$, and the inverses (4.6)<sub>6</sub>, (4.7)<sub>6</sub>, when $`\text{}=so(3)`$. Apart from this minimal completion, another possible one is the socalled $`h`$-adic, namely the ring of formal power series in $`h=\mathrm{log}q`$ with coefficients in $`𝒟_{ϵ,\text{}}`$. Other completions, e.g. in operator norms, can be considered according to the needs. One can easily show that the extention of the action of $`H`$ to any such completion is uniquely determined (we omit to write down its explicit expression, since we don’t need it).
According to the main theorem, we set
$$\begin{array}{ccc}& & y^{1,i}x^{1,i}\hfill \\ & & _{y,1,a}_{1,a}\hfill \\ & & y^{\alpha ,i}\chi ^{}(x^{\alpha ,i})=\phi _1(^{}{}_{j}{}^{i})x^{\alpha ,j}\alpha >1\hfill \\ & & _{y,\alpha ,a}\chi ^{}(_{\alpha ,a})=\phi _1(S^{}{}_{a}{}^{d})_{\alpha ,d}\alpha >1\hfill \end{array}$$
(4.9)
and we find
###### Corollary 2
$`𝒫_a{}_{hk}{}^{ij}y_{}^{\alpha ,h}y^{\alpha ,k}=0`$
$`𝒫_a{}_{ij}{}^{hk}_{y\alpha ,k}^{}_{y,\alpha ,h}=0`$ (4.10)
$`_{y,\alpha ,i}y^{\alpha ,j}=\delta _i^j+(q\widehat{R})^{ϵ_\alpha }{}_{im}{}^{jl}y_{}^{\alpha ,m}_{y,\alpha ,l}`$
for all $`\alpha =1,\mathrm{},M`$, together with
$$\begin{array}{cc}[y^{1,i},y^{\alpha ,j}]=0\hfill & [_{y,1,i},y^{\alpha ,j}]=0\hfill \\ [_{y,\alpha ,i},y^{1,j}]=0\hfill & [_{y,1,i},_{y,\alpha ,j}]=0\hfill \end{array}$$
(4.11)
when $`\alpha >1`$, and
$$\begin{array}{c}y^{\alpha ,i}y^{\beta ,j}=\widehat{R}_{hk}^{ij}y^{\beta ,h}y^{\alpha ,k}\hfill \\ _{y,\alpha ,i}_{y,\beta ,j}=\widehat{R}_{ji}^{kh}_{y,\beta ,h}_{y,\alpha ,k}\hfill \\ _{y,\alpha ,i}y^{\beta ,j}=\widehat{R}^1{}_{ik}{}^{jh}y_{}^{\beta ,k}_{y,\alpha ,h}\hfill \\ _{y,\beta ,i}y^{\alpha ,j}=\widehat{R}_{ik}^{jh}y^{\alpha ,k}_{y,\beta ,h}\hfill \end{array}$$
(4.12)
when $`1<\beta <\alpha `$.
By (4.11) $`y^{1,i}x^{1,i}`$ and $`_{y,1,i}_{1,i}`$ commute with the subalgebra generated by $`y^{2,i},\mathrm{},y^{M,i}`$ and $`_{y,2,i},\mathrm{},_{y,M,i}`$ which we shall call $`\stackrel{~}{𝒜}^+`$. This was the first step of the unbraiding procedure. Now we can reiterate the latter for $`\stackrel{~}{𝒜}^+`$, with $`y^{2,i},_{y,2,i}`$ playing the role of $`x^{1,i},_{1,i}`$. After $`M1`$ steps, we shall have determined $`M`$ independent commuting subalgebras of $`𝒜^+`$ which we shall call $`\stackrel{~}{𝒜}_\alpha ^+`$.
For the sake of brevity we omit the unbraiding procedure for the alternative braided tensor product algebra stemming from prescription (1.10), which can be found following arguments completely analogous to the ones presented at the end of Section 3. We summarize the results of this section by
###### Proposition 4
Let $`𝒜_1,𝒜_2,\mathrm{},𝒜_M`$ be $`M`$ copies of the $`U_q\text{g}`$ -covariant Heisenberg algebra $`𝒟_{ϵ,\text{}}`$, $`\text{}=sl(N)`$, $`so(N)`$. Then $`𝒜_1\underset{¯}{}^\pm 𝒜_2\underset{¯}{}^\pm \mathrm{}\underset{¯}{}^\pm 𝒜_M=𝒜_1\stackrel{~}{𝒜}_2^\pm \mathrm{}\stackrel{~}{𝒜}_M^\pm `$, where $`\stackrel{~}{𝒜}_2^\pm ,\mathrm{},\stackrel{~}{𝒜}_M^\pm `$ are subalgebras of the lhs isomorphic to $`𝒟_{ϵ,\text{}}`$.
Relations (A.1.28), (2.2), (3.2), (4.2) and (4.1) fix the $``$-structure of $`𝒜_1`$ to be
$$(x^i)^{}=x^i,(_i)^{}=_i\{\begin{array}{cc}q^{\pm 2(Ni+1)}\text{ if }H=U_qsl(N)\hfill & \\ q^{\pm N+\rho _i}\text{ if }H=U_qso(N)\hfill & \end{array}$$
(4.13)
if $`|q|=1`$, and
$$(x^h)^{}=x^kg_{kh},(_i)^{}=\frac{\mathrm{\Lambda }^{\pm 2}}{q^{\pm N}+q^{\pm 2}}[(g^{jh}_h_j),x^i]$$
(4.14)
if $`H=U_qso(N)`$ and $`q^+`$. The upper or lower sign respectively refer to the choices $`ϵ=1,1`$ in (4.1)<sub>3</sub>, and $`\mathrm{\Lambda }^{\pm 2}`$ are respectively defined by
$$\mathrm{\Lambda }^{\pm 2}:=\left[1+(q^{\pm 2}1)x^i_i+\frac{(q^{\pm 2}1)^2}{\omega _n^2}r^2(g^{ji}_i_j)\right]^1.$$
(4.15)
The map $`\phi `$ is a $``$-homomorphism both for $`q`$ real and $`|q|=1`$. If we denote by $`\phi ^\pm `$ its restrictions to $`𝒜>H^\pm `$, then they are $``$-homomorphisms when $`|q|=1`$ (see appendix A.4), and fulfill the relation
$$[\phi ^\pm (g)]^{}=\phi ^{}(g^{})$$
(4.16)
when $`q^+`$ . Since these relations are of the type considered in proposition 2.12, the claims of the latter for $`\chi ^\pm `$ and their consequences hold. In particular, when $`|q|=1`$ one has a well-defined $``$ on the braided tensor product of $`𝒜_1,\mathrm{},𝒜_M`$ mapping each of the independent, commuting subalgebras $`\stackrel{~}{𝒜}_\alpha ^\pm `$ into itself.
Finally, the above results have an important corollary. According to Hochschild cohomology arguments developed by Gerstenhaber and applicable to Heisenberg algebras because of the results found by Du Cloux in Ref. , any deformed Heisenberg algebra, in particular the braided tensor products of $`𝒜_1,\mathrm{},𝒜_M`$ considered in this section, can be realized simply by a change of generators in the $`h`$-adic completion, $`h=\mathrm{log}q`$, of its undeformed counterpart (but in general not in other, e.g. operator-norm, completions). However explicit realizations are not provided by these results. The results presented here, combined to some older ones, allow to determine one such realization. In Ref. Ogievetsky found an explicit realization $`\varphi `$ or ‘deforming map’ of the elements of $`𝒟_{ϵ,\text{}}`$ in terms of formal power series in $`h=\mathrm{log}q`$ with coefficients in the corresponding undeformed Heisenberg algebra. Another, less explicit, one was found in Ref. . The composition of the unbraiding map found in this section, which allows to ‘decouple’ $`M`$ different copies of $`𝒟_{ϵ,\text{}}`$ from each other, with the map $`\varphi `$ provides an explicit realization or ‘deforming map’ of the larger Heisenberg algebra $`𝒜`$ (what we have called the ‘braided chain of Heisenberg algebras’), in the $`h`$-adic completion of the undeformed $`(NM)`$-dimensional Heisenberg algebra.
## 5 Formulae for the left action
For psychological reasons we often prefer to work with a left action rather than with a right one. In this section we give the analogs for left $`H`$-module algebras of the main results found so far for right $`H`$-module algebras. The left action of $`gH`$ on a product fulfills
$`(gg^{})a=g(g^{}a),`$ (5.1)
$`g(aa^{})=(g_{(1)}a)(g_{(2)}a^{}).`$ (5.2)
The product laws in the braided tensor product algebras $`\widehat{𝒜}^+,\widehat{𝒜}^{}`$ are respectively given by
$`a_2a_1=(^{1(1)}a_1)(^{1(2)}a_2).`$ (5.3)
$`a_2a_1=(^{(2)}a_1)(^{(1)}a_2),`$ (5.4)
The analog of Theorem 1 reads
###### Theorem 2
Let $`\{H,\}`$ be a quasitriangular Hopf algebra and $`H^+,H^{}`$ be Hopf subalgebras of $`H`$ such that $`H^+H^{}`$. Let $`\widehat{𝒜}_1,\widehat{𝒜}_2`$ be respectively a (left) $`H^+`$\- and a $`H^{}`$-module algebra, so that we can define $`\widehat{𝒜}^+`$ as in (5.3), and $`\widehat{\phi }_1^+:H^+<\widehat{𝒜}_1\widehat{𝒜}_1`$ be an algebra homomorphism fulfilling (1.12), so that we can define a map $`\widehat{\chi }^+:\widehat{𝒜}_2\widehat{𝒜}^+`$ by
$$\widehat{\chi }^+(a_2):=(^{(2)}a_2)\widehat{\phi }_1^+(^{(1)}).$$
(5.5)
Alternatively, let $`\widehat{𝒜}_1,\widehat{𝒜}_2`$ be respectively a (left) $`H^{}`$\- and a $`H^+`$-module algebra, so that we can define $`\widehat{𝒜}^{}`$ as in (5.4), and $`\widehat{\phi }_1^+:H^+<\widehat{𝒜}_1\widehat{𝒜}_1`$ be an algebra homomorphism fulfilling (1.12), so that we can define a map $`\widehat{\chi }^{}:\widehat{𝒜}_2\widehat{𝒜}^{}`$ by
$$\widehat{\chi }^{}(a_2):=(^{1(1)}a_2)\widehat{\phi }_1^{}(^{1(2)}).$$
(5.6)
In either case $`\widehat{\chi }^\pm `$ are then injective algebra homomorphisms and
$$[\widehat{\chi }^\pm (a_2),\widehat{𝒜}_1]=0,$$
(5.7)
namely the subalgebras $`\stackrel{~}{\widehat{𝒜}}_2^\pm :=\widehat{\chi }^\pm (\widehat{𝒜}_2)\widehat{𝒜}_2`$ commute with $`\widehat{𝒜}_1`$. Moreover $`\widehat{𝒜}^\pm =\widehat{𝒜}_1\stackrel{~}{\widehat{𝒜}}_2^\pm `$.
The results of section 2 apply without modifications (one just has to place a $`\widehat{}`$ in the appropriate places).
To enumerate the generators of the algebras considered in Sections 3,4 we shall exchange lower with upper indices, so the generators will read $`x_{\alpha ,i},^{\alpha ,i}`$. This is necessary if we wish the $`x`$’s to carry what we shall consider the fundamental (vector) representation $`\rho `$ of $`U_q\text{g}`$ ,
$$gx_i=x_j\rho _i^j(g),$$
(5.8)
rather than its contragredient $`\rho ^TS`$, because this follows from the row$`\times `$column multiplication law $`\rho _h^i(gg^{})=\rho _j^i(g)\rho _h^j(g^{})`$. Apart from this replacement, all the commutation relations remain the same, but can be rephrased in an equivalent way exchanging lower with upper indices also in the braid matrices and in the projectors $`𝒫_a`$, because $`\widehat{R}^T=\widehat{R}`$, $`𝒫_a{}_{}{}^{T}=𝒫_a`$. For instance, the analog of (3.1) will read
$$𝒫_a{}_{ij}{}^{hk}x_{h}^{}x_k=0.$$
(5.9)
The analogs of (3.2), (4.2) read
$`gx_i=\rho _i^j(g)x_j`$ (5.10)
$`g^i=^h\rho _h^i(Sg).`$ (5.11)
Algebra homomorphisms $`\widehat{\phi }_1^\pm `$ for the algebras considered in Sections 3,4 are immediately obtained in terms of the $`\phi _1^\pm `$ described there, according to the rule
$$\widehat{\phi }_1^\pm (^\pm {}_{j}{}^{h})=U^1{}_{a}{}^{j}\phi _{1}^{}(^{}{}_{b}{}^{a})U_h^b.$$
(5.12)
Here
$$U_c^b:=\rho _c^b(u),$$
(5.13)
$`uH`$ is a special element as in (A.1.6 ), and at the rhs the correct expression in the new notation has lower and upper indices exchanged. If $`\widehat{𝒜}_1`$ is the quantum Euclidean space $`_q^N`$ one finds, for instance,
$$\widehat{\phi }_1^{}(^{}{}_{j}{}^{h})=U^1{}_{a}{}^{j}g_{ac}^{}[\overline{\mu }^c,x_k]_{q^1}g^{kb}U_b^h\stackrel{(\text{A.1.30})}{=}g_{cj}[\overline{\mu }^c,x_k]_qg^{hk},$$
(5.14)
where $`\mu ^c`$ is the same as $`\mu _c`$ \[see A.3.3)\], but in the new notation. For instance, when $`|c|>1`$ it reads
$$\overline{\mu }^c=\overline{\gamma }_cr_{|c|}^1r_{|c|1}^1x_c,$$
(5.15)
with $`\gamma _c`$ defined as in (A.3) and $`r_a`$ ($`a0`$) defined by the condition
$$r_a^2=\underset{h=a}{\overset{a}{}}x_hx^h=\underset{h=a}{\overset{a}{}}g^{hk}x_hx_k.$$
The analog of (3.7) is therefore (with $`\alpha >1`$)
$`y_{1,i}:=x_{1,i}`$ (5.16)
$`y^{\alpha ,i}:=\widehat{\chi }^{}(x_{\alpha ,i})=x_{\alpha ,j}\widehat{\phi }_1(^{}{}_{i}{}^{j})=x_{\alpha ,j}g_{hi}[\overline{\mu }^{1,h},x_{1,k}]_{q^1}g^{jk}.`$ (5.17)
## 6 Unbraiding ‘chains’ of fuzzy quantum spheres
As a last example, we consider the braided tensor product of $`M`$ copies $`𝒜_1,\mathrm{},𝒜_M`$ of the $`q`$–deformed fuzzy sphere $`\widehat{S}_{q,N}^2`$ To relate this to our conventions, the $`q`$ in should be replaced by $`q^{1/2}`$, which we consider as a left $`U_qso(3)`$ module algebra. It is generated by $`x_i`$ fulfilling the relations
$`\epsilon _k^{ij}x_ix_j`$ $`=`$ $`\mathrm{\Lambda }_Nx_k,`$
$`g^{ij}x_ix_j`$ $`=`$ $`R^2.`$ (6.1)
Here $`R>0`$,
$$C_N=\frac{[N]_q[N+2]_q}{[2]_q^2},\mathrm{\Lambda }_N=R\frac{[2]_{q^{N+1}}}{\sqrt{[N]_q[N+2]_q}}$$
(6.2)
where $`[n]_q:=\frac{q^{n/2}q^{n/2}}{q^{1/2}q^{1/2}}`$, and
$$\begin{array}{cc}\epsilon _1^{10}=q^{1/2},\hfill & \epsilon _1^{01}=q^{1/2},\hfill \\ \epsilon _0^{00}=q^{1/2}q^{1/2},\hfill & \epsilon _0^{11}=1=\epsilon _0^{11},\hfill \\ \epsilon _1^{01}=q^{1/2},\hfill & \epsilon _1^{10}=q^{1/2}\hfill \end{array}$$
(6.3)
are the spin 1 Clebsch–Gordan coefficients. The multiplet $`(x_i)`$ carries the fundamental vector representation $`\rho `$ of $`H=U_qso(3)`$:
$$gx_i=x_j\rho _i^j(g).$$
(6.4)
There is no obvious generalization to higher dimensions, but this algebra appears to be relevant e.g. to $`D`$–branes on the $`SU(2)`$ WZW model . It has a unique irreducible representation, which is equivalent to $`Mat(N+1)`$. Here we only consider the case $`q^+`$, where the star structure is given by $`x_i^{}=g^{ij}x_j`$. Then $`\widehat{S}_{q,N}^2`$ is simply the “discrete series” of Podles’s spheres . It was shown in that there is a star–algebra homomorphism $`\widehat{\phi }:H<\widehat{S}_{q,N}^2\widehat{S}_{q,N}^2`$, which takes a particularly simple form
$`\widehat{\phi }(E^+)`$ $`=`$ $`{\displaystyle \frac{1}{R}}\sqrt{q^1[2]_qC_N}x_1,\widehat{\phi }(E^{})={\displaystyle \frac{1}{R}}\sqrt{q[2]_qC_N}x_1,`$
$`\widehat{\phi }(q^{H/2})`$ $`=`$ $`{\displaystyle \frac{[2]_{q^{N+1}}}{[2]_q}}\left(1{\displaystyle \frac{q^{1/2}q^{1/2}}{\mathrm{\Lambda }_N}}x_0\right)`$ (6.5)
where $`E^\pm =X^\pm q^{H/4}U_qso(3)`$. Note that $`(1\frac{q^{1/2}q^{1/2}}{\mathrm{\Lambda }_N}x_0)`$ is invertible since the eigenvalues of $`q^{H/2}`$ are positive (assuming $`q>0`$), therefore $`\widehat{\phi }(q^{H/2})\widehat{S}_{q,N}^2`$ is well–defined also. Hence the algebra homomorphisms $`\widehat{\phi }`$ is defined on the entire algebra $`U_qso(3)`$. Using the definition (A.1.15) and the explicit form for the universal $``$ (see e.g. ), one finds
$$[^{}{}_{j}{}^{i}]=\left[\begin{array}{ccc}q^{H/2},& 0,& 0\\ (1q^1)\sqrt{[2]_q}E^{},& 1,& 0\\ q^{1/2}(1q^1)^2q^{H/2}(E^{})^2,& (1q^1)\sqrt{[2]_q}q^{H/2}E^{},& q^{H/2}\end{array}\right]$$
(6.6)
and
$$[^+{}_{j}{}^{i}]=\left[\begin{array}{ccc}q^{H/2},& (q1)\sqrt{[2]_q}q^{H/2}E^+,& (q1)^2q^{H/2}(E^+)^2,\\ 0,& 1,& q^{1/2}(q1)\sqrt{[2]_q}E^+\\ 0,& 0,& q^{H/2}\end{array}\right].$$
(6.7)
The unbraiding procedure then works as in Theorem 2. To be specific, assume that the braided tensor product algebra is as in (5.3). Then we set
$`y_{1,i}:=x_{1,i}`$ (6.8)
$`y_{\alpha ,i}:=\widehat{\chi }(x_{\alpha ,i})=x_{\alpha ,j}\widehat{\phi }_1(^+{}_{i}{}^{j}),\alpha >1,`$ (6.9)
without spelling out these expressions further. According to Theorem 2, they satisfy
###### Corollary 3
$`\epsilon _k^{ij}y_{\alpha ,i}y_{\alpha ,j}=\mathrm{\Lambda }_Ny_{\alpha ,k},`$
$`g^{ij}y_{\alpha ,i}y_{\alpha ,j}=R^2`$
for all $`\alpha =1,\mathrm{},M`$, together with
$`[y_{1,i},y_{\alpha ,j}]=0`$ (6.10)
$`y_{\alpha ,i}y_{\beta ,j}=\widehat{R}_{ij}^{hk}y_{\beta ,h}y_{\alpha ,k}`$ (6.11)
when $`1<\alpha `$ and $`\alpha \beta `$.
Iterating this procedure as before, we find
###### Proposition 5
Let $`𝒜_1,𝒜_2,\mathrm{},𝒜_M`$ be $`M`$ copies of the $`U_qso(3)`$–covariant fuzzy quantum sphere. Then $`𝒜_1\underset{¯}{}^\pm 𝒜_2\underset{¯}{}^\pm \mathrm{}\underset{¯}{}^\pm 𝒜_M=𝒜_1\stackrel{~}{𝒜}_2^\pm \mathrm{}\stackrel{~}{𝒜}_M^\pm `$, where $`\stackrel{~}{𝒜}_2^\pm ,\mathrm{},\stackrel{~}{𝒜}_M^\pm `$ are subalgebras of the lhs isomorphic to $`𝒜_1`$.
## A Appendix
### A.1 The universal $`R`$-matrix
In this appendix we recall the basics about the universal $`R`$-matrix of the quantum groups $`U_q\text{g}`$ , while fixing our conventions. Recall that the universal $`R`$-matrix $``$ is a special element
$$^{(1)}^{(2)}U_q\text{g}\text{ }U_q\text{g}\text{ }$$
(A.1.1)
intertwining between $`\mathrm{\Delta }`$ and opposite coproduct $`\mathrm{\Delta }^{op}`$, and so does also $`_{21}^1`$:
$$\begin{array}{c}(g_{(1)}g_{(2)})=(g_{(2)}g_{(1)}),\hfill \\ _{21}^1(g_{(1)}g_{(2)})=(g_{(2)}g_{(1)})_{21}^1.\hfill \end{array}$$
(A.1.2)
In (A.1.1) we have used a Sweedler notation with upper indices: the right-hand side is a short-hand notation for a sum $`_I_I^{(1)}_I^{(2)}`$ of infinitely many terms. We recall some useful formulae
$`(\mathrm{\Delta }\text{id})=_{13}_{23},`$ (A.1.3)
$`(\text{id}\mathrm{\Delta })=_{13}_{12},`$ (A.1.4)
$`(S\text{id})=^1=(\text{id}S^1),`$ (A.1.5)
$`S^1(g)=u^1S(g)u.`$ (A.1.6)
Here $`u`$ is any of the elements $`u_1,u_2,..u_8`$ defined below:
$$\begin{array}{cc}u_1:=(S^{(2)})^{(1)}\hfill & u_2:=(S^1{}_{}{}^{(1)})^1^{(2)}\hfill \\ u_3:=^{(2)}S^1^{(1)}\hfill & u_4:=^1{}_{}{}^{(1)}S_{}^{1}^1^{(2)}\hfill \\ (u_5)^1:=^{(1)}S^{(2)}\hfill & (u_6)^1:=(S^1^{(1)})^{(2)}\hfill \\ (u_7)^1:=^1{}_{}{}^{(2)}S^1^{(1)}\hfill & (u_8)^1:=(S^1^1{}_{}{}^{(2)})^1^{(1)}\hfill \end{array}$$
(A.1.7)
In fact, using the results of Drinfel’d one can show that
$$u_1=u_3=u_7=u_8=vu_2=vu_4=vu_5=vu_6,$$
(A.1.8)
where $`v`$ is a suitable element belonging to the center of $`U_q\text{so}(N)`$.
From (A.1.2) and (A.1.3,A.1.4) it follows the universal Yang-Baxter relation
$$_{12}_{13}_{23}=_{23}_{13}_{12},$$
(A.1.9)
whence the other two relations follow
$`^1{}_{12}{}^{}_{}^{1}{}_{13}{}^{}_{}^{1}_{23}`$ $`=`$ $`^1{}_{23}{}^{}_{}^{1}{}_{13}{}^{}_{}^{1}{}_{12}{}^{},`$ (A.1.10)
$`_{13}_{23}^1_{12}`$ $`=`$ $`^1{}_{12}{}^{}_{23}^{}_{13}.`$ (A.1.11)
As before, let $`\rho `$ be the fundamental $`N`$-dimensional representation of $`\text{}=sl(N),so(N),sp(N)`$ By applying $`\text{id}\rho _c^a\rho _d^b`$ to (A.1.9), $`\rho _c^a\rho _d^b\text{id}`$ to (A.1.10) and $`\rho _c^a\text{id}\rho _d^b`$ to (A.1.11) we respectively find the commutation relations
$`\widehat{R}_{cd}^{ab}^+{}_{f}{}^{d}_{}^{+}{}_{e}{}^{c}=^+{}_{c}{}^{b}_{}^{+}{}_{d}{}^{a}\widehat{R}_{ef}^{dc},`$ (A.1.12)
$`\widehat{R}_{cd}^{ab}^{}{}_{f}{}^{d}_{}^{}{}_{e}{}^{c}=^{}{}_{c}{}^{b}_{}^{}{}_{d}{}^{a}\widehat{R}_{ef}^{dc},`$ (A.1.13)
$`\widehat{R}_{cd}^{ab}^+{}_{f}{}^{d}_{}^{}{}_{e}{}^{c}=^{}{}_{c}{}^{b}_{}^{+}{}_{d}{}^{a}\widehat{R}_{ef}^{dc},`$ (A.1.14)
where $`^\pm _l^a`$ are the Faddeev-Reshetikin-Takhtadjan generators of $`U_q\text{g}`$ , defined by
$$^+{}_{l}{}^{a}:=^{(1)}\rho _l^a(^{(2)})^{}{}_{l}{}^{a}:=\rho _l^a(^1{}_{}{}^{(1)})^1{}_{}{}^{(2)}.$$
(A.1.15)
It is known that $`\{^+{}_{j}{}^{i},^{}{}_{j}{}^{i}\}`$ and the square roots of the elements $`^\pm _i^i`$ provide a (overcomplete) set of generators of $`U_q\text{g}`$ . Since in our conventions
$$H^+H^{},$$
(A.1.16)
then $`^+{}_{l}{}^{a}H^+`$ and $`^{}{}_{l}{}^{a}H^{}`$. Beside (A.1.12-A.1.14) these generators fulfill
$`^+{}_{j}{}^{i}=0,\text{if }i>j`$ (A.1.17)
$`^{}{}_{j}{}^{i}=0,\text{if }i<j`$ (A.1.18)
$`^{}{}_{i}{}^{i}_{}^{+}{}_{i}{}^{i}=^+{}_{i}{}^{i}_{}^{}{}_{i}{}^{i}=1,i`$ (A.1.19)
$`^\pm {}_{n}{}^{n}\mathrm{}.^\pm {}_{n}{}^{n}=1,`$ (A.1.20)
and, when $`\text{}=so(N),sp(N)`$, some additional relations. When $`\text{}=so(N)`$ the latter read
$$^\pm {}_{j}{}^{i}_{}^{\pm }{}_{k}{}^{h}g_{}^{kj}=g^{hi}^\pm {}_{i}{}^{j}_{}^{\pm }{}_{h}{}^{k}g_{kj}^{}=g_{hi},$$
(A.1.21)
where $`g_{ij}`$ has been defined in (3.5). The braid matrix $`\widehat{R}`$ is related to $``$ by $`\widehat{R}_{hk}^{ij}R_{hk}^{ji}:=(\rho _h^j\rho _k^i)`$. With the indices’ convention described in sections 3, 4 $`\widehat{R}`$ is given by
$$\widehat{R}=q^{\frac{1}{N}}\left[q\underset{i}{}e_i^ie_i^i+\underset{ij}{}e_i^je_j^i+k\underset{i<j}{}e_i^ie_j^j\right]$$
(A.1.22)
when $`\text{}=sl(N)`$, and by
$`\widehat{R}`$ $`=`$ $`q{\displaystyle \underset{i0}{}}e_i^ie_i^i+{\displaystyle \underset{\stackrel{ij,j}{\text{ or }i=j=0}}{}}e_i^je_j^i+q^1{\displaystyle \underset{i0}{}}e_i^ie_i^i`$
$`+k({\displaystyle \underset{i<j}{}}e_i^ie_j^j{\displaystyle \underset{i<j}{}}q^{\rho _i+\rho _j}e_i^je_i^j)`$
when $`\text{}=so(N)`$. Here $`e_j^i`$ is the $`N\times N`$ matrix with all elements equal to zero except for a $`1`$ in the $`i`$th column and $`j`$th row. The braid matrix of $`sl(N)`$ admits the orthogonal projector decomposition
$$\widehat{R}=q𝒫_Sq^1𝒫_a,\text{}=sl(N);$$
(A.1.24)
$`𝒫_a,𝒫_S`$ are the $`U_qsl(N)`$-covariant deformed antisymmetric and symmetric projectors. The braid matrix of $`so(N)`$ admits the orthogonal projector decomposition
$$\widehat{R}=q𝒫_sq^1𝒫_a+q^{1N}𝒫_t\text{}=so(N);$$
(A.1.25)
$`𝒫_a,𝒫_t,𝒫_s`$ are the $`q`$-deformed antisymmetric, trace, trace-free symmetric projectors.
The compact section of $`U_q\text{g}`$ requires $`q^+`$ if $`\text{}=so(N)`$, $`q`$ if $`\text{}=sl(N)`$ and is characterized by the $``$-structure
$$(^\pm {}_{j}{}^{i})^{}=S^{}{}_{i}{}^{j}.$$
(A.1.26)
For $`\text{}=so(N)`$ this amounts to
$$(^\pm {}_{j}{}^{i})^{}=g_{ih}^{}{}_{k}{}^{h}g_{}^{kj}.$$
(A.1.27)
The non-compact sections of $`U_q\text{g}`$ require $`|q|=1`$ and are characterized by the $``$-structure
$$(^\pm {}_{j}{}^{i})^{}=U^1{}_{r}{}^{i}_{}^{\pm }{}_{s}{}^{r}U_{j}^{s}=u^\pm {}_{j}{}^{i}u_{}^{1}.$$
(A.1.28)
This can be checked using the property $`(\widehat{R}_{hk}^{ij})^{}=\widehat{R}^1_{kh}^{ji}`$. Here we have defined
$$U_j^i=\rho _j^i(u)$$
(A.1.29)
with $`u`$ any of the elements defined in (A.1.7). For $`\text{}=so(N)`$ one can take
$$U_j^i:=g^{ih}g_{jh}.$$
(A.1.30)
From formulae (A.1.3), (A.1.4) in the Appendix A.1 one finds that the coproducts are given by
$$\mathrm{\Delta }(^+{}_{j}{}^{i})=^+{}_{h}{}^{i}^+{}_{j}{}^{h}\mathrm{\Delta }(^{}{}_{j}{}^{i})=^{}{}_{h}{}^{i}^{}{}_{j}{}^{h}.$$
(A.1.31)
### A.2 Proof of Proposition 1
We make use of the identity
$$\phi ^\pm (g^\pm )h^\pm =\phi ^\pm (g^\pm h^\pm ),$$
(A.2.1)
for any $`g^\pm ,h^\pm H^\pm `$, which we prove in Ref. . The right action appearing at the rhs is the (right) adjoint action on itself
$$hg=Sg_{(1)}hg_{(2)},g,hH;$$
(A.2.2)
where $`S`$ denotes the antipode of the Hopf algebra $`H`$. We shall also need the inverse of (1.9),
$$a_1a_2=(a_2^{1(2)})(a_1^{1(1)}).$$
(A.2.3)
Now,
$`\chi ^+(a_2)`$ $`\stackrel{(\text{1.13})}{=}`$ $`\phi _1^+(^{(1)})(a_2^{(2)})`$
$`\stackrel{(\text{A.2.3})}{=}`$ $`(a_2^{(2)}^{1(2^{})})[\phi _1^+(^{(1)})^{1(1^{})}]`$
$`\stackrel{(\text{A.2.1})}{=}`$ $`(a_2^{(2)}^{1(2^{})})\phi _1^+(^{(1)}^{1(1^{})})`$
$`\stackrel{(\text{A.2.2})}{=}`$ $`(a_2^{(2)}^{1(2^{})})\phi _1^+(S_{(1)}^{1(1^{})}^{(1)}_{(2)}^{1(1^{})})`$
$`\stackrel{(\text{A.1.3})}{=}`$ $`(a_2^{(2)}^{1(2^{})}^{1(2\mathrm{"})})\phi _1^+(S^{1(1\mathrm{"})}^{(1)}^{1(1^{})})`$
$`=`$ $`(a_2^{1(2\mathrm{"})})\phi _1^+(S^{1(1\mathrm{"})}),`$
which proves (1.21). Similarly one proves (1.22).
### A.3 The maps $`\phi ^\pm `$ for the quantum Euclidean spaces or spheres
We introduce the short-hand notation
$$[A,B]_x=ABxBA.$$
(A.3.1)
In Ref. we have found algebra homomorphisms $`\phi ^\pm :_q^N>U_q^\pm so(N)_q^N`$. The images of $`\phi ^{}`$ on the negative FRT generators read
$$\phi ^{}(^{}{}_{j}{}^{i})=g^{ih}[\mu _h,x^k]_qg_{kj},$$
(A.3.2)
where
$$\begin{array}{cc}\mu _0=\gamma _0(x^0)^1\hfill & \text{for }N\text{ odd,}\hfill \\ \mu _{\pm 1}=\gamma _{\pm 1}(x^{\pm 1})^1^\pm _1^1\hfill & \text{for }N\text{ even,}\hfill \\ \mu _a=\gamma _ar_{|a|}^1r_{|a|1}^1x^a\hfill & \text{otherwise,}\hfill \end{array}$$
(A.3.3)
and $`\gamma _a`$ are normalization constants fulfilling the conditions
$$\begin{array}{cc}\gamma _0=q^{\frac{1}{2}}h^1\hfill & \text{for }N\text{ odd,}\hfill \\ \gamma _1\gamma _1=\{\begin{array}{c}q^1h^2\hfill \\ k^2\hfill \end{array}\hfill & \begin{array}{c}\text{for }N\text{ odd,}\hfill \\ \text{for }N\text{ even,}\hfill \end{array}\hfill \\ \gamma _a\gamma _a=q^1k^2\omega _a\omega _{a1}\hfill & \text{for }a>1.\hfill \end{array}$$
$`h,k,\omega _a`$ are defined as in Sections 3, 4. On the other hand, the images of $`\phi ^+`$ on the positive FRT generators read
$$\phi ^+(^+{}_{j}{}^{i})=g^{ih}[\overline{\mu }_h,x^k]_{q^1}g_{kj},$$
(A.3.4)
where
$$\begin{array}{cc}\overline{\mu }_0=\overline{\gamma }_0(x^0)^1\hfill & \text{for }N\text{ odd,}\hfill \\ \overline{\mu }_{\pm 1}=\overline{\gamma }_{\pm 1}(x^{\pm 1})^1^{}_1^1\hfill & \text{for }N\text{ even,}\hfill \\ \overline{\mu }_a=\overline{\gamma }_ar_{|a|}^1r_{|a|1}^1x^a\hfill & \text{otherwise,}\hfill \end{array}$$
(A.3.5)
and $`\overline{\gamma }_a`$ normalization constants fulfilling the conditions
$$\begin{array}{cc}\overline{\gamma }_0=q^{\frac{1}{2}}h^1\hfill & \text{for }N\text{ odd,}\hfill \\ \overline{\gamma }_1\overline{\gamma }_1=\{\begin{array}{c}qh^2\hfill \\ k^2\hfill \end{array}\hfill & \begin{array}{c}\text{for }N\text{ odd,}\hfill \\ \text{for }N\text{ even,}\hfill \end{array}\hfill \\ \overline{\gamma }_a\overline{\gamma }_a=qk^2\omega _a\omega _{a1}\hfill & \text{for }a>1.\hfill \end{array}$$
Incidentally, for odd $`N`$ one can choose the free parameters $`\gamma _a,\overline{\gamma }_a`$ in such a way that $`\phi ^+,\phi ^{}`$ can be ‘glued’ into an algebra homomorphism $`\phi :_q^N>U_q\text{so}(N)_q^N`$ .
We give the explicit expression for $`\phi ^\pm (^\pm {}_{j}{}^{i})`$ in the case $`N=3`$:
$$[\phi ^{}(^{}{}_{j}{}^{i})]=\left[\begin{array}{ccc}qh\gamma _1(x^0)^1r& & \\ q^{\frac{1}{2}}(q+1)(x^0)^1x^+& 1& \\ q^{\frac{1}{2}}(q+1)(h\gamma _1rx^0)^1(x^+)^2& (1+q^1)(h\gamma _1r)^1& (qh\gamma _1r)^1x^0\end{array}\right]$$
(A.3.6)
and
$$[\phi ^+(^+{}_{j}{}^{i})]=\left[\begin{array}{ccc}h\overline{\gamma }_1r^1x^0& q^{\frac{1}{2}}\overline{\gamma }_1kr^1x^{}& q^2k\overline{\gamma }_1(rx^0)^1(x^{})^2\\ & 1& q^{\frac{1}{2}}(q^1+1)(x^0)^1x^{}\\ & & (h\overline{\gamma }_1x^0)^1r\end{array}\right]$$
(A.3.7)
When $`q^+`$ the real structure of $`_q^N`$ is given by
$$(x^i)^{}=x^jg_{ji}.$$
(A.3.8)
Note that when $`N`$ is odd $`\mu _0,\overline{\mu }_0`$, which are completely determined by their definitions, are such that $`\mu _0^{}=q^1\overline{\mu }_0`$. We fix the other $`\gamma _a,\overline{\gamma }_a`$ so that for any $`a`$
$$\mu _a^{}=q^1g_{ab}\overline{\mu }_b.$$
(A.3.9)
This was already considered in Ref. and requires
$$\begin{array}{cc}\gamma _{\pm 1}^{}=\overline{\gamma }_1\hfill & \text{if }N\text{ even}\hfill \\ \gamma _a^{}=\overline{\gamma }_a\{\begin{array}{cc}\text{1 if }a<0\hfill & \\ q^2\text{ if }a>0\hfill & \end{array}\hfill & \text{otherwise.}\hfill \end{array}$$
(A.3.10)
As a consequence,
$`\left[\phi ^{}(^{}{}_{j}{}^{i})\right]^{}`$ $`\stackrel{(\text{A.3.2})}{=}`$ $`\left(g^{ih}[\mu _h,x^k]_qg_{kj}\right)^{}`$
$`\stackrel{(\text{A.3.8})}{=}`$ $`g^{ih}[x^j,\mu _h^{}]_q`$
$`\stackrel{(\text{A.3.9})}{=}`$ $`[\overline{\mu }_i,x^j]_{q^1}`$
$`\stackrel{(\text{A.3.4})}{=}`$ $`g_{ih}\phi ^+(^+{}_{k}{}^{h})g^{kj}`$
$`\stackrel{(\text{A.1.27})}{=}`$ $`\phi ^+\left[(^{}{}_{j}{}^{i})^{}\right]`$
In other words
$$[\phi ^\pm (g)]^{}=\phi ^{}(g^{}).$$
(A.3.11)
When $`|q|=1`$
$$(x^i)^{}=x^i$$
(A.3.12)
Note that when $`N`$ is odd $`\mu _0,\overline{\mu }_0`$, which are completely determined by their definitions, are such that $`\mu _0^{}=q\mu _0=\overline{\mu }_0`$. We fix the other $`\gamma _a,\overline{\gamma }_a`$ so that for any $`a`$
$$\mu _a^{}=q\mu _a,\overline{\mu }_a^{}=q^1\overline{\mu }_a.$$
(A.3.13)
This requires
$$\begin{array}{cc}\gamma _{\pm 1}^{}=\gamma _{\pm 1}\hfill & \text{if }N\text{ even}\hfill \\ \gamma _a^{}=\gamma _a\{\begin{array}{cc}\text{1 if }a<0\hfill & \\ q^2\text{ if }a>0\hfill & \end{array}\hfill & \text{otherwise.}\hfill \end{array}$$
(A.3.14)
As a consequence,
$`\left[\phi ^{}(^{}{}_{j}{}^{i})\right]^{}`$ $`\stackrel{(\text{A.3.2})}{=}`$ $`\left(g^{ih}[\mu _h,x^k]_qg_{kj}\right)^{}`$
$`\stackrel{(\text{A.3.12})}{=}`$ $`q^1g^{hi}[\mu _h^{},x^k]_qg_{jk}`$
$`\stackrel{(\text{A.3.13})}{=}`$ $`g^{hi}[\mu _h,x^k]_qg_{jk}`$
$`\stackrel{(\text{A.1.30}),(\text{A.3.2})}{=}`$ $`U^1{}_{r}{}^{i}\phi _{}^{}(^{}{}_{s}{}^{r})U_r^i`$
$`\stackrel{(\text{A.1.28})}{=}`$ $`\phi ^{}\left[(^{}{}_{j}{}^{i})^{}\right].`$
Similarly one proves that $`[\phi ^{}(^{}{}_{j}{}^{i})]^{}=\phi ^{}[(^{}{}_{j}{}^{i})^{}]`$. In other words, $`\phi ^\pm `$ are $``$-homomorphisms.
### A.4 The maps $`\phi `$ for the deformed Heisenberg algebras
In Ref. we constructed an algebra homomorphism $`\phi :U_q\text{so}(N)<𝒜_1𝒜_1`$, where $`𝒜_1`$ denotes the $`U_q\text{so}(N)`$-covariant (deformed) Heisenberg algebra, such that $`\phi `$ is a $``$-homomorphism
$$\phi (g^{})=\phi (g)^{}$$
(A.4.1)
on the compact section of $`U_q\text{so}(N)`$(what requires $`q^+`$). One can easily prove the same result also for the noncompact section (A.1.28) of $`\text{}=so(N)`$ as well as the compact and noncompact sections of $`\text{}=sl(N)`$. This can be done maybe most rapidly using as a set of generators the socalled ”vector fields” $`Z_j^i`$ , which are related to the FRT generators by
$$Z_j^i=^+{}_{h}{}^{i}S^{}{}_{j}{}^{h}.$$
(A.4.2)
From (A.1.26), (A.1.28) one immediately finds
$`(Z_j^i)^{}=Z_i^j\text{if }q^+`$ (A.4.3)
$`(Z_j^i)^{}=U^1{}_{a}{}^{i}(S^1^{}{}_{b}{}^{h})^+{}_{h}{}^{a}U_{j}^{b}\text{if }|q|=1;`$ (A.4.4)
if $`\text{}=so(N)`$ the second relation reduces to
$$(Z_j^i)^{}=U^1{}_{b}{}^{a}Z_{c}^{b}\widehat{R}^1{}_{aj}{}^{ci}.$$
(A.4.5)
In Ref. the explicit expression of $`\phi (Z_j^i)`$ in terms of the $`x`$’s and $``$’s is given both for $`g=sl(N)`$ and $`\text{}=so(N)`$, and it is not difficult to show that on these generators (and therefore on all of $`U_q\text{g}`$ ) (A.4.1) is satisfied. In performing the calculations one has to keep in mind that the authors of Ref. work with the left action, rather than with the right, so one has to switch to the conventions described in section 5, but, as explained there, this wil not modify the result (A.4.1). As an intermediate step, we give the action of the $``$-structure on the coordinates and derivatives for the case $`\text{}=so(N)`$, in the notation used there:
$`(x_h)^{}=g^{hk}x_k,(^i)^{}=q^N\widehat{}_i\text{if }q^+`$ (A.4.6)
$`(x_h)^{}=x_h,(^i)^{}=q^NU^1{}_{j}{}^{i}_{}^{j},(\widehat{}_i)^{}=q^N_i\text{if }|q|=1.`$ (A.4.7)
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# Spontaneous magnetisation in the plane
##
The Arak process is a solvable stochastic process which generates coloured patterns in the plane. Patterns are made up of a variable number of random non-intersecting polygons. We show that the distribution of Arak process states is the Gibbs distribution of its states in thermodynamic equilibrium in the grand canonical ensemble. The sequence of Gibbs distributions form a new model parameterised by temperature. We prove that there is a phase transition in this model, for some non-zero temperature. We illustrate this conclusion with simulation results. We measure the critical exponents of this off-lattice model and find they are consistent with those of the Ising model in two dimensions.
## Figure Captions
(A) A state $`\chi `$ of the Arak process (B) The discontinuity set $`\gamma `$ of (A).
(A) A set of lines $`\mathrm{}`$ intersecting $`𝒟`$ (B) an admissible graph drawn on the set $`\mathrm{}`$ (C) one of the two colourings of $`𝒟`$ with discontinuity set given by the graph in (B).
Updates in the Markov Chain Monte Carlo. Dashed and solid edges are exchanged by the moves, which are reversible. (A) Interior vertex birth and death (B) move a vertex, and (C) recolour a region by swapping a pair of edges. In an extra move, not shown, a small triangle may be created or deleted. Further move types are used to update boundary structures.
Binder parameter $`U_d`$ (see text), regressed with cubic polynomials. Curves correspond to distinct box-side lengths $`d`$. The maximum likelihood fit, constrained to intersect at a point, is shown. Error bars in this and all other graphs are $`1\sigma `$.
The Binder parameter data of Figure 4 rescaled with Ising critical exponents. The regression is a cubic polynomial. $`\chi _{434}^2=38.5`$ for the fit is acceptable.
The magnetisation $`\overline{m}_d(T)`$, regressed with cubic polynomials.
The magnetisation data of Figure 6 rescaled with Ising critical exponents. The regression is a quartic polynomial. The value of the $`\chi ^2`$ statistic shows that the fit is a poor one.
A selection of states equilibrated in a box of side $`d=12`$ at temperatures below and above the estimated critical temperature $`T_c0.6665(5)`$.
## 1 Introduction
The Widom-Rowlinson model, with two species of discs and hard-core interactions between discs of unlike species, is sometimes referred to as the “continuum Ising model”. However there is another continuum model which might share the title. In 1982 Arak presented a stochastic process in the plane with realisations of the kind shown in Figure 1A. States are composed of a variable number of coloured non-intersecting random polygons. Remarkably, the normalising constant is available as an explicit function of the area and boundary length of the region in which the process is realised. We present rigorous results and simulation based measurements related to critical phenomena in a two dimensional “continuum Ising model” derived from the Arak process.
There are few rigorous results for continuum models of interacting extended two dimensional objects. Moreover, relatively few Monte Carlo simulation studies have been made, perhaps on account of the complexity of the simulation algorithms required. The Widom-Rowlinson model has a phase transition . Its critical exponents have been measured and put it in the Ising universality class . Critical phenomena are known to occur in a range of related models with $`q2`$ species and certain soft-core interactions . Where critical exponents have been measured the universality class seems to be the class of the corresponding $`q`$species Potts model. For single-species models rigorous existence results for phase transitions have been given only in certain restricted models having area interactions .
In the model we consider the interface between black and white regions summarises the state in the same way that Peierls’ contours parameterise an Ising system. The energy associated with a state is proportional to the length of the interface. In contrast to the Ising model, the vertices of the polygon forming the interface take positions in the plane continuum. At a temperature $`T=1`$, the model we consider corresponds to the Arak process. For this value of the temperature the partition function equals the normalising constant of the corresponding Arak process. At smaller values of the temperature we are no longer dealing with an Arak process. We no longer have a closed form for the partition function. However the model remains well defined, and two phases coexist at temperatures bounded away from zero.
Besides this result, which we prove, we estimate the critical exponents of the temperature-modified Arak model, using Markov chain simulation to generate realisations of the process. Values (obtained by “data-collapsing”) are consistent with the corresponding critical exponents of the Ising model. This is in accord with what we expect from the hypothesis of universality, since the ground state of the temperature-modified Arak model is two-fold degenerate, and states are two dimensional.
Although there is no high temperature limit for polygonal models (a class of models including the Arak process) consistent polygonal models might play this role (this point is made in ). We give no rigorously determined upper bound on the critical temperature, although it is clear, from our simulations, that the consistent Arak process has a single phase.
## 2 The Arak process
We now define the Arak process, following . A state is a colouring map $`\chi :𝒟𝒥`$ from each point in an open convex set $`𝒟\mathrm{}^2`$, onto a set $`𝒥`$ of possible colours. See Figure 1A. We write $`𝒟`$ for the set of points in the boundary of $`𝒟`$. We consider the simplest case, $`𝒥=\{\mathrm{𝚋𝚕𝚊𝚌𝚔},\mathrm{𝚠𝚑𝚒𝚝𝚎}\}`$, of two colours.
Let $`X_𝒟`$ be the class of all finite subsets $`x`$ of $`𝒟𝒟`$. Let $`X_𝒟^{(n)}`$ ($`n1`$) be the set of point-sets $`x`$ composed of $`n`$ points, so that $`X_𝒟=_{n=0}^{\mathrm{}}X_𝒟^{(n)}`$, with $`X_𝒟^{(0)}=\{\varphi \}`$ the subset $`x=\varphi `$ containing no points. Let $`dx_i`$ be the element of area in $`𝒟`$ and length on $`𝒟`$. A measure $`d\nu (x)`$ is defined on $`X_𝒟`$ by
$$d\nu (x)=dx_1dx_2\mathrm{}dx_n.$$
This is the measure of an independent pair of Poisson point processes of unit intensity, on the boundary and interior.
Let $`\mathrm{\Gamma }_𝒟(x)`$ be the set of all “polygon graphs” $`\gamma `$ which can be drawn on the point-set $`x`$, $`ie`$ the set of all graphs which can be drawn in $`𝒟`$ with edges non-intersecting straight lines, with the points in $`x`$ as vertices. All interior vertices must have degree 2 (they are $`V`$-vertices), and all boundary vertices degree 1 ($`I`$-vertices). $`\gamma `$ is composed of a number of separate polygons which may be chopped off by the boundary. See Figure 1B.
The space of all allowed polygon graphs is the union over vertex sets $`x`$ of the polygon graphs of $`x`$:
$$\mathrm{\Gamma }_𝒟\underset{xX_𝒟}{}\mathrm{\Gamma }_𝒟(x).$$
We define a measure on $`\mathrm{\Gamma }_𝒟`$ by
$`d\lambda (\gamma )`$ $`=`$ $`\kappa (\gamma )d\nu (x(\gamma )),`$ (1)
$`\kappa (\gamma )`$ $`=`$ $`{\displaystyle \underset{<i,j>}{}}{\displaystyle \frac{1}{e_{ij}}}{\displaystyle \underset{i=1}{\overset{n}{}}}\mathrm{sin}(\psi _i),`$ (2)
for a pattern $`\gamma `$ with vertices at $`x(\gamma )=(x_1,x_2\mathrm{}x_n)`$. $`\psi _i`$ is the smaller angle at vertex $`i`$ in $`𝒟`$, and the smaller angle made with the boundary tangent at $`x_i`$ for vertices on $`𝒟`$. The product over $`<i,j>`$ runs over vertex pairs $`i,j`$ connected by an edge in $`\gamma `$, with $`e_{ij}=|x_ix_j|`$ the length of the edge between vertices $`i`$ and $`j`$. A counting measure is taken on the graphs of a fixed point set. The significance of $`\kappa `$ is sketched at the end of this section.
Arak’s probability measure on $`\mathrm{\Gamma }_𝒟`$ is
$$P_𝒟\{d\gamma \}=\frac{1}{𝒵_𝒟}\mathrm{exp}(2L(\gamma ))d\lambda (\gamma ),$$
(3)
with $`L(\gamma )`$ the summed length of all edges in $`\gamma `$, and $`𝒵_𝒟`$ a normalising constant. Remarkably, $`𝒵_𝒟`$ has a simple closed form (ie the model is solvable),
$$𝒵_𝒟=\mathrm{exp}(L(𝒟)+\pi A(𝒟)),$$
where $`L(𝒟)`$ and $`A(𝒟)`$ are respectively the perimeter length and area of $`𝒟`$. Certain expectation values have been calculated (see ). Some examples are given in Table 1.
A colouring map $`\chi :𝒟𝒥`$ is a function assigning a colour, black or white, to each point in $`𝒟`$. See Figure 1A. Let a colouring map $`\chi `$ be given and let $`B_\chi `$ be the set of points $`x𝒟`$ with a black point, $`ie`$ some $`y𝒟`$ such that $`\chi (y)=\mathrm{𝚋𝚕𝚊𝚌𝚔}`$, in every $`ϵ`$-neighbourhood. Let $`W_\chi `$ be similarly defined for white points. Let $`\gamma (\chi )=B_\chi W_\chi `$ denote the discontinuity set of this colouring. For each polygon graph $`\gamma `$ we consider two colouring maps $`\chi :𝒟𝒥`$ each having discontinuity set $`\gamma (\chi )=\gamma `$. The two distinct colourings of a given polygon graph are assigned equal probability, so the probability measure for colour maps is just $`P_𝒟\{d\gamma (\chi )\}/2`$.
The probability measure (3) has a number of beautiful properties, besides solvability. Striking are consistency and the Markov property. Consider an open region $`S`$ of $`𝒟`$ with $`ϵ`$-neighbourhood $`(S)_ϵ𝒟`$; the probability measure for events in $`S`$, given full information about $`\chi `$ on $`(S)_ϵ`$, is independent of any further information about the state in $`𝒟S`$. That is the Markov property. Next, let $`S`$ be an open convex region $`S𝒟`$ and let $`\gamma _S\mathrm{\Gamma }_S`$ denote the restriction of a state $`\gamma \mathrm{\Gamma }_𝒟`$ to $`S`$. The probability measure for events simulated in $`𝒟`$ from $`P_𝒟\{d\gamma \}`$ but observed in the subset $`S`$ is equal to $`P_S\{d\gamma \}`$, in other words $`P_S\{d\gamma \}=P_𝒟\{d\gamma _S\}`$. That is consistency. The Arak process shares these properties with a much larger family of probability measures called the consistent polygonal models. See for the general picture.
We will now explain in brief how $`\kappa (\gamma )`$ arises, following closely. Consider a number of straight lines drawn in the plane. Let $`l_i=(\rho _i,\varphi _i)`$ where $`\rho _i`$ is the perpendicular distance from the line to an origin and $`\varphi _i`$ is the angle the line makes to the $`x`$-axis. The parameter space of $`l_i`$ is $`L=[0,\mathrm{})\times [0,\pi )`$. Let $`L_D`$ be that subset of $`L`$ consisting of all lines intersecting $`𝒟`$. Let $`dl=d\rho d\varphi `$ be Lebesgue measure of $`L_D`$. Let $`L_D^n`$ be the set of all line sets $`\mathrm{}=\{l_1,l_2\mathrm{}l_n\}`$ made up of $`n`$ lines, each in $`L_D`$. In this parameterisation $`_𝒟=_nL_D^n`$ is the set of all sets of lines in the plane intersecting $`𝒟`$, and
$$d\stackrel{~}{\nu }(\mathrm{})=dl_1dl_2\mathrm{}dl_n$$
is the element of measure of a line process in $`𝒟`$, corresponding to a Poisson point process of unit intensity in $`L_D`$. Referring to Figure 2, we define an admissible graph on a line set $`\mathrm{}`$ to be a graph with edges coinciding with lines in $`\mathrm{}`$, such that each line in $`\mathrm{}`$ contributes a single closed segment of non-zero length to the graph. All interior vertices are $`V`$ vertices, all boundary vertices are $`I`$ vertices. The set of all admissible graphs which can be drawn on some line set in $`_𝒟`$ is identical to $`\mathrm{\Gamma }_𝒟`$. Let $`\gamma `$ be some legal graph drawn on the line set $`\mathrm{}`$. Define a measure $`d\stackrel{~}{\lambda }(\gamma )=d\stackrel{~}{\nu }(\mathrm{})`$ in $`\mathrm{\Gamma }_𝒟`$ using the line process as our base measure, and taking counting measure over the legal graphs of a line set. We now have two parameterisations of the graph: from its line set $`\mathrm{}`$, or from its vertex set $`x`$. The authors of have shown that $`d\stackrel{~}{\lambda }(\gamma )=d\lambda (\gamma )`$, $`ie`$, $`\kappa (\gamma )`$ arises as the Jacobian of the transformation between $`x`$ and $`\mathrm{}`$.
## 3 Properties of a temperature modified Arak process
We choose to modify the measure (3), and consequently loose solvability. Consider a system of non-overlapping polygonal chains of fluctuating number, length and vertex composition, confined to a planar region $`𝒟`$. The chains may be attached in some places to the boundary of $`𝒟`$. The state is described by a graph $`\gamma \mathrm{\Gamma }_𝒟`$. Micro-states are associated with elements of volume $`d\stackrel{~}{\lambda }(\gamma )`$ in $`\mathrm{\Gamma }_𝒟`$, so that in the Gibbs ensemble edge segments are isotropic in orientation (a rather unnatural choice). However, the Gibbs distribution $`\mathrm{Q}_𝒟\{d\gamma \}`$ of this system is just the Arak distribution above, modified by the addition of a temperature parameter, as we now show.
The Gibbs distribution $`\mathrm{Q}_𝒟\{d\gamma \}`$ has a density, $`g(\gamma )`$ say, with respect to $`d\stackrel{~}{\lambda }(\gamma )`$, the line measure. The Boltzmann entropy of the system is
$$H[g]=_{\mathrm{\Gamma }_𝒟}g(\gamma )\mathrm{ln}(g(\gamma ))𝑑\stackrel{~}{\lambda }(\gamma )$$
In the grand canonical ensemble, the energy and dimension of the system state fluctuate about fixed average values. We suppose that the state energy $`E(\gamma )`$ is given by the total length of the chains, $`E(\gamma )=cL(\gamma )`$, with $`c`$ a positive constant. The dimension of the vertex position vector $`x`$ is $`dim(x)=2n_i+n_b`$ with $`n_i`$ ($`n_b`$) the number of interior (boundary) vertices in $`\gamma `$. Maximising the entropy subject to constraints on the mean energy and mean dimension of the state, we obtain the distribution of systems of chains,
$$\mathrm{Q}_𝒟\{d\gamma \}=\frac{1}{𝒵_{𝒟}^{}{}_{T}{}^{}}\mathrm{exp}(cL(\gamma )/T)q^{n_e}d\lambda (\gamma ),$$
where $`T`$ and $`q`$ are Lagrange multipliers, and $`n_e`$ is the number of edges in $`\gamma `$ ($`n_e=n_i+n_b/2`$). Under the change of scale $`x_i\widehat{q}x_i`$, the measure transforms as $`d\lambda (\gamma )\widehat{q}^{n_e}d\lambda (\gamma )`$. We therefore set $`q=1`$ without loss of generality. Setting $`c=2`$ we obtain a “temperature-modified” Arak process
$$P_{T,𝒟}\{d\gamma \}=\frac{1}{𝒵_{𝒟}^{}{}_{T}{}^{}}\mathrm{exp}(2L(\gamma )/T)d\lambda (\gamma ).$$
(4)
The function $`2L(\gamma )/T`$ is a potential, ($`ie`$ $`𝒵_{𝒟}^{}{}_{T}{}^{}`$ is finite), at least when $`0T1`$, and, by Theorem 8.1 of , the temperature-modified measure keeps the spatial Markov property of the Arak measure.
Let $`\mu _D^B(T)`$ be the mean proportion of $`𝒟`$ coloured black (and $`\mu _D^W(T)`$ white),
$$\mu _D^B(T)=𝖤_{T,𝒟}\left\{A(B_\chi )/A(D)\right\}.$$
The magnetisation of a state
$$m(\chi )=|A(B_\chi )A(W_\chi )|/A(D)$$
measures the colour asymmetry in that state. In our simulations (reported below) we see a qualitatively Ising-like temperature dependence in the mean magnetisation. We prove, in an Appendix, that there is long range order ($`ie`$ phase coexistence) in magnetisation, at all sufficiently small temperatures. We have translated Griffiths’ version of Peierls’ proof of phase coexistence in the Ising lattice model to this continuum case.
Let $`\mu _D^{B|W}(T)`$ be the expected proportion of $`𝒟`$ coloured black given that the boundary is white, that is
$$\mu _D^{B|W}(T)=𝖤_{T,𝒟}\left\{A(B_\chi )/A(D)DB_\chi =\mathrm{}\right\}.$$
TheoremFor the temperature modified Arak process in an open convex region $`𝒟\mathrm{}^2`$ there exists a temperature $`T_{\mathrm{cold}}`$, $`0<T_{\mathrm{cold}}<1`$ and a constant $`a`$, $`a>0`$, such that
$$\mu _D^{B|W}(T_{\mathrm{cold}})\frac{1}{2}a$$
independent of the area $`A(𝒟)`$ of the region. Surgailis has shown that, for an open convex set $`S𝒟`$, the thermodynamic limit $`𝒟\mathrm{}^2`$ of $`P_{T,𝒟}\{d\gamma _S\}`$ exists, for a class of measures including $`P_{T,𝒟}\{d\gamma \}`$, for all temperatures below some small fixed positive value. With the theorem above,
$$\mu ^{B|W}(T)=\underset{𝒟\mathrm{}^2}{lim}\mu _D^{B|W}(T)$$
exists and satisfies $`\mu ^{B|W}(T_{\mathrm{cold}})<1/2a`$ for some $`a>0`$. Hence, there is phase coexistence at all temperatures $`T<T_{\mathrm{cold}}`$.
In fact it follows from the result stated in the Appendix that
$$\mu ^{B|W}(T)\frac{1}{4\pi ^2}\left(\frac{1}{z^3}+\frac{4}{z^2}+\frac{8}{z}\right),$$
(5)
where $`z=(1/(\pi T)1)`$. Sketching the function of $`T`$ on the right hand side of Equation (5), we see that $`T_{\mathrm{cold}}>0.18`$, though this bound is not at all sharp. Simulation (see below) shows that the model has a phase transition with critical temperature very close to $`T=2/3`$.
The proof of the theorem is in two parts. We are after an upper bound on the expected area coloured black. The area of black in a state with white boundaries is not more than the summed area of the polygons it contains, and is maximised when they are not nested. This observation leads to a simplified bound on the expected area coloured black, Equation (10). This first result is obtained by an obvious translation of the Griffiths calculation into the terms of a continuum process. In that case the next step, bounding the number ways a polygon can be drawn on a lattice of fixed size, using a fixed number of links, is straightforward. In the continuum, the analogous problem is to bound the volume of the parameter space of a polygon of fixed length, where volume is measured using $`\stackrel{~}{\lambda }`$, the line-based measure. The main difficulty lies in the fact that there are unbounded, but integrable, functions in the measure which arise, for example, when an edge length goes to zero; these would be absent if there were no polygon closure constraint; as a consequence the closure constraint may not be relaxed as simply as it is in the Ising case.
## 4 Simulation Results
The probability measure $`P_{T,𝒟}\{d\gamma \}`$ may be sampled using the Metropolis-Hastings algorithm, and Markov chain Monte Carlo. In our simulations we take $`𝒟`$ to be a square box of side length $`d`$. Note that the number of vertices is not fixed. Since the dimension of a state depends on the number of vertices in it, the Markov chain must make jumps, corresponding to vertex addition and deletion, between states of unequal dimension. Simulation algorithms of this kind are widely used in physical chemistry and statistics . Although there exist vertex birth and death moves sufficient for ergodicity, we allow a number of other moves in order to reduce the correlation time of the chain. See Figure 3. At each update we generate a candidate state $`\gamma ^{}`$, by selecting one of the moves, and applying it to a randomly selected part of the graph. The candidate state becomes the current state ($`ie`$ it is accepted) with a probability given by the Metropolis-Hastings prescription. Otherwise it is rejected and the current state is not changed. In this way a reversible Markov chain is simulated. The chain is ergodic, with equilibrium measure $`P_{T,𝒟}\{d\gamma \}`$. Full details of our algorithm, including explicit detailed balance calculations for all the Markov chain updates, are given in .
The sampling algorithm is quite complex, but because the model is solvable at $`T=1`$, it is possible to debug the code, by comparing a range of estimated expectations with predicted values. In Table 1 we present a selection of system statistics at $`T=1`$. Quantities in brackets are one standard deviation in the place of the last quoted digit. These data show how the simulation code was tested in comparisons using analytically derived expectation values. Let $`\widehat{f}`$ equal the average of some statistic $`f(\chi )`$ over an output sequence of length $`N`$, let $`\rho _f(t)`$ equal the normalised autocorrelation (or ACF) of $`f`$ at lag $`t`$ and, for $`M>0`$, let $`\tau _f=1+2_{t=1}^M\rho (t)`$ estimate the normalised autocorrelation time of $`f(\chi )`$ in the output, so that the variance of $`\widehat{f}`$ is estimated by $`\tau _f\mathrm{var}(f)/N`$. We used Geyer’s initial monotone indicator to determine $`M`$, the lag at which the ACF is truncated. The asymptotic variance $`\sigma _\rho ^2`$ of the ACF as $`t\mathrm{}`$ was estimated and used as a consistency check on each measurement: the estimated ACF should fall off to zero smoothly, and at large lag should stay within $`2\sigma _\rho `$ bounds of zero. As usual we cannot show the Markov simulation process has converged, but it is at least stationary.
Run parameters for the measurements at $`T<1`$ are summarised in Table 2. Autocorrelations reported are for $`T=0.66`$, near the critical temperature. We estimate the integrated autocorrelation time $`\tau _m`$ of the state magnetisation, along with its standard error and present these alongside the total run length. Referring to Table 2, the autocorrelation time is fitted within standard error by $`\tau _md^{4.6}`$. Our Metropolis Hastings algorithm is a local update algorithm and this places practical limits on the size of the largest system we can explore.
We now report our measurements of the mean magnetisation, $`\overline{m}_d(T)=𝖤_{T,d}\{m(\chi )\}`$, and the Binder parameter
$$U_d(T)=1\frac{𝖤_{T,d}\{m(\chi )^4\}}{3𝖤_{T,d}\{m(\chi )^2\}^2}.$$
Under the scaling hypothesis, the various curves $`U_d(T)`$ indexed by $`d`$ all intersect at a single $`T`$-value, the critical temperature , $`T=T_c`$ say. A Bayesian estimate $`\widehat{T}_c`$ may be given for the intersection point. Let $`\widehat{U}`$ denote the ordered set of independent $`U`$-measurements we made (43 in all), let $`T_U`$ denote the ordered set of $`T`$ values at which measurements were made, and let $`\widehat{\mathrm{\Sigma }}_U`$ denote the ordered set of estimated standard errors for the measurements in $`\widehat{U}`$. These data are represented by the error bars in Figure 4. Each measurement is an independent measurement. For each $`d=6,8,12,16`$, we model the unknown true curve $`U_d(T)`$ using a cubic
$$U_d^{}(T)=U^{}+(TT^{})\underset{p=0}{\overset{2}{}}a_p^{(d)}x^p.$$
The parameterisation constrains the regression in such a way that the four curves intersect at a point $`(T^{},U^{})`$. We simulate the joint posterior distribution of the random variables
$$a^{(6)},a^{(8)},a^{(12)},a^{(16)},T^{},U^{}|\widehat{U},T_U,\mathrm{\Sigma }_U,$$
conditioning the slope to be negative in the region containing the data, and conditioning the lines to intersect at a point, but otherwise taking an improper prior equal to a constant for all vectors of parameter values. Again MCMC simulation was used. The marginal posterior distribution of $`T^{}`$ is very nearly Gaussian. Our estimate of the critical temperature is then
$$\widehat{T}_c=0.6665(5).$$
The quoted standard error is the standard deviation of $`T^{}`$ in its marginal posterior distribution.
The Bayesian inference scheme used to estimate $`T_c`$ above is attractive for several reasons. Above all it quantifies the uncertainty in our estimate of $`T_c`$, taking full account of the complex constraints applying in the regression (though taking no account of possible errors due to violations of scaling). The sensitivity of the outcome to the orders of the regressing polynomials was explored. The chosen orders were the smallest that gave an acceptable likelihood. The posterior mode, which is the maximum likelihood estimate for $`T_c`$, on account of our flat prior, occurs at $`T^{}=0.6663`$. Metric factors weight the mass of probability in the posterior distribution only slightly away from the maximum of the likelihood.
Because the energy has a discrete two-fold symmetry, and states are two dimensional, we expect the model to lie in the universality class of the Ising model. Finite size scaling under the scaling hypothesis leads to a system size dependence of the form
$`\overline{m}_d(\tau )`$ $`=`$ $`d^{\beta /\nu }g(d^{1/\nu }\tau )`$
$`U_d(\tau )`$ $`=`$ $`f(d^{1/\nu }\tau )`$
with $`f`$ and $`g`$ unknown functions, $`\tau `$ the reduced temperature $`(T/T_c1)`$, and $`\beta `$ and $`\nu `$ critical exponents. If we plot $`U_d(\tau )`$ or $`d^{\beta /\nu }\overline{m}_d(\tau )`$ against $`d^{1/\nu }\tau `$, we expect to see no significant dependence on system size $`d`$ for $`\tau `$ near zero. Using the Ising critical exponents $`\nu =1`$ and $`\beta =0.125`$ and our estimate $`\widehat{T}_c`$ for the critical temperature, we show, in Figures 5 and 7, the maximum likelihood fit to the transformed data. The transformed $`U_d`$-data lies on a smooth curve. The transformed $`\widehat{m}_d`$-data does not give a satisfactory $`\chi ^2`$ (all of the misfit comes from points at $`T>T_c`$), but this is to be expected: we are seeing scaling violations (a satisfactory fit to a quartic can be obtained ($`\chi _{295}^2=30`$) by dropping points at large $`T`$ from the $`d=6`$ and $`d=8`$ data). If this is so, then the critical exponents of the Ising model the temperature dependent Arak process are equal at the precision of our simulation analysis.
Sample realisations from the model, taken at temperatures around the critical temperature are shown in Figure 8.
## Acknowledgements
It is a pleasure to thank Bruce Calvert (Mathematics, Auckland University) for his advice and ideas.
## Appendix: long range order
We now give the proof. Condition on a white boundary. There can be no boundary vertices. Let $`\mathrm{\Gamma }_{W𝒟}`$ be the subspace of $`\mathrm{\Gamma }_𝒟`$ of polygon graphs with no boundary vertices. Let $`\mathrm{\Theta }_𝒟`$ be the subspace of $`\mathrm{\Gamma }_{W𝒟}`$ of graphs made up of just one polygon. Each point in $`\mathrm{\Theta }_𝒟`$ corresponds to a single polygon, lying wholly in $`𝒟`$. We begin by proving the inequality Equation (10) below.
Among states built from a given set of polygons, with no edge connected to the boundary, the black area is largest when the polygons are arranged so that none are nested. It follows that the area of black in a state $`\chi `$ with a white boundary is less than or equal to the sum of the areas of all the polygons in that state. The area of a polygon $`\theta `$ of perimeter length $`L(\theta )`$ is smaller than the area of a circle with the same perimeter, so $`A(\theta )<L(\theta )^2/4\pi `$ and
$$A(B_\chi )\underset{\theta \gamma (\chi )}{}\frac{L(\theta )^2}{4\pi }.$$
(6)
We want to take expectations of either side of Equation (6) so we clear $`\gamma `$ from the domain of the sum, using
$$\underset{\theta \gamma (\chi )}{}f(\theta )_{\mathrm{\Theta }_𝒟}f(\theta )\delta (\theta \gamma (\chi ))𝑑\nu (x(\theta )).$$
$`\delta (\theta \gamma (\chi ))`$ puts a delta function at each point in $`\mathrm{\Theta }`$ corresponding to a polygon in $`\gamma `$. Each of these is a product of delta functions in $`𝒟`$ for the vertices of $`\theta `$ to coincide with those of $`\gamma `$, with an indicator function for the edge connections to coincide. $`x(\theta )`$ is the set of vertex coordinate variables of the polygon $`\theta `$.
Now take the expectation of $`A(B_\chi )/A(𝒟)`$ over patterns $`\chi `$ in $`\mathrm{\Gamma }_{W𝒟}`$. We have
$$\mu _𝒟^{B|W}_{\mathrm{\Theta }_𝒟}\frac{L(\theta )^2}{4\pi A(𝒟)}𝖤\{\delta (\theta \gamma (\chi ))|𝒟B_\chi =\mathrm{}\}𝑑\nu (x(\theta )).$$
The expectation of the delta function is by definition
$$𝖤\{\delta (\theta \gamma )|𝒟B_\chi =\mathrm{}\}=\frac{_{\mathrm{\Gamma }_{W𝒟}}\delta (\theta \gamma )\times e^{2L(\gamma )/T}𝑑\lambda (\gamma )}{_{\mathrm{\Gamma }_{W𝒟}}e^{2L(\gamma )/T}𝑑\lambda (\gamma )}.$$
(7)
Simplify the denominator by restricting the integral to those graphs to which the polygon $`\theta `$ could be added without intersecting an edge of a polygon already in place. That is, if
$$\mathrm{\Gamma }_{W𝒟}^\theta \{\gamma \mathrm{\Gamma }_{W𝒟}:\gamma \theta \}$$
is the set of polygon graphs containing the polygon $`\theta `$, then
$$\stackrel{~}{\mathrm{\Gamma }}_{W𝒟}^\theta \underset{\gamma \mathrm{\Gamma }_{W𝒟}^\theta }{}\{\gamma \theta \}$$
is the sub-domain of interest. We have
$$_{\mathrm{\Gamma }_{W𝒟}}e^{2L(\gamma )/T}𝑑\lambda (\gamma )_{\stackrel{~}{\mathrm{\Gamma }}_{W𝒟}^\theta }e^{2L(\gamma )/T}𝑑\lambda (\gamma )$$
(8)
We now turn to the numerator of Equation (7). Carrying out the integration over vertices in $`\theta `$ using the $`\delta `$-function,
$`{\displaystyle _{\mathrm{\Gamma }_{W𝒟}}}\delta (\theta \gamma )\times e^{2L(\gamma )/T}𝑑\lambda (\gamma )`$ $`=`$ $`{\displaystyle _{\mathrm{\Gamma }_{W𝒟}^\theta }}e^{2L(\gamma )/T}\kappa (\theta )𝑑\lambda (\gamma \theta )`$ (9)
$`=`$ $`\kappa (\theta )e^{2L(\theta )/T}{\displaystyle _{\stackrel{~}{\mathrm{\Gamma }}_{W𝒟}^\theta }}e^{2L(\gamma )/T}𝑑\lambda (\gamma ),`$
since $`\gamma `$ does not contain $`\theta `$ in the second line. Substituting with (8) and (9) in (7), and cancelling,
$$𝖤\{\delta (\theta \gamma )|𝒟B_\chi =\mathrm{}\}\kappa (\theta )e^{2L(\theta )/T},$$
and consequently,
$$\mu _D^{B|W}\frac{1}{4\pi A(𝒟)}_{\mathrm{\Theta }_𝒟}L(\theta )^2e^{2L(\theta )/T}𝑑\lambda (\theta ).$$
In close analogy with Griffiths’ proof, we obtain
$$\mu _D^{B|W}\frac{1}{4\pi A(𝒟)}_0^{\mathrm{}}b^2e^{2b/T}\left[_{\mathrm{\Theta }_𝒟}\delta (bL(\theta ))𝑑\lambda (\theta )\right]𝑑b$$
(10)
The integral over $`b`$ is an integral over polygon perimeter lengths. The problem is now to bound the integral over $`\mathrm{\Theta }_𝒟`$ without introducing more than one factor of $`A(𝒟)`$, or too rapidly growing a function of $`b`$. This is done by the following Lemma. Let $`\mathrm{\Theta }_𝒟^{(n)}`$ be the subset of $`\mathrm{\Theta }_𝒟`$ of polygons with $`n`$ vertices.
LemmaLet
$$J_n_{\mathrm{\Theta }_𝒟^{(n)}}\delta (bL(\theta ))𝑑\lambda (\theta )$$
(11)
so that
$$_{\mathrm{\Theta }_𝒟}\delta (bL(\theta ))𝑑\lambda (\theta )=\underset{n=3}{\overset{\mathrm{}}{}}J_n$$
in (10). Then
$$J_nA(𝒟)n^2(n1)\frac{(2\pi )^{n1}b^{n3}}{(n2)!},$$
and consequently
$$_{\mathrm{\Theta }_𝒟}\delta (bL(\theta ))𝑑\lambda (\theta )(2\pi )^2A(𝒟)(4+2\pi b)^2e^{2\pi b}.$$
(12)
Proof of the Lemma: Start with $`J_n`$ defined in Equation (11). Use a standard labelling with $`x_1`$ the variable corresponding to the vertex in $`\theta `$ with the smallest x-coordinate, (smallest y-coordinate in case of ties) and vertex number increasing clockwise around $`\theta `$. In the first step we break the polygon at $`x_1`$ to make a chain. Consider the set $`\stackrel{~}{\mathrm{\Theta }}_𝒟^{(n)}`$ of distinct non-intersecting chains $`\stackrel{~}{\theta }`$ of $`n`$ edges linking $`n+1`$ vertices, labeled with variables $`x_1`$ to $`x_{n+1}`$. All the vertices in a chain lie entirely to the right of the first vertex (or directly above). Polygons are chains, $`\mathrm{\Theta }_𝒟^{(n)}\stackrel{~}{\mathrm{\Theta }}_𝒟^{(n)}`$, since the first and last vertices in a chain may coincide. Transform variables from $`\{x_i\}_{i=1}^n`$ to $`\{x_1,\{\underset{¯}{e}_i\}_{i=1}^n\}`$, where $`\underset{¯}{e}_i`$ is a Cartesian vector with origin $`x_i`$ corresponding to the edge from the $`i`$’th to the $`(i+1)`$’th vertex. When we switch to integrating over chains, we constrain $`\underset{¯}{e}_1+\underset{¯}{e}_2+\mathrm{}+\underset{¯}{e}_n`$ to be zero, so that the polygon closes. Equation (11) becomes
$$J_n_{\stackrel{~}{\mathrm{\Theta }}_𝒟^{(n)}}\delta (bL(\stackrel{~}{\theta }))\delta ^{(2)}\left(\mathrm{\Sigma }_k\underset{¯}{e}_k\right)\frac{d\underset{¯}{e}_1d\underset{¯}{e}_2\mathrm{}d\underset{¯}{e}_n}{e_1e_2\mathrm{}e_n}𝑑x_1,$$
with $`e_iL(\underset{¯}{e}_i)`$ and using $`\mathrm{sin}(\psi _i)1`$.
The integrand is unbounded. We partition the space into regions, and impose the constraints $`b=L(\stackrel{~}{\theta })`$ and $`\underset{¯}{e}_1+\underset{¯}{e}_2+\mathrm{}+\underset{¯}{e}_n=0`$ by integration over different variables in each region. For any particular region, the variables eliminated by the constraints are chosen so that the integrand is bounded in that region.
Our second step then is to fix, by an integration in some $`d\underset{¯}{e}_i`$, the closure constraint. We will need to be able to bound below the length of at least one edge of the chain. So define
$`\stackrel{~}{\mathrm{\Theta }}_{𝒟,ϵ}^{(n)}`$ $`=`$ $`\{\stackrel{~}{\theta }\stackrel{~}{\mathrm{\Theta }}_𝒟^{(n)}|\mathrm{abs}(L(\stackrel{~}{\theta })b)ϵ\}`$
$`\stackrel{~}{\mathrm{\Theta }}_{𝒟,ϵ}^{(n,i)}`$ $`=`$ $`\{\stackrel{~}{\theta }\stackrel{~}{\mathrm{\Theta }}_𝒟^{(n)}|\mathrm{abs}(L(\stackrel{~}{\theta })b)<ϵ,e_i(bϵ)/n\},`$
with $`i\{1,2,\mathrm{},n\}`$, and $`ϵ`$ a small positive constant, $`0<ϵ<b`$, depending on $`b`$. Each chain in $`\stackrel{~}{\mathrm{\Theta }}_{𝒟,ϵ}^{(n,i)}`$ has the property that its $`i`$th edge has length at least $`(bϵ)/n`$. Any chain, with $`n`$ edges and a total length differing from $`b`$ by not more than $`ϵ`$, must have such an edge. The sets $`\stackrel{~}{\mathrm{\Theta }}_{𝒟,ϵ}^{(n,i)}`$, $`i=1,2\mathrm{}n`$ are not disjoint, but combine with $`\stackrel{~}{\mathrm{\Theta }}_{𝒟,ϵ}^{(n)}`$ to cover $`\stackrel{~}{\mathrm{\Theta }}_𝒟^{(n)}`$. Chains in $`\mathrm{\Theta }_{𝒟,ϵ}^{(n)}`$ will not contribute to the integral. It follows that
$`J_n`$ $``$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle _{\stackrel{~}{\mathrm{\Theta }}_{𝒟,ϵ}^{(n,i)}}}\delta (bL(\stackrel{~}{\theta }))\delta ^{(2)}\left(\mathrm{\Sigma }_k\underset{¯}{e}_k\right){\displaystyle \frac{d\underset{¯}{e}_1d\underset{¯}{e}_2\mathrm{}d\underset{¯}{e}_n}{e_1e_2\mathrm{}e_n}}𝑑x_1`$ (13)
$``$ $`{\displaystyle \frac{n}{(bϵ)}}{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle _{\mathrm{\Theta }_{𝒟,ϵ}^{(n,i)}}}\delta (bL(\theta )){\displaystyle \frac{d\underset{¯}{e}_1d\underset{¯}{e}_2\mathrm{}d\underset{¯}{e}_i\mathrm{}d\underset{¯}{e}_n}{e_1e_2\mathrm{}e_i\mathrm{}e_n}}𝑑x_1.`$
where a $`i`$ subscript indicates that element is left out of a product or sum. $`\mathrm{\Theta }_{𝒟,ϵ}^{(n,i)}`$ is the set of polygons with a long $`i`$th edge (that is, the set of chains in $`\stackrel{~}{\mathrm{\Theta }}_{𝒟,ϵ}^{(n,i)}`$ with $`x_{n+1}=x_1`$). We have carried out the integral $`d\underset{¯}{e}_i\delta ^{(2)}\left(\mathrm{\Sigma }_k\underset{¯}{e}_k\right)`$ and used the bound on $`e_i`$.
The third step is to eliminate an edge length parameter, using $`b=L(\stackrel{~}{\theta })`$, the length constraint. Let $`\varphi _i`$ denote the angle made by edge $`\underset{¯}{e}_i`$ to a fixed direction in the plane. In polar coordinates Equation (13) is
$$J_n\frac{n}{(bϵ)}\underset{i=1}{\overset{n}{}}_{\mathrm{\Theta }_{𝒟,ϵ}^{(n,i)}}\delta (bL(\theta ))𝑑e_1𝑑\varphi _1\mathrm{}𝑑e_i𝑑\varphi _i\mathrm{}𝑑e_n𝑑\varphi _n𝑑x_1.$$
(14)
For the polygon to close
$`e_i\mathrm{sin}(\varphi _i)`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{k=1}{ki}}{\overset{n}{}}}e_k\mathrm{sin}(\varphi _k),`$ (15)
$`e_i\mathrm{cos}(\varphi _i)`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{k=1}{ki}}{\overset{n}{}}}e_k\mathrm{cos}(\varphi _k),`$ (16)
and consequently
$$bL(\theta )=b\underset{\genfrac{}{}{0pt}{}{k=1}{ki}}{\overset{n}{}}e_ke_k\mathrm{cos}(\varphi _k\varphi _i).$$
Integrating $`de_j`$ for some $`j`$ may lead to an unbounded integrand. In order to control this, we partition $`\mathrm{\Theta }_{𝒟,ϵ}^{(n,i)}`$ on its angle variables. Let
$$\mathrm{\Theta }_{𝒟,ϵ}^{(n,i,j)}=\{\theta \mathrm{\Theta }_{𝒟,ϵ}^{(n,i)}|\frac{\pi }{2}<|\varphi _j\varphi _i|<\frac{3\pi }{2}\}$$
A polygon in $`\mathrm{\Theta }_{𝒟,ϵ}^{(n,i,j)}`$ has the property that the $`j`$th edge “turns back” from the direction of the long $`i`$th edge. There must be at least one such edge for the polygon to close. The sets $`\mathrm{\Theta }_{𝒟,ϵ}^{(n,i,j)}`$, $`j=1,2\mathrm{}n,ji`$ are not disjoint but their union covers $`\mathrm{\Theta }_{𝒟,ϵ}^{(n,i)}`$. From Equation (14)
$$J_n\frac{n}{(bϵ)}\underset{i=1}{\overset{n}{}}\underset{\genfrac{}{}{0pt}{}{j=1}{ji}}{\overset{n}{}}_{\mathrm{\Theta }_{𝒟,ϵ}^{(n,i,j)}}\delta (bL(\theta ))𝑑e_1𝑑\varphi _1\mathrm{}𝑑e_i𝑑\varphi _i\mathrm{}𝑑e_n𝑑\varphi _n𝑑x_1.$$
We may now apply the integral $`de_j`$ to the delta-function $`\delta (bL(\theta ))`$. We transform from $`e,\varphi `$ to $`e^{},\varphi ^{}`$ where $`\varphi _k^{}=\varphi _k`$ and $`e_k^{}=e_k`$ for $`1kn`$, $`kj`$, and $`\varphi _j^{}=\varphi _j`$ and
$$e_j^{}=e_je_j\mathrm{cos}(\varphi _j\varphi _i).$$
The Jacobian of the full transformation $`e,\varphi e^{},\varphi ^{}`$ is just
$`𝒥^1(e,\varphi e^{},\varphi ^{})`$ $`=`$ $`{\displaystyle \frac{e_j^{}}{e_j}}`$ (17)
$`=`$ $`1\mathrm{cos}(\varphi _j\varphi _i)e_j\mathrm{sin}(\varphi _j\varphi _i){\displaystyle \frac{\varphi _i}{e_j}}.`$
Repeated use of Equations (15) and (16) gives
$$\frac{\varphi _i}{e_j}=\frac{\mathrm{sin}(\varphi _j\varphi _i)}{e_i},$$
in Equation (17) and then using $`\pi /2<|\varphi _j\varphi _i|<3\pi /2`$, we have $`𝒥^1>1`$. The angle partition was needed to control this function. We can replace $`\delta (bL(\theta ))de_j`$ by one, and restrict the integration domain to polygons of length $`b`$, $`ie`$ set $`ϵ=0`$. We obtain the simplified bound
$$J_n\frac{n}{b}\underset{i=1}{\overset{n}{}}\underset{\genfrac{}{}{0pt}{}{j=1}{ji}}{\overset{n}{}}_{\mathrm{\Theta }_{𝒟,ϵ=0}^{(n,i,j)}}𝑑e_1𝑑\varphi _1\mathrm{}𝑑e_j𝑑\varphi _j\mathrm{}𝑑e_i𝑑\varphi _i\mathrm{}𝑑e_n𝑑\varphi _n𝑑x_1.$$
(18)
The last step is to bound the integral in Equation (18). Enlarge $`\mathrm{\Theta }_{𝒟,ϵ=0}^{(n,i,j)}`$ to allow each variable to range independently over its full domain, keeping only the bound on total edge length, $`L(\theta )=b`$, and requiring $`x_1`$ to remain in $`𝒟`$. This will include polygons with crossing edges and allow the polygon to overlap the border of $`𝒟`$. The integral $`dx_1`$ gives a factor $`A(𝒟)`$. Each angle variable ranges over $`0`$ to $`2\pi `$ contributing $`(2\pi )^{n1}`$. The edge integrals are over the $`(n2)`$-dimensional tetrahedron
$$e_1+e_2+\mathrm{}e_j+\mathrm{}e_i+\mathrm{}+e_nbb/n$$
of volume less than $`b^{n2}/(n2)!`$. Combining these factors with a factor of $`(n1)`$ from the sum over $`j`$, we obtain the bound on $`J_n`$ given in the Lemma. This is the end of the proof of the Lemma.
Equation (5) is obtained by evaluating the integral over $`b`$ in Equation (10) with the bound from Equation (12), and the Theorem follows directly from Equation (5).
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# 1 Introduction
## 1 Introduction
The issue of the cosmological constant has got renewed thrust from the observational evidence of an acceleration in the expansion of our Universe, recently reported by two different groups . There has been some controversy on the reliability of the results obtained from those observations and on its precise interpretation, by a number of different reasons . Anyway, there is presently reasonable consensus among the community of cosmologists that it certainly could happen that there is, in fact, an acceleration, and that it has the order of magnitude obtained in the above mentioned observations. In support of this consensus, the recently issued analysis of the data taken by the BOOMERANG and MAXIMA-1 balloons have been correspondigly crossed with those from the just mentioned observations, to conclude that the results of BOOMERANG and MAXIMA-1 can perfectly account for an accelerating universe and that, taking together both kinds of observations, one inferes that we most probably live in a flat universe. As a consequence, many theoretists have urged to try to explain this fact, and also to try to reproduce the precise value of the cosmological constant coming from these observations, in the available models .
Now, as crudely stated by Weinberg in a recent review , it is even more difficult to explain why the cosmological constant is so small but non-zero, than to build theoretical models where it exactly vanishes . Rigorous calculations performed in quantum field theory on the vacuum energy density, $`\rho _V`$, corresponding to quantum fluctuations of the fields we observe in nature, lead to values that are over 120 orders of magnitude in excess of the values allowed by observations of the space-time around us.
Rather than trying to understand the fine-tuned cancellation of such enormous values at this local level (a very difficult question that we are going to leave unanswered, and even unattended, here), in this paper we will elaborate on a quite simple and primitive idea (but, for the same reason, of far reaching, inescapable consequences), related with the global topology of the universe and in connection with the possibility that a very faint, massless scalar field pervading the universe could exist. Fields of this kind are ubiquitous in inflationary models, quintessence theories, and the like. In other words, we do not pretend here to solve the old problem of the cosmological constant, not even to contribute significantly to its understanding, but just to present an extraordinarily simple model which shows that the right order of magnitude of (some contributions to) $`\rho _V`$, in the precise range deduced from the astrophysical observations , e.g. $`\rho _V10^{10}`$ erg/cm<sup>3</sup>, are not difficult to obtain. To say it in different words, we only address here what has been termed by Weinberg the new cosmological constant problem.
In short, we shall assume the existence of a scalar field background extending throught the universe and shall calculate the contribution to the cosmological constant coming from the Casimir energy density corresponding to this field for some typical boundary conditions. The ultraviolet contributions will be safely set to zero by some mechanism of a fundamental theory. Another hypothesis will be the existence of both large and small dimensions (the total number of large spatial coordinates will be always three), some of which (from each class) may be compactified, so that the global topology of the universe will play an important role, too. There is by now a quite extense literature both in the subject of what is the global topology of spatial sections of the universe and also on the issue of the possible contribution of the Casimir effect as a source of some sort of cosmic energy, as in the case of the creation of a neutron star . There are arguments that favor different topologies, as a compact hiperbolic manifold for the spatial section, what would have clear observational consequences . Other interesting work along these lines was reported in and related ideas have been discussed very recently in . However, our paper differs from all those in several respects. To begin, the emphasis is put now in obtaining the right order of magnitude for the effect, e.g., one that matches the recent observational results. At the present stage, in view of the observational precission, it has no sense to consider the whole amount of possibilities concerning the nature of the field, the different models for the topology of the universe, and the different boundary conditions possible.
At this level, from our previous experience in these calculations and from the many tables (see, e.g., where precise values of the Casimir effect corresponding to a number of different configurations have been reported), we realize that the range of orders of magnitude of the vacuum energy density for the most common possibilities is not so widespread, and may only differ by at most a couple of digits. This will allow us, both for the sake of simplicity and universality, to deal with a most simple situation, which is the one corresponding to a scalar field with periodic boundary conditions. Actually, as explained in in detail, all other cases for parallel plates, with any of the usual boundary conditions, can be reduced to this one, from a mathematical viewpoint.
## 2 Two basic space-time models
Let us thus consider a universe with a space-time of one of the following types: $`𝐑^{𝐝+\mathrm{𝟏}}\times 𝐓^p\times 𝐓^q`$, $`𝐑^{𝐝+\mathrm{𝟏}}\times 𝐓^p\times 𝐒^q,\mathrm{}`$, which are actually plausible models for the space-time topology. A (nowadays) free scalar field pervading the universe will satisfy
$`(\mathrm{}+M^2)\varphi =0,`$ (1)
restricted by the appropriate boundary conditions (e.g., periodic, in the first case considered). Here, $`d0`$ stands for a possible number of non-compactified dimensions.
Recall now that the physical contribution to the vacuum or zero-point energy $`<0|H|0>`$ (where $`H`$ is the Hamiltonian corresponding to our massive scalar field and $`|0>`$ the vacuum state) is obtained on subtracting to these expression —with the vacuum corresponding to our compactified spatial section with the assumed boundary conditions— the vacuum energy corresponding to the same situation with the only change that the compactification is absent (in practice this is done by conveniently sending the compactification radii to infinity). As well known, both of these vacuum energies are in fact infinite, but it is its difference
$`E_C=<0|H|0>|_R<0|H|0>|_R\mathrm{}`$ (2)
(where $`R`$ stands here for a typical compactification length) that makes physical sense, giving rise to the finite value of the Casimir energy $`E_C`$, which will depend on $`R`$ (after a well defined regularization/renormalization procedure is carried out). In fact we will discuss the Casimir (or vacuum) energy density, $`\rho _C=E_C/V`$, which can account for either a finite or an infinite volume of the spatial section of the universe (from now on we shall assume that all diagonalizations already correspond to energy densities, and the volume factors will be replaced at the end). In terms of the spectrum $`\{\lambda _n\}`$ of $`H`$:
$`<0|H|0>={\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}\lambda _n,`$ (3)
where the sum over $`n`$ is a sum over the whole spectrum, which involves, in general, several continuum and several discrete indices. The last appear tipically when compactifying the space coordinates (much in the same way as time compactification gives rise to finite-temperature field theory), as in the cases we are going to consider. Thus, the cases treated will involve integration over $`d`$ continuous dimensions and multiple summations over $`p+q`$ indices (for a pedagogical description of this procedure, see ).
To be precise, the physical vacuum energy density corresponding to our case, where the contribution of a scalar field, $`\varphi `$ in a (partly) compactified spatial section of the universe is considered, will be denoted by $`\rho _\varphi `$ (note that this is just the contribution to $`\rho _V`$ coming from this field, there might be other, in general). It is given by
$`\rho _\varphi ={\displaystyle \frac{1}{2}}{\displaystyle \underset{\text{k}}{}}{\displaystyle \frac{1}{\mu }}\left(k^2+M^2\right)^{1/2},`$ (4)
where the sum $`_\text{k}`$ is a generalized one (as explained above) and $`\mu `$ is the usual mass-dimensional parameter to render the eigenvalues adimensional (we take $`\mathrm{}=c=1`$ and shall insert the dimensionfull units only at the end of the calculation). The mass $`M`$ of the field will be here considered to be arbitrarily small and will be kept different from zero, for the moment, for computational reasons —as well as for physical ones, since a very tiny mass for the field can never be excluded. Some comments about the choice of our model are in order. The first seems obvious: the coupling of the scalar field to gravity should be considered. This has been done in all detail in, e.g., (see also the references therein). The conclusion is that taking it into account does not change the results to be obtained here. Of course, the renormalization of the model is rendered much more involved, and one must enter a discussion on the orders of magnitude of the different contributions, which yields, in the end, an ordinary perturbative expansion, the coupling constant being finally re-absorbed into the mass of the scalar field. In conclusion, we would not gain anything new by taking into account the coupling of the scalar field to gravity. Owing, essentially, to the smallness of the resulting mass for the scalar field, one can prove that, quantitatively, the difference in the final result is at most of a few percent.
Another important consideration is the fact that our model is stationary, while the universe is expanding. Again, careful calculations show that this effect can actually be dismissed at the level of our order of magnitude calculation, since its value cannot surpass the one that we will get (as is seen from the present value of the expansion rate $`\mathrm{\Delta }R/R10^{10}`$ per year or from direct consideration of the Hubble coefficient). As before, for the sake of simplicity, and in order to focus just on the essential issues of our argument, we will perform a (momentaneously) static calculation. As a consequence, the value of the Casimir energy density, and of the cosmological constant, to be obtained will correspond to the present epoch, and are bound to change with time.
The last comment at this point would be that (as shown by the many references mentioned above), the idea presented here is not entirely new. However, the simplicity and the generality of its implementation below are indeed brand new. The issue at work here is absolutely independent of any specific model, the only assumptions having been clearly specified before (e.g., existence of a very light scalar field and of some reasonably compactified scales, see later). Secondly, it will turn out, in the end, that the only ‘free parameter’ to play with (the number of compactified dimensions) will actually not be that ‘free’ but, on the contray, very much constrained to have an admissible value. This will become clear after the calculations below. Thirdly, although the calculation may seem easy to do, in fact it is not so. Recently derived reflection identities will allow us to to perform it analitically, for the first time.
## 3 The vacuum energy density and its regularization
To exhibit explicitly a couple of the wide family of cases considered, let us write down in detail the formulas corresponding to the two first topologies, as described above. For a ($`p,q`$)-toroidal universe, with $`p`$ the number of ‘large’ and $`q`$ of ‘small’ dimensions:
$`\rho _\varphi ={\displaystyle \frac{\pi ^{d/2}}{2^d\mathrm{\Gamma }(d/2)_{j=1}^pa_j_{h=1}^qb_h}}{\displaystyle _0^{\mathrm{}}}𝑑kk^{d1}{\displaystyle \underset{\text{n}_p=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\text{m}_q=\mathrm{}}{\overset{\mathrm{}}{}}}\left[{\displaystyle \underset{j=1}{\overset{p}{}}}\left({\displaystyle \frac{2\pi n_j}{a_j}}\right)^2+{\displaystyle \underset{h=1}{\overset{q}{}}}\left({\displaystyle \frac{2\pi m_h}{b_h}}\right)^2+M^2\right]^{1/2}`$ (5)
$`{\displaystyle \frac{1}{a^pb^q}}{\displaystyle \underset{\text{n}_p,\text{m}_q=\mathrm{}}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{1}{a^2}}{\displaystyle \underset{j=1}{\overset{p}{}}}n_j^2+{\displaystyle \frac{1}{b^2}}{\displaystyle \underset{h=1}{\overset{q}{}}}m_h^2+M^2\right)^{(d+1)/2+1},`$ (6)
where the last formula corresponds to the case when all large (resp. all small) compactification scales are the same. In this last expression the squared mass of the field should be divided by $`4\pi ^2\mu ^2`$, but we have renamed it again $`M^2`$ to simplify the ensuing formulas (as $`M`$ is going to be very small, we need not keep track of this change). We also will not take care for the moment of the mass-dim factor $`\mu `$ in other places $``$as is usually done$``$ since formulas would get unnecesarily complicated and there is no problem in recovering it at the end of the calculation. For a ($`p`$-toroidal, $`q`$-spherical)-universe, the expression turns out to be
$`\rho _\varphi `$ $`=`$ $`{\displaystyle \frac{\pi ^{d/2}}{2^d\mathrm{\Gamma }(d/2)_{j=1}^pa_jb^q}}{\displaystyle _0^{\mathrm{}}}𝑑kk^{d1}{\displaystyle \underset{\text{n}_p=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}P_{q1}(l)\left[{\displaystyle \underset{j=1}{\overset{p}{}}}\left({\displaystyle \frac{2\pi n_j}{a_j}}\right)^2+{\displaystyle \frac{Q_2(l)}{b^2}}+M^2\right]^{1/2}`$ (7)
$``$ $`{\displaystyle \frac{1}{a^pb^q}}{\displaystyle \underset{\text{n}_p=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}P_{q1}(l)\left({\displaystyle \frac{4\pi ^2}{a^2}}{\displaystyle \underset{j=1}{\overset{p}{}}}n_j^2+{\displaystyle \frac{l(l+q)}{b^2}}+M^2\right)^{(d+1)/2+1},`$ (8)
where $`P_{q1}(l)`$ is a polynomial in $`l`$ of degree $`q1`$, and where the second formula corresponds to the similar situation as the second one before. On dealing with our observable universe, in all these expression we assume that $`d=3p`$, the number of non-compactified, ‘large’ spatial dimensions (thus, no $`d`$ dependence will remain).
As is clear, all these expressions for $`\rho _\varphi `$ need to be regularized. We will use zeta function regularization, taking advantage of the very powerful equalities that have been derived recently , which reduce the enormous burden of such computations to the easy application of some formulas. For the sake of completeness, let us very briefly summarize how this works . We deal here only with the case when the spectrum of the Hamiltonian operator is known explicitly. Going back to the most general expressions of the Casimir energy corresponding to this case, namely Eqs. (LABEL:c1)-(4), we replace the exponents in them with a complex variable, $`s`$, thus obtaining the zeta function associated with the operator as:
$`\zeta (s)={\displaystyle \frac{1}{2}}{\displaystyle \underset{\text{k}}{}}\left({\displaystyle \frac{k^2+M^2}{\mu ^2}}\right)^{s/2}.`$ (9)
The next step is to perform the analytic continuation of the zeta function from a domain of the complex $`s`$-plane with Re $`s`$ big enough (where it is perfectly defined by this sum) to the point $`s=1`$, to obtain:
$`\rho _\varphi =\zeta (1).`$ (10)
The effectiveness of this method has been sufficiently described before (see, e.g., ). As we know from precise Casimir calculations in those references, no further subtraction or renormalization is needed in the cases here considered, in order to obtain the physical value for the vacuum energy density (there is actually a subtraction at infinity taken into account, as carefully described above, but it is of null value, and no renormalization, not even a finite one, very common to other frameworks, applies here).
Using the recent formulas that generalize the well-known Chowla-Selberg expression to the situations considered above, Eqs. (5) and (7) —namely, multidimensional, massive cases— we can provide arbitrarily accurate results for different values of the compactification radii. However, as argued above we can only aim here at matching the order of magnitude of the Casimir value and, thus, we shall just deal with the most simple case of Eq. (6) (or (5), which yield the same orders of magnitude as the rest). Also in accordance with this observation, we notice that among the models here considered and which lead to the values that will be obtained below, there are in particular the very important typical cases of isotropic universes with the spherical topology. As all our discussion here is in terms of orders of magnitude and not of precise values with small errors, all these cases are included on equal footing. But, on the other hand, it has no sense to present a lengthy calculation dealing in detail with all the possible spatial geometries. Anyhow, all these calculations are very similar to the one to be carried out here, as has been described in detail elsewhere .
For the analytic continuation of the zeta function corresponding to (5), we obtain :
$`\zeta (s)`$ $`=`$ $`{\displaystyle \frac{2\pi ^{s/2+1}}{a^{p(s+1)/2}b^{q(s1)/2}\mathrm{\Gamma }(s/2)}}{\displaystyle \underset{\text{m}_q=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{h=0}{\overset{p}{}}}({}_{h}{}^{p})2^h{\displaystyle \underset{\text{n}_h=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{_{j=1}^hn_j^2}{_{k=1}^qm_k^2+M^2}}\right)^{(s1)/4}`$ (11)
$`\times K_{(s1)/2}\left[{\displaystyle \frac{2\pi a}{b}}\sqrt{{\displaystyle \underset{j=1}{\overset{h}{}}}n_j^2\left({\displaystyle \underset{k=1}{\overset{q}{}}}m_k^2+M^2\right)}\right],`$
where $`K_\nu (z)`$ is the modified Bessel function of the second kind. Having performed already the analytic continuation, this expression is ready for the substitution $`s=1`$, and yields
$`\rho _\varphi ={\displaystyle \frac{1}{a^pb^{q+1}}}{\displaystyle \underset{h=0}{\overset{p}{}}}({}_{h}{}^{p})2^h{\displaystyle \underset{\text{n}_h=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\text{m}_q=\mathrm{}}{\overset{\mathrm{}}{}}}\sqrt{{\displaystyle \frac{_{k=1}^qm_k^2+M^2}{_{j=1}^hn_j^2}}}K_1\left[{\displaystyle \frac{2\pi a}{b}}\sqrt{{\displaystyle \underset{j=1}{\overset{h}{}}}n_j^2\left({\displaystyle \underset{k=1}{\overset{q}{}}}m_k^2+M^2\right)}\right].`$ (12)
Now, from the behaviour of the function $`K_\nu (z)`$ for small values of its argument,
$`K_\nu (z){\displaystyle \frac{1}{2}}\mathrm{\Gamma }(\nu )(z/2)^\nu ,z0,`$ (13)
we obtain, in the case when $`M`$ is very small,
$`\rho _\varphi `$ $`=`$ $`{\displaystyle \frac{1}{a^pb^{q+1}}}\{MK_1\left({\displaystyle \frac{2\pi a}{b}}M\right)+{\displaystyle \underset{h=0}{\overset{p}{}}}({}_{h}{}^{p})2^h{\displaystyle \underset{\text{n}_h=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{M}{\sqrt{_{j=1}^hn_j^2}}}K_1\left({\displaystyle \frac{2\pi a}{b}}M\sqrt{{\displaystyle \underset{j=1}{\overset{h}{}}}n_j^2}\right)`$ (14)
$`+𝒪\left[q\sqrt{1+M^2}K_1\left({\displaystyle \frac{2\pi a}{b}}\sqrt{1+M^2}\right)\right]\}.`$
At this stage, the only presence of the mass-dim parameter $`\mu `$ is as $`M/\mu `$ everywhere. This does not conceptually affect the small-$`M`$ limit, $`M/\mu <<b/a`$. Using (13) and inserting now in the expression the $`\mathrm{}`$ and $`c`$ factors, we finally get
$`\rho _\varphi ={\displaystyle \frac{\mathrm{}c}{2\pi a^{p+1}b^q}}[1+{\displaystyle \underset{h=0}{\overset{p}{}}}({}_{h}{}^{p})2^h\alpha ]+𝒪\left[qK_1\left({\displaystyle \frac{2\pi a}{b}}\right)\right],`$ (15)
where $`\alpha `$ is some finite constant, computable and under control, which is obtained as an explicit geometrical sum in the limit $`M0`$. It is remarkable that we here obtain such a well defined limit, independent of $`M^2`$, provided that $`M^2`$ is small enough. In other words, a physically very nice situation turns out to correspond, precisely, to the mathematically rigorous case. This is moreover, let me repeat, the kind of expression that one gets not just for the model considered, but for many general cases, corresponding to different fields, topologies, and boundary conditions —aside from the sign in front of the formula, that may change with the number of compactified dimensions and the nature of the boundary conditions (in particular, for Dirichlet boundary conditions one obtains a value in the same order of magnitude but of opposite sign).
## 4 Numerical results
For the most common variants, the constant $`\alpha `$ in (15) has been calculated to be of order $`10^2`$, and the whole factor, in brackets, of the first term in (15) has a value of order $`10^7`$. This shows the value of a precise calculation, as the one undertaken here, together with the fact that just a naive consideration of the dependences of $`\rho _\varphi `$ on the powers of the compactification radii, $`a`$ and $`b`$, is not enough in order to obtain the correct result. Notice, moreover, the non-trivial change in the power dependencies from going from Eq. (14) to Eq. (15).
For the compactification radii at small scales, $`b`$, we shall simply take the magnitude of the Planck length, $`bl_{P(lanck)}`$, while the typical value for the large scales, $`a`$, will be taken to be the present size of the observable universe, $`aR_U`$. With this choice, the order of the quocient $`a/b`$ in the argument of $`K_1`$ is as big as: $`a/b10^{60}`$. Thus, we see immediately that, in fact, the final expression for the vacuum energy density is completely independent of the mass $`M`$ of the field, provided this is very small (eventually zero). In fact, since the last term in Eq. (15) is exponentially vanishing, for large arguments of the Bessel function $`K_1`$, this contribution is zero, for all practical purposes, what is already a very nice result. Taken in ordinary units (and after tracing back all the transformations suffered by the mass term $`M`$) the actual bound on the mass of the scalar field is $`M1.2\times 10^{32}`$ eV, that is, physically zero, since it is lower by several orders of magnitude than any bound comming from the more usual SUSY theories $``$where in fact scalar fields with low masses of the order of that of the lightest neutrino do show up , which may have observable implications.
By replacing all these values in Eq. (15), we obtain the results listed in Table 1, for the orders of magnitude of the vacuum energy density corresponding to a sample of different numbers of compactified (large and small) dimensions and for different values of the small compactification length in terms of the Planck length. Notice again that the total number of large space dimensions is three, as corresponds to our observable universe. As we see from the table, good coincidence with the observational value for the cosmological constant is obtained for the contribution of a massless scalar field, $`\rho _\varphi `$, for $`p`$ large compactified dimensions and $`q=p+1`$ small compactified dimensions, $`p=0,\mathrm{},3`$, and this for values of the small compactification length, $`b`$, of the order of 100 to 1000 times the Planck length $`l_P`$ (what is actually a very reasonable conclusion, according also to other approaches). To be noticed is the fact that full agreement is obtained only for cases where there is exactly one small compactified dimension in excess of the number of large compactified dimensions. We must point out that the $`p`$ large and $`q`$ small dimensions are not all that are supposed to exist (in that case $`p`$ should be at least, and at most, 3 and the other cases would lack any physical meaning). In fact, as we have pointed out before, $`p`$ and $`q`$ refer to the compactified dimensions only, but there may be other, non-compactifed dimensions (exactly $`3p`$ in the case of the ‘large’ ones), what translates into a slight modification of the formulas above, but does not change the order of magnitude of the final numbers obtained, assuming the most common boundary conditions for the non-compactified dimensions (see e.g. for an explanation of this technical point). In particular, the cases of pure spherical compactification and of mixed toroidal (for small magnitudes) and spherical (for big ones) compactification can be treated in this way and yield results in the same order of magnitude range. Both these cases correspond to (observational) isotropic spatial geometries. Also to be remarked again is the non-triviality of these calculations, when carried out exactly, as done here, to the last expression, what is apparent from the use of the generalized Chowla-Selberg formula. Simple power counting is absolutely unable to provide the correct order of magnitude of the results.
## 5 Conclusions
Dimensionally speaking, within the global approach adopted in the present paper everything is dictated, in the end, by the two basic lengths in the problem, which are its Planck value and the radius of the observable Universe. Just by playing with these numbers in the context of our (very precise) calculation of the Casimir effect, we have shown that the observed value of $`\rho _V`$ may be remarkably well fitted, under general hypothesis, for the most common models of the space-time topology. Notice also that the most precise fits with the observational value of the cosmological constant are obtained for $`b`$ between $`b=100l_P`$ and $`b=1000l_P`$, with (1,2) and (2,3) compactified dimensions, respectively. The fact that the value obtained for the cosmological constant is so sensitive to the input may be viewed as a drawback but also, on the contrary, as a very positive feature of our model. For one, the table has a sharp discriminating power. In other words, there is in fact no tuning of a ‘free parameter’ in our model and the number of large compactified dimensions could have been fixed beforehand, to respect what we know already of our observable universe.
Also, it proves that the observational value is not easy at all to obtain. The table itself proves that there is only very little chance of getting the right figure (a truly narrow window, since very easily we are off by several orders of magnitude). In fact, if we trust this value with the statistics at hand, we can undoubtedly claim $``$through use of our model$``$ that the ones so clearly picked up by Table 1 are the only two possible configurations of our observable universe (together with a couple more coming from corresponding spherical compactifications). And all them correspond to a marginally closed universe, in full agreement too with other completely independent analysis of the observational data .
Many questions may be posed to the simple models presented here, as concerning the dynamics of the scalar field, its couplings with gravity and other fields, a possible non-symmetrical behaviour with respect to the large and small dimensions, or the relevance of vacuum polarization (see , concerning this last point). Above we have already argued that they can be proven to have little influence on the final numerical result (cf., in particular, the mass obtained for the scalar field in Ref. , extremely close to our own result, and the corresponding discussion there). From the very existence and specific properties of the cosmic microwave radiation (CMB) $``$which mimics somehow the situation described (the ‘mass’ corresponding to the CMB is also in the sub-lightest-neutrino range)$``$ we are led to the conclusion that such a field could be actually present, unnoticed, in our observable universe. In fact, the existence of scalar fields of very low masses is also demanded by other frameworks, as SUSY models, where the scaling behaviour of the cosmological constant has been considered .
Let us finally recall that the Casimir effect is an ubiquitous phenomena. Its contribution may be small (as it seems to be the case, yet controverted, to sonoluminiscence), of some 10-30$`\%`$ (that is, of the right order of magnitude, as in wetting phenomena involving He in condensed matter physics), or even the whole thing (as in recent, dedicated experimental confirmations of the effect). Here we have seen that it provides a contribution of the right order of magnitude, corresponding to our present epoch in the evolution of the universe. The implication that this calculation bears for the early universe and inflation is not clear from the final result, since it should be adapted to the situation and boundary conditions corresponding to those primeval epochs, what cannot be done straightforwardly. Work along this line is in progress.
Acknowledgments
I am grateful to Robert Kirshner, Tom Mongan, Varun Sahni and Joan Solà for important comments. Thanks are also given to the referee for interesting suggestions that have led to an improvement of the paper. This investigation has been supported by DGICYT (Spain), project PB96-0925 and by CIRIT (Generalitat de Catalunya), grants 1997SGR-00147 and 1999SGR-00257.
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# Supergravity Brane Cosmologies
## I Introduction
There has been considerable interest recently in the possibility that our observable universe may be viewed as a $`p`$–brane embedded in a higher–dimensional spacetime. In this picture, the gauge interactions are confined to the brane, but gravity may propagate in the bulk. This change in viewpoint has been partially motivated by advances in our understanding of non–perturbative string theory. For example, the strongly–coupled, field theoretic limit of the $`\mathrm{E}_8\times \mathrm{E}_8`$ heterotic string has been described by Hořava and Witten as $`D=11`$ supergravity on an orbifold $`\mathrm{S}^1/\mathrm{Z}_2`$, where the two sets of $`\mathrm{E}_8`$ gauge fields are confined to the orbifold fixed planes . Compactification of this theory on a Calabi–Yau three–fold results in an effective five–dimensional theory containing a superpotential that supports a pair of parallel 3–branes (domain walls) .
A related five–dimensional model with an extra orbifold dimension was recently proposed by Randall and Sundrum. In this model, there are two 3–branes with equal and opposite tensions at the orbifold fixed points and our universe is identified as the positive tension brane. The existence of a negative cosmological constant in the five–dimensional bulk results in a curved background. This supports a bound state of the higher–dimensional graviton that is localized to the 3–brane and, consequently, the size of the extra dimension can be arbitrarily large.
This picture differs significantly from traditional Kaluza–Klein compactification, where the higher–dimensional universe is represented as a product space and the four–dimensional Planck mass is determined by the volume of the extra dimensions. In the braneworld scenario, the geometry is non–factorizable because the brane tension induces a ‘warp factor’ in the metric.
A crucial question that must be addressed is whether the braneworld scenario is consistent with our understanding of early universe cosmology. A central paradigm of the early universe is that of cosmological inflation, where the universe undergoes an epoch of accelerated expansion. It is therefore important to develop brane cosmologies where inflation may proceed.
The cosmological implications of the Randall–Sundrum model have been considered by a number of authors and other examples of curved branes were recently presented in Ref. . In view of the above developments, we show in this paper that a wide class of higher–dimensional supergravity theories admit solutions that may be interpreted as (non–supersymmetric) brane cosmologies, where the dilaton field varies non–trivially over the world–volume. The effect of this variation is formally equivalent, after appropriate field redefinitions, to the introduction of a massless, minimally coupled scalar field on the brane. Hence, these solutions are relevant to inflationary models based on string theory such as the pre–big bang scenario , since, in this model, the accelerated expansion is driven by the kinetic energy of the dilaton field.
The paper is organised as follows. The class of brane cosmologies is derived in Section II. In Secion III, we discuss some of the models from an eleven–dimensional perspective. Brane cosmologies in the Hořava–Witten heterotic theory are found in Section IV and we conclude with a discussion on dilaton–driven inflation in Section V.
## II Brane Cosmology in Supergravity Theories
### A Ricci–Flat Branes
We consider the class of $`D`$–dimensional effective actions, where the graviton, $`g_{MN}`$, is coupled to a dilaton field, $`\mathrm{\Phi }`$, and the $`q`$–form field strength, $`F_{[q]}=dA_{[q1]}`$, of an antisymmetric gauge field, $`A_{[q1]}`$:
$$S=d^Dx\sqrt{|g|}\left[R\frac{1}{2}\left(\mathrm{\Phi }\right)^2\frac{1}{2q!}e^{\alpha \mathrm{\Phi }}F_{[q]}^2\right],$$
(2.1)
where $`R`$ is the Ricci curvature scalar of the spacetime, $`g\mathrm{det}g_{AB}`$ and the coupling parameter, $`\alpha `$, is a constant. Given appropriate conditions on the form fields, action (2.1) represents a consistent truncation of the bosonic sectors of the $`D=10`$, $`N=2`$ supergravity theories<sup>*</sup><sup>*</sup>*Throughout this paper, the Chern–Simons terms that also arise in the effective actions are trivial for the solutions we consider and we do not present them here.. The toroidal compactification of the type II theories to $`D<10`$ also results in an action of the form given in Eq. (2.1) . The effective action for M–theory is given by Eq. (2.1), where $`D=11`$, $`q=0`$ and $`\mathrm{\Phi }=0`$ .
The field equations derived by varying the action (2.1) areIn this paper, upper case, Latin indices take values in the range $`A=(0,1,\mathrm{},D1)`$, lower case, Greek indices vary from $`\mu =(0,1,\mathrm{},d1)`$ and lower case, Latin indices from $`a=(d,\mathrm{},D1)`$. The dimensionality of the world–volume of the brane is denoted $`d_W`$, the dimensionality of the transverse space is $`d_T`$ and the spacetime metric has signature $`(,+,\mathrm{},+)`$.
$`R_{AB}={\displaystyle \frac{1}{2}}_A\mathrm{\Phi }_B\mathrm{\Phi }+{\displaystyle \frac{1}{2q!}}e^{\alpha \mathrm{\Phi }}(qF_{AC_2\mathrm{}C_q}F_{B}^{}{}_{}{}^{C_2\mathrm{}C_q}`$ (2.2)
$`\left({\displaystyle \frac{q1}{D2}}\right)F_{[q]}^2g_{AB})`$ (2.3)
$`^2\mathrm{\Phi }={\displaystyle \frac{\alpha }{2q!}}e^{\alpha \mathrm{\Phi }}F_{[q]}^2`$ (2.4)
$`_A\left(e^{\alpha \mathrm{\Phi }}F^{AB_2\mathrm{}B_q}\right)=0.`$ (2.5)
When $`q=0`$, the field strength may be interpreted as a cosmological constant. For $`q=1`$, it represents the gradient of a massless axion field.
In $`D`$ dimensions, a solitonic $`(Dq2)`$–brane is supported by the ‘magnetic’ charge of a $`q`$-form field strength. Thus, M–theory contains a 5–brane due to the four–form field strength (M5–brane). Moreover, both ten–dimensional type II theories admit a 5–brane supported by the Neveu–Schwarz/Neveu–Schwarz (NS–NS) three–form field strength (NS5–brane) . These theories also admit branes supported by Ramond-Ramond (RR) fields (D$`p`$–branes). (For a review, see, e.g., Ref. ). The RR sector of the type IIB theory contains a one–form, a three–form and a (self–dual) five–form that result in a D7–, D5– and D3–brane, respectively. The massless type IIA theory, on the other hand, admits a D6– and D4–brane, whereas the massive type IIA theory due to Romans also admits a D8–brane supported by a 0–form. In general, the coupling of the RR $`q`$—forms to the dilaton in the type II theories is given by $`\alpha =(5q)/2`$.
Recently, Brecher and Perry showed that Eqs. (2.2)–(2.5) admit solitonic D$`p`$– and M–brane solutions with a Ricci–flat world–volume, $`f_{\mu \nu }=f_{\mu \nu }(x^\rho )`$:
$`ds_D^2=H^mf_{\mu \nu }dx^\mu dx^\nu +H^n\delta _{ij}dy^idy^j`$ (2.6)
$`e^\mathrm{\Phi }=H^{\alpha /2}`$ (2.7)
$`F_{a_1\mathrm{}a_q}=\lambda ϵ_{a_1\mathrm{}a_qb}{\displaystyle \frac{y^b}{r^{q+1}}},`$ (2.8)
where
$$m=\frac{q1}{D2},n=\frac{Dq1}{D2}$$
(2.9)
and $`d_T=q+1`$. The coordinates on the space transverse to the brane are $`\{y^i\}`$ and $`H=H(r)`$ is an harmonic function on this space:
$$\delta ^{ij}_i_jH=0,$$
(2.10)
where $`r`$ represents the radial coordinate. The components of the alternating tensor in Eq. (2.8) are $`\pm 1`$ and $`\lambda `$ is a constant. The solutions preserve some fraction of the supersymmetry if the world–volume admits parallel spinors . For the M5–brane, $`m=1/3`$, $`n=2/3`$ and $`\alpha =0`$.
### B Curved Branes
In this paper we allow the transverse space to depend directly on the world–volume coordinates of the brane by introducing a scalar function $`B=B(x^\mu )`$:
$$ds_D^2=H^mf_{\mu \nu }dx^\mu dx^\nu +e^{2B}H^n\delta _{ij}dy^idy^j,$$
(2.11)
where $`(m,n)`$ and $`H`$ are defined in Eqs. (2.9) and (2.10), respectively. In the standard Kaluza–Klein picture, the degree of freedom, $`B`$, would represent the ‘breathing mode’ of the internal dimensions and would play the role of a modulus field in the lower–dimensional theory.
The components of the Ricci tensor for the metric (2.11) are given by
$`R_{\mu \nu }=\overline{R}_{\mu \nu }d_T\left(\overline{}_{\mu \nu }B+\overline{}_\mu B\overline{}_\nu B\right)`$ (2.12)
$`+e^{2B}f_{\mu \nu }\widehat{Q}`$ (2.13)
$`R_{\mu b}={\displaystyle \frac{m}{2}}(D2)\overline{}_\mu B{\displaystyle \frac{_bH}{H}}`$ (2.14)
$`R_{ab}=\widehat{R}_{ab}e^{2B}H^{nm}\delta _{ab}\left(\overline{}^2B+d_T\left(\overline{}B\right)^2\right),`$ (2.15)
where an overbar identifies terms that are calculated with the world–volume metric, $`f_{\mu \nu }`$. The quantity, $`\widehat{Q}`$, represents a sum of terms depending on the function $`H`$ and its first and second derivatives. This sum is identical to the one that is obtained in the Ricci–flat limit, where $`B=0`$. Likewise, $`\widehat{R}_{ab}`$ represents the transverse components of the Ricci tensor calculated from the metric (2.6).
We proceed to search for solutions to the field equations (2.2)–(2.5) for the ansatz (2.11). We assume that the dilaton field has a separable form such that $`\mathrm{\Phi }(x,y)=\mathrm{\Phi }_1(x)+\mathrm{\Phi }_2(y)`$, where the transverse–dependent part is given by the right–hand side of Eq. (2.7). Moreover, we assume that the field strength satisfies Eq. (2.8). The introduction of a modulus field, $`B`$, leads to a non–trivial, off–diagonal component of the $`D`$–dimensional Ricci tensor. Nevertheless, the $`(\mu b)`$–component of the Einstein equations (2.2) can be directly integrated to yield the constraint
$$\alpha \mathrm{\Phi }_1=2(q1)B$$
(2.16)
relating the dilaton and modulus fields.
The question that now arises is whether Eq. (2.16) is compatible with the remaining field equations. We deduce by direct substitution that Eq. (2.5) is solved by Eqs. (2.7), (2.8) and (2.11). Furthermore, by imposing the constraint
$$\overline{}^2B+d_T\left(\overline{}B\right)^2=0$$
(2.17)
on the modulus field, we find that the $`(ab)`$–component of the Einstein equations (2.2) and the dilaton field equation (2.4) are also solved by Eqs. (2.7), (2.8) and (2.11). Finally, the $`(\mu \nu )`$–component of Eq. (2.2) is solved by the same conditions if the Ricci tensor of the world–volume metric satisfies
$$\overline{R}_{\mu \nu }=d_T\overline{}_{\mu \nu }B+\left(d_T+\frac{2(q1)^2}{\alpha ^2}\right)\overline{}_\mu B\overline{}_\nu B.$$
(2.18)
Eqs. (2.17) and (2.18) may be expressed in a more familiar form by performing the conformal transformation
$$\stackrel{~}{f}_{\mu \nu }=\mathrm{\Omega }^2f_{\mu \nu },\mathrm{\Omega }^2e^{2d_TB/(d_W2)}$$
(2.19)
on the world–volume metric and rescaling the modulus field, $`BQ^1\chi `$, where
$$Q\sqrt{2}\left(q+1+\frac{2(q1)^2}{\alpha ^2}+\frac{(q+1)^2}{d_W2}\right)^{1/2}.$$
(2.20)
This implies that
$`\stackrel{~}{R}_{\mu \nu }={\displaystyle \frac{1}{2}}\stackrel{~}{}_\mu \chi \stackrel{~}{}_\nu \chi `$ (2.21)
$`\stackrel{~}{}^2\chi =0`$ (2.22)
and Eqs. (2.21) and (2.22) represent the $`d_W`$–dimensional field equations for a massless scalar field minimally coupled to Einstein gravity.
Thus, modulo a solution to Eqs. (2.21) and (2.22), we have found a class of solutions to the supergravity field equations (2.2)–(2.5) that reduce to the Ricci–flat branes (2.6)–(2.8) in the limit where the dilaton field is constant on the world–volume. Since the dependence of these solutions on the transverse coordinates is identical to that of the Ricci–flat limit, they may be interpreted as brane cosmologies, where the curvature of the brane is induced by the variation of the dilaton field over the world–volume.
In particular, the D$`p`$–brane cosmologies have a metric given in the Einstein frame by
$$ds_{\mathrm{D}p}^2=H^{(1q)/8}f_{\mu \nu }dx^\mu dx^\nu +H^{(9q)/8}e^{ϵ\mathrm{\Phi }_1}\delta _{ij}dy^idy^j,$$
(2.23)
where $`ϵ(5q)/[2(q1)]`$, $`\mathrm{\Phi }_1=2B/ϵ`$, $`e^{\mathrm{\Phi }_2}=H^{(q5)/4}`$ and $`\{f_{\mu \nu },B\}`$ solve Eqs. (2.17) and (2.18). The corresponding metric in the string frame, $`g_{AB}^{(s)}=\mathrm{\Theta }^2g_{AB}`$, where $`\mathrm{\Theta }^2e^{\mathrm{\Phi }/2}`$, is given by
$`ds_s^2=H^{1/2}e^{\mathrm{\Phi }_1/2}f_{\mu \nu }dx^\mu dx^\nu `$ (2.24)
$`+H^{1/2}e^{2\mathrm{\Phi }_1/(q1)}{\displaystyle \underset{j=1}{\overset{d_T}{}}}dy_j^2.`$ (2.25)
In effect, solutions of this type exist because Eqs. (2.2)–(2.5) can each be separated into a sector that depends only on the world–volume coordinates and a sector that depends only on the transverse coordinates. If the separation constants are then set to zero, the latter sector reduces to the field equations that arise in the Ricci–flat limit.
The dilaton and modulus fields must vary in direct proportion to one other and the constant of proportionality depends on the degree of the form field and its coupling to the dilaton. However, it is independent of the dimensionality of spacetime. There are two cases where the world–volume must remain Ricci–flat, however, at least within the context of the assumptions made above. The dilaton field must depend only on the radial coordinate of the transverse space in the case of a $`(D3)`$–brane $`(q=1)`$ or when it is not directly coupled to the form field $`(\alpha =0)`$.
In the following Section, we consider some of the above brane cosmologies from an eleven–dimensional, M–theoretic perspective.
## III Eleven–Dimensional Interpretations
### A D8–Brane Cosmology
An important brane that has received considerable attention is the D8–brane of Romans’ massive IIA theory . This domain wall is supported by a $`0`$–form coupled to the dilaton in Eq. (2.1) by $`\alpha =5/2`$. The cosmological version of this brane is given by
$$ds_{\mathrm{D8}}^2=H^{1/8}f_{\mu \nu }dx^\mu dx^\nu +H^{9/8}e^{5\mathrm{\Phi }_1/2}dy^2,$$
(3.1)
where $`e^{\mathrm{\Phi }_2}=H^{5/4}`$ and $`\mathrm{\Phi }_1=5B/4`$.
The eleven–dimensional origin of the D8–brane is presently unclear, although Hull has shown that it can be obtained by reducing M–theory on a torus bundle over a circle in the limit where the bundle size vanishes . We now derive a solution to eleven–dimensional supergravity that can be related to a cosmological version of the D8–brane. Standard compactification of vacuum M–theory on a non–dynamical two–torus leads to a nine–dimensional theory of the form (2.1), where the field strength corresponds to that of a massless axion field, $`F_A=_A\sigma `$. The dilaton and axion parametrize the $`\mathrm{SL}(2,R)/\mathrm{U}(1)`$ coset. The existence of this non–compact global symmetry of the action implies that a generalized Scherk–Schwarz compactification on a circle may then be performed , where the axion field has a linear dependence on the circle’s coordinate. This introduces a mass parameter (cosmological constant) in the eight–dimensional theory. After a suitable rescaling of the moduli fields, the reduced action takes the form of Eq. (2.1), where $`D=8`$ and the coupling between the scalar field and $`0`$–form is given by $`\alpha =\sqrt{19/3}`$ .
Thus, the corresponding domain wall (6–brane) cosmology is of the form
$$ds_8=H^{1/6}f_{\mu \nu }dx^\mu dx^\nu +e^{\sqrt{19/3}\mathrm{\Phi }_1}H^{7/6}dy^2,$$
(3.2)
where $`H(y)=1+m|y|`$, $`\mathrm{\Phi }_2=\sqrt{19/12}\mathrm{ln}H`$ and $`m`$ is a constant representing the slope parameter of the nine–dimensional axion, $`\sigma (x^\mu ,y)=\sigma (x^\mu )+my`$. The dilaton field, $`\mathrm{\Phi }`$, is a linear combination of the three moduli fields, $`\stackrel{}{\phi }=(\phi _1,\phi _2,\phi _3)`$, originating from the diagonal components of the compactifying metric, i.e., $`\mathrm{\Phi }=\sqrt{3/19}\stackrel{}{b}_{123}.\stackrel{}{\phi }`$, where $`\stackrel{}{b}_{123}=(3/2,\sqrt{7/4},\sqrt{7/3})`$ .
Following the prescription of Lü and Pope , the solution may be oxidised back to eleven dimensions. We find that
$`ds_{11}^2=e^{\mathrm{\Phi }_1/3\alpha }f_{\mu \nu }dx^\mu dx^\nu +He^{6\mathrm{\Phi }_1/\alpha }dy^2`$ (3.3)
$`+He^{2\mathrm{\Phi }_1/\alpha }\left(dz_2^2+dz_3^2\right)`$ (3.4)
$`+H^1e^{2\mathrm{\Phi }_1/\alpha }\left(dz_1+mz_2dz_3\right)^2.`$ (3.5)
The compactifying dimensions in Eq. (3.3) form a torus bundle :
$$ds_B^2=R^2dz_2^2+\frac{1}{\mathrm{Im}\tau }\left|dz_1+\tau dz_3\right|^2,$$
(3.6)
where the $`T^2`$ fibre is spanned by the periodic coordinates $`\{z_1,z_3\}`$, $`R=H^{1/2}e^{\mathrm{\Phi }_1/\alpha }`$ is the circumference of the circular base space and $`\tau mz_2+iHe^{2\mathrm{\Phi }_1/\alpha }`$ is the complex structure. These degrees of freedom depend on both the world–volume and transverse coordinates.
We now compactify the eleven–dimensional metric (3.3) in the $`z_1`$ direction, producing a type IIA D6–brane cosmology supported by the magnetic charge, $`m`$, of the two–form, $`F_2=mdz_2dz_3`$. Conformally transforming to the string frame and employing a standard T–duality transformation in the $`z_3`$ direction leads to the corresponding D7 type IIB solution. Finally, applying the massive T–duality rules of Ref. in the $`z_2`$ direction produces a D8–brane cosmology given, in the Einstein frame, by
$`ds_{\mathrm{D8}}^2=H^{1/8}[e^{\overline{\sigma }_1/30}f_{\mu \nu }dx^\mu dx^\nu `$ (3.7)
$`+e^{\overline{\sigma }_1/10}(dz_2^2+dz_3^2)]+H^{9/8}e^{5\overline{\sigma }_1/2}dy^2,`$ (3.8)
where $`\overline{\sigma }_1=5\mathrm{\Phi }_1/(2\alpha )`$ is the world–volume dependent part of the ten–dimensional dilaton field, $`\mathrm{\Phi }_1`$ is given in Eq. (3.2) and $`f_{\mu \nu }`$ is the six–dimensional metric solving Eqs. (2.17) and (2.18), where $`B=\sqrt{19/12}\mathrm{\Phi }_1`$. The solution (3.7) has at least two abelian isometries on the world–volume and reduces to the Ricci–flat D8–brane in the limit where $`\overline{\sigma }_1`$ vanishes .
### B NS5–Brane Cosmology
Another important brane of the type IIA theory is the NS5–brane supported by the NS-NS three–form. This is coupled to the dilaton field such that $`\alpha =1`$. The corresponding NS5–brane cosmology is therefore of the form
$`ds_{\mathrm{NS5}}^2=H^{1/4}f_{\mu \nu }dx^\mu dx^\nu `$ (3.9)
$`+H^{3/4}e^{\mathrm{\Phi }_1/2}\left(dy_1^2+\mathrm{}+dy_4^2\right),`$ (3.10)
where $`e^{\mathrm{\Phi }_2}=H^{1/2}`$ and $`\mathrm{\Phi }_1=4B`$. It is well known that the type IIA theory may be derived by compactifying $`D=11`$ supergravity on a circle, where the radius of the circle is related to the string coupling by $`r_{11}=g_s^{2/3}=e^{2\mathrm{\Phi }/3}`$ . Thus, the ten–dimensional brane cosmology (3.9) may be oxidized to eleven dimensions to yield
$`ds_{\mathrm{M5}}^2=H^{1/3}\left(e^{\mathrm{\Phi }_1/6}f_{\mu \nu }dx^\mu dx^\nu \right)`$ (3.11)
$`+H^{2/3}\left(e^{2\mathrm{\Phi }_1/3}\delta _{ij}dy^idy^j+e^{4\mathrm{\Phi }_1/3}dz^2\right),`$ (3.12)
where $`z`$ is the coordinate of the eleventh dimension.
Eq. (3.11) represents a new solution to the $`D=11`$ supergravity equations of motion and may be interpreted as a M5–brane cosmology, where both the world–volume and transverse spaces are curved due to the dilaton’s dependence on the world–volume coordinates. Indeed, the transverse space is no longer conformally flat in this case. Since the eleventh dimension becomes large in the strongly coupled limit, an equivalent interpretation of this solution is given in terms of a strongly–coupled NS5–brane cosmology where the extra dimension is part of the transverse space.
The NS5–solution is also related to a D5–brane cosmology of the type IIB theory by S–duality . In type IIB supergravity, the dilaton and RR scalar field, $`\lambda `$, parametrize the $`\mathrm{SL}(2,R)/\mathrm{U}(1)`$ coset. Consequently, the theory exhibits a global $`\mathrm{SL}(2,R)`$ symmetry . The transformation is equivalent to the complex scalar field $`\kappa \lambda +ie^\mathrm{\Phi }`$ undergoing a fractional linear transformation: $`\overline{\kappa }=(A\kappa +B)/(C\kappa +D)`$, where $`ADBC=1`$. The Einstein–frame metric is a singlet under this transformation and the two–form potentials transform as a doublet. The NS5 and D5 solutions are related by the special transformation $`A=D=0`$ and $`C=B=1`$ and this relates a strongly–coupled solution to a weakly coupled one since the sign of the dilaton field is reversed. It follows, therefore, that a more general type IIB brane cosmology may be generated from a seed NS5–brane solution by applying a global $`\mathrm{SL}(2,R)`$ symmetry transformation. This produces a non–trivial scalar and 2–form potential in the RR sector. Furthermore, the dilaton field of the dual solution is given by $`e^{\overline{\mathrm{\Phi }}}=C^2e^\mathrm{\Phi }+D^2e^\mathrm{\Phi }`$ and cannot be separated into world–volume and transverse–dependent parts.
## IV Hořava–Witten Cosmology
Thus far, we have considered brane cosmologies within the context of the type II theories. However, the $`\mathrm{E}_8\times \mathrm{E}_8`$ heterotic string theory has been favoured from a phenomenological perspective and it is therefore important to discuss its cosmological consequences. The strongly coupled limit of this theory is M–theory on an orbifold, $`S^1/\mathrm{Z}_1`$, and compactification on a Calabi–Yau three–fold leads to a gauged, five–dimensional supergravity theory with two four–dimensional boundaries. For the purposes of the present discussion, it is sufficient to consider a consistent truncation of this theory that includes the breathing mode of the Calabi–Yau space, $`\mathrm{\Phi }`$, and a massless scalar field, $`\sigma `$, arising from the universal hypermultiplet. The action is given by
$`S={\displaystyle }d^5x\sqrt{|g|}[R{\displaystyle \frac{1}{2}}(\mathrm{\Phi })^2{\displaystyle \frac{1}{2}}e^\mathrm{\Phi }(\sigma )^2`$ (4.1)
$`\mathrm{\Lambda }e^{2\mathrm{\Phi }}]+{\displaystyle }_{i=1}^2(1)^i\sqrt{24}\mathrm{\Lambda }{\displaystyle }d^4x\sqrt{\left|g_i\right|}e^\mathrm{\Phi }.`$ (4.2)
The potential term in Eq. (4.1) is due to the non–trivial flux of the four–form field strength on four–cycles of the Calabi–Yau space and it supports a solitonic 3–brane (domain wall) solution . This 3–brane has an eleven–dimensional interpretation in terms of 5–branes that are located on the ten–dimensional orbifold planes, where two of the dimensions are wrapped around a Calabi–Yau two–cycle.
Cosmological brane solutions in Hořava–Witten theory have been found previously for a trivial axion field . The five–dimensional metric is given by
$$ds_{\mathrm{HW}}^2=Hf_{\mu \nu }dx^\mu dx^\nu +H^4e^{2B}dy^2,$$
(4.3)
where $`H=1+(2\mathrm{\Lambda }/3)^{1/2}|y|`$, the breathing mode, $`\mathrm{\Phi }=\mathrm{\Phi }_1(x)+\mathrm{\Phi }_2(y)`$, is given by
$$\mathrm{\Phi }_1=B,\mathrm{\Phi }_2=3\mathrm{ln}H$$
(4.4)
and the world–volume metric, $`f_{\mu \nu }`$, is determined by the effective field equations
$`\overline{R}_{\mu \nu }=\overline{}_{\mu \nu }B+{\displaystyle \frac{3}{2}}\overline{}_\mu B\overline{}_\nu B`$ (4.5)
$`\overline{}^2B+\left(\overline{}B\right)^2=0.`$ (4.6)
We now consider the effects of introducing the axion field, $`\sigma `$. We assume that the metric and breathing mode are given by Eqs. (4.3) and (4.4), respectively. The $`(\mu y)`$–components of the Einstein field equations are still solved by the separable ansatz (4.4) if the axion field is independent of the world–volume coordinates. We therefore assume that it depends only on the transverse dimension, $`y`$. Its field equation then admits the first integral
$$\sigma ^{}=AH^3,$$
(4.7)
where a prime denotes differentiation with respect to $`y`$ and $`A`$ is an arbitrary constant of integration. The $`(\mu \nu )`$–components of the Einstein field equations are solved as before provided that the Ricci tensor of the world–volume satisfies Eq. (4.5). However, the equation of motion for the breathing mode acquires an additional term due to the axion field. It is solved if
$$\overline{}^2B+\left(\overline{}B\right)^2=\frac{A^2}{2}e^{3B}$$
(4.8)
and it can be shown that the $`(yy)`$–component of the Einstein equations is also solved if Eq. (4.8) is satisfied.
Hence, the compactified heterotic M–theory action (4.1) admits a curved domain wall cosmology of the form given by Eq. (4.3), where the axion field satisfies Eq. (4.7). The cosmological expansion of the brane is determined by the conditions (4.5) and (4.8). Performing the conformal transformation
$$\stackrel{~}{f}_{\mu \nu }=\mathrm{\Theta }^2f_{\mu \nu },\mathrm{\Theta }^2e^B$$
(4.9)
and field redefinition $`B=\chi /2`$ implies that these conditions are equivalent to
$`\stackrel{~}{R}_{\mu \nu }={\displaystyle \frac{1}{2}}\stackrel{~}{}_\mu \chi \stackrel{~}{}_\nu \chi +{\displaystyle \frac{A^2}{4}}\stackrel{~}{f}_{\mu \nu }e^{2\chi }`$ (4.10)
$`\stackrel{~}{}^2\chi =A^2e^{2\chi }.`$ (4.11)
Eqs. (4.10) and (4.11) may be interpreted as the four–dimensional field equations for a minimally coupled scalar field, $`\chi `$, that self–interacts through an exponential potential $`V=(A^2/2)e^{Q\chi }`$, where $`Q=2`$. The momentum of the axion field in the orbifold direction manifests itself to an observer on the brane as a self–interaction potential for the breathing mode of the Calabi–Yau space . For an exponential potential of this type, the late–time attractor for the spatially flat Friedmann–Robertson–Walker (FRW) cosmology is a power law, $`at^{1/Q^2}`$, for $`Q^23`$, otherwise it is $`at^{1/3}`$ . In this Hořava–Witten model, $`Q^2=4`$, and the latter situation therefore arises. Hence, the unique late–time attractor is non–inflationary, although it is interesting that a potential for the breathing mode can be generated in this fashion.
## V Discussion and Conclusion
The solitonic D$`p`$– and M–branes of string and M–theory have played a central role in establishing the duality relationships that exist between the different theories. A necessary condition for a brane to be interpreted cosmologically is that its world–volume should be non–static and curved due to the existence of matter fields varying dynamically on the brane. We have found that at the level of the supergravity field equations, the world–volume of many of these branes becomes curved when the dilaton has a non–trivial dependence on the world–volume coordinates and is related to the transverse dimensions in an appropriate way. In particular, we have presented a cosmological version of the M5–brane, where both the world–volume and transverse spaces are curved. This solution represents the strongly–coupled limit of an NS5–brane cosmology. An eleven–dimensional interpretation was also given for a cosmological D8–brane of the massive type IIA theory. Finally, a class of strongly–coupled braneworlds was found in heterotic M–theory compactified on a Calabi–Yau space.
Moreover, the geometry of the brane world–volume was kept arbitrary in the analysis and was not restricted to the spatially isotropic FRW metrics. This is important since the effects of spatial anisotropy and inhomogeneity may have been significant in the very early universe. The problem of solving the field equations (2.2)–(2.5) was reduced to solving Einstein gravity minimally coupled to a massless scalar field and this system has been extensively studied in the literature.
We conclude by considering the possibility that the kinetic energy of the dilaton field can drive an epoch of inflationary expansion on the D$`p`$–branes. In the standard, pre–big bang scenario, the simplest solution is that of the spatially flat, homogeneous Bianchi I model defined over $`t<0`$. (For a review, see, e.g., Ref. ). This is the time–reversal of the ‘rolling radii’ solution of Mueller and represents a generalization of the Kasner solution . The string frame metric and dilaton field are given by
$$ds^2=dt^2+\underset{i=1}{\overset{9}{}}(t)^{2\beta _i}dz_idz^i$$
(5.1)
and
$$\mathrm{\Phi }=\frac{1}{2}\left(1\underset{i=1}{\overset{9}{}}\beta _i\right)\mathrm{ln}(t),$$
(5.2)
respectively, where the constants, $`\{\beta _i\}`$, satisfy the constraint equation
$$\underset{i=1}{\overset{9}{}}\beta _i^2=1.$$
(5.3)
Given the nature of Eq. (2.24), we consider D$`p`$–brane cosmologies of the form
$`ds^2=H^{1/2}\left(dt^2+{\displaystyle \underset{i=1}{\overset{d_W1}{}}}(t)^{2\beta _i}dx_i^2\right)`$ (5.4)
$`+H^{1/2}(t)^{2\gamma }{\displaystyle \underset{j=1}{\overset{d_T}{}}}dy_j^2,`$ (5.5)
where $`\mathrm{\Phi }_1`$ depends only on cosmic time. The constraints on the exponents $`\{\beta _i,\gamma \}`$ may be deduced by noting that the $`q`$–form field strength supporting the brane is non–trivial only in the transverse–dependent sector of the field equations (2.2)–(2.5). Thus, the time–dependence of the metric and dilaton is given by the rolling radii solution (5.1) and (5.2):
$`\mathrm{\Phi }_1={\displaystyle \frac{1}{2}}\left(1{\displaystyle \underset{i=1}{\overset{d_W1}{}}}\beta _id_T\gamma \right)\mathrm{ln}(t)`$ (5.6)
$`{\displaystyle \underset{i=1}{\overset{d_W1}{}}}\beta _i^2+d_T\gamma ^2=1.`$ (5.7)
However, there is an additional constraint on the exponents because the dilaton field is directly related to the transverse dimensions by the separability condition (2.16). Comparison of Eqs. (2.24) and (5.4) implies that
$$1\underset{i=1}{\overset{d_W1}{}}\beta _i=(3q)\gamma $$
(5.8)
and thus, for $`q3`$,
$$\underset{i=1}{\overset{d_W1}{}}\beta _i^2+\frac{q+1}{(q3)^2}\left(1\underset{i=1}{\overset{d_W1}{}}\beta _i\right)^2=1.$$
(5.9)
During dilaton–driven inflation, the string coupling increases in value. Eqs. (5.6) and (5.8) imply that a necessary condition for inflation is $`\gamma (q1)<0`$ and there is a wide region of parameter space where inflation of this type can proceed on the brane. For example, let us consider the D8–brane cosmology (3.1) of the massive type IIA theory and, for simplicity, assume that five of the world–volume dimensions are independent of time and that the remaining three are isotropic, $`\beta _i=\beta `$. Eq. (5.9) then implies that $`\beta =(1\pm \sqrt{33})/12`$ and $`\gamma =(3\sqrt{33})/12`$ and the negative root therefore leads to accelerated expansion as $`t0^{}`$.
Finally, we remark that the Ricci–flat branes (2.6)–(2.8) are also directly relevant to cosmology as a consequence of a powerful embedding theorem due to Campbell . This theorem states that any analytic Riemannian space of dimension $`n`$ and signature $`(1,n1)`$ can be locally and isometrically embedded in a Ricci–flat, Riemannian space of dimension $`n+1`$ and signature $`(1,n)`$. The embedding is established by solving a set of constraint equations that are compatible with the Gauss–Codazzi equations . In particular, perfect fluid FRW cosmologies can be embedded in this fashion . ¿From a five–dimensional point of view, the solution is interpreted as a shock wave travelling through time and the fifth dimension and the non–trivial energy–momentum tensor is induced on the four–dimensional hypersurface by relaxing the cylinder condition of Kaluza–Klein theory . Since Campbell’s theorem is independent of the dimensionality of the space, the procedure may be repeated an arbitrary number of times to embedded four–dimensional cosmologies in Ricci–flat branes of higher dimensions, such as the M5–brane. It would be interesting to explore such embeddings further.
###### Acknowledgements.
The author is supported by the Royal Society. We thank J. Gauntlett for helpful comments.
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# Implicit Integration of the Time-Dependent Ginzburg–Landau Equations of Superconductivity
## 1 Introduction
At the macroscopic level, the state of a superconductor can be described in terms of a complex-valued order parameter and a real vector potential. These variables, which determine the superconducting and electromagnetic properties of the system at equilibrium, are found as solutions of the Ginzburg–Landau (GL) equations of superconductivity. Because they correspond to critical points of the GL energy functional , they can, at least in principle, be determined by minimizing a functional. In practice, one introduces a time-like variable and computes equilibrium states by integrating the time-dependent Ginzburg–Landau (TDGL) equations. The TDGL equations, first formulated by Schmid and subsequently derived from microscopic principles by Gor’kov and Éliashberg , are nontrivial generalizations of the (time-independent) GL equations, as the time rate of change must be introduced in such a manner that gauge invariance is preserved at all times. The TDGL equations have been analyzed by several authors; see, for example, the articles and the references cited there.
We are interested, in particular, in vortex solutions of the GL equations. These are singular solutions, where the phase of the order parameter changes by $`2\pi `$ along any closed contour surrounding a vortex point. Vortices are of critical importance in technological applications of superconductivity.
Computing vortex solutions of the GL equations by integrating the TDGL equations to equilibrium has the advantage that the solutions thus found are stable. At the same time, one obtains information about the transient behavior of the system. Integrating the TDGL equations to equilibrium is, however, a time-consuming process requiring considerable computing resources. In simulations of vortex dynamics in superconductors, which were performed on an IBM SP with tens of processors in parallel, using a simple one-step Euler integration procedure, we routinely experienced equilibration times on the order of one hundred hours . Incremental changes would gradually drive the system to lower energy levels. These very long equilibration times arise, of course, because we are dealing with large physical systems undergoing a phase transition. The energy landscape for such systems is a broad, gently undulating plain with many shallow local minima. It is therefore important to develop efficient integration techniques that remain stable and accurate as the time step increases.
In this article we present four integration techniques ranging from fully explicit to fully implicit for problems on rectangular domains in two dimensions. These two-dimensional domains should be viewed as cross sections of three-dimensional systems that are infinite and homogeneous in the direction of the magnetic field, which is orthogonal to the plane of the cross section. The algorithms are scalable in a multiprocessing environment and generalize to three dimensions. We evaluate the performance of each algorithm on the same benchmark problem, namely, the equilibration of a vortex configuration in a system consisting of a superconducting core embedded in a blanket of insulating material (air) and undergoing a transition from the Meissner state to the vortex state under the influence of an externally applied magnetic field. We determine the maximum allowable time step for stability, the number of time steps needed to reach the equilibrium configuration, and the CPU cost per time step.
Different algorithms correspond to different dynamics through state space, so the eventual equilibrium vortex configuration may differ from one algorithm to another. Hence, once we have the equilibrium configurations, we need some measure to assess their accuracy. For this purpose we use three parameters: the number of vortices, the mean intervortex distance (bond length), and the mean bond angle taken over nearest-neighbor pairs of bonds. When each of these parameters differs less than a specified tolerance, we say that the corresponding vortex configurations are the same.
Our investigations show that one can increase the time step by almost two orders of magnitude, without losing stability, by going from the fully explicit to the fully implicit algorithm. The fully implicit algorithm has a higher cost per time step, but the wall clock time needed to compute the equilibrium solution (the most important measure for practical purposes) is still significantly less. The wall clock time can be reduced further by using a multi-timestepping procedure.
In §2, we present the Ginzburg–Landau model of superconductivity, first in its formulation as a system of partial differential equations, then as a system of ordinary differential equations after the spatial derivatives have been approximated by finite differences. In §3, we give four algorithms to integrate the system of ordinary equations: a fully explicit, a semi-implit, an implicit, and a fully implicit algorithm. In §4, we present and evaluate the results of the investigation. In §5, we further evaluate the fully implicit algorithm from the point of view of parallelism and multi-timestepping. The conclusions are summarized in §6.
## 2 Ginzburg–Landau Model
The time-dependent Ginzburg–Landau (TDGL) equations of superconductivity are two coupled partial differential equations for the complex-valued order parameter $`\psi =|\psi |\mathrm{e}^{i\varphi }`$ and the real vector-valued vector potential $`𝐀`$,
$`{\displaystyle \frac{\mathrm{}^2}{2m_sD}}\left({\displaystyle \frac{}{t}}+{\displaystyle \frac{ie_s}{\mathrm{}}}\mathrm{\Phi }\right)\psi `$ $`=`$ $`{\displaystyle \frac{1}{2m_s}}\left({\displaystyle \frac{\mathrm{}}{i}}{\displaystyle \frac{e_s}{c}}𝐀\right)^2\psi +a\psi b|\psi |^2\psi ,`$ (1)
$`\nu \left({\displaystyle \frac{1}{c}}{\displaystyle \frac{𝐀}{t}}+\mathrm{\Phi }\right)`$ $`=`$ $`{\displaystyle \frac{c}{4\pi }}\times \times 𝐀+𝐉_s.`$ (2)
Here, $`𝐉_s`$ is the supercurrent density, which is a nonlinear function of $`\psi `$ and $`𝐀`$,
$$𝐉_s=\frac{e_s\mathrm{}}{2im_s}(\psi ^{}\psi \psi \psi ^{})\frac{e_s^2}{m_sc}|\psi |^2𝐀=\frac{e_s}{m_s}|\psi |^2\left(\mathrm{}\varphi \frac{e_s}{c}𝐀\right).$$
(3)
The real scalar-valued electric potential $`\mathrm{\Phi }`$ is a diagnostic variable. The constants in the equations are $`\mathrm{}`$, Planck’s constant divided by $`2\pi `$; $`a`$ and $`b`$, two positive constants; $`c`$, the speed of light; $`m_s`$ and $`e_s`$, the effective mass and charge, respectively, of the superconducting charge carriers (Cooper pairs); $`\nu `$, the electrical conductivity; and $`D`$, the diffusion coefficient. As usual, $`i`$ is the imaginary unit, and denotes complex conjugation.
The quantity $`|\psi |^2`$ represents the local density of Cooper pairs. The local time rate of change $`_t𝐀`$ of $`𝐀`$ determines the electric field, $`𝐄=(1/c)_t𝐀+\mathrm{\Phi }`$, its spatial variation the (induced) magnetic field, $`𝐁=\times 𝐀`$.
The TDGL equations describe the gradient flow for the Ginzburg–Landau energy functional. This functional is zero in the normal state, when $`\psi =0`$ and the externally applied magnetic field penetrates the superconductor everywhere, $`\times 𝐀=𝐇`$. In the superconducting state, it is given by the expression
$$E=\left[\frac{1}{2m_s}\left|\left(\frac{\mathrm{}}{i}\frac{e_s}{c}𝐀\right)\psi \right|^2+\left(a|\psi |^2+\frac{b}{2}|\psi |^4\right)+|\times 𝐀𝐇|^2\right]dx.$$
(4)
The three terms represent the kinetic energy, the condensation energy, and the field energy, respectively. A thermodynamic equilibrium configuration corresponds to a critical point of $`E`$.
The energy functional (4) assumes that there are no defects in the superconductor. Material defects can be naturally present or artifically induced and can be in the form of point, planar, or columnar defects (quenched disorder). A material defect results in a local reduction of the depth of the well of the condensation energy. A simple way to include material defects in the Ginzburg–Landau model is by assuming that the parameter $`a`$ depends on position and has a smaller value wherever a defect is present.
### 2.1 Dimensionless Form
Let $`\psi _{\mathrm{}}^2=a/b`$, and let $`\lambda `$, $`\xi `$, and $`H_c`$ denote the London penetration depth, the coherence length, and the thermodynamic critical field, respectively,
$$\lambda =\left(\frac{m_sc^2}{4\pi \psi _{\mathrm{}}^2e_s^2}\right)^{1/2},\xi =\left(\frac{\mathrm{}^2}{2m_sa}\right)^{1/2},H_c=(4\pi a\psi _{\mathrm{}}^2)^{1/2}.$$
(5)
In this study, we render the TDGL equations dimensionless by measuring lengths in units of $`\xi `$, time in units of the relaxation time $`\xi ^2/D`$, fields in units of $`H_c\sqrt{2}`$, and energy densities in units of $`(1/4\pi )H_c^2`$. The nondimensional TDGL equations are
$`\left({\displaystyle \frac{}{t}}+i\mathrm{\Phi }\right)\psi `$ $`=`$ $`\left({\displaystyle \frac{i}{\kappa }}𝐀\right)^2\psi +\tau \psi |\psi |^2\psi ,`$ (6)
$`\sigma \left({\displaystyle \frac{𝐀}{t}}+\kappa \mathrm{\Phi }\right)`$ $`=`$ $`\times \times 𝐀+𝐉_s,`$ (7)
where
$$𝐉_s=\frac{1}{2i\kappa }(\psi ^{}\psi \psi \psi ^{})\frac{1}{\kappa ^2}|\psi |^2𝐀=\frac{1}{\kappa }|\psi |^2\left(\varphi \frac{1}{\kappa }𝐀\right).$$
(8)
Here, $`\kappa =\lambda /\xi `$ is the Ginzburg–Landau parameter and $`\sigma `$ is a dimensionless resistivity, $`\sigma =(4\pi D/c^2)\nu `$. The coefficient $`\tau `$ has been inserted to account for defects; $`\tau (x)<1`$ if $`x`$ is in a defective region; otherwise $`\tau (x)=1`$. The nondimensional TDGL equations are associated with the dimensionless energy functional
$$E=\left[\left|\left(\frac{i}{\kappa }𝐀\right)\psi \right|^2+\left(\tau |\psi |^2+\frac{1}{2}|\psi |^4\right)+|\times 𝐀𝐇|^2\right]dx.$$
(9)
### 2.2 Gauge Choice
The (nondimensional) TDGL equations are invariant under a gauge transformation,
$$𝒢_\chi :(\psi ,𝐀,\mathrm{\Phi })(\psi \mathrm{e}^{i\chi },𝐀+\kappa \chi ,\mathrm{\Phi }_t\chi ).$$
(10)
Here, $`\chi `$ can be any real scalar-valued function of position and time. We maintain the zero-electric potential gauge, $`\mathrm{\Phi }=0`$, at all times, using the link variable $`𝐔`$,
$$𝐔=\mathrm{exp}\left(\frac{i}{\kappa }𝐀\right).$$
(11)
This definition is componentwise: $`U_x=\mathrm{exp}(i\kappa ^1^xA_x(x^{},y,z)dx^{})`$, $`\mathrm{}`$. The gauged TDGL equations can now be written in the form
$`{\displaystyle \frac{\psi }{t}}`$ $`=`$ $`{\displaystyle \underset{\mu =x,y,z}{}}U_\mu ^{}{\displaystyle \frac{^2}{\mu ^2}}(U_\mu \psi )+\tau \psi |\psi |^2\psi ,`$ (12)
$`\sigma {\displaystyle \frac{𝐀}{t}}`$ $`=`$ $`\times \times 𝐀+𝐉_s,`$ (13)
where
$$J_{s,\mu }=\frac{1}{\kappa }\text{ Im }\left[(U_\mu \psi )^{}\frac{}{\mu }(U_\mu \psi )\right],\mu =x,y,z.$$
(14)
### 2.3 Two-Dimensional Problems
From here on we restrict the discussion to problems on a two-dimensional rectangular domain (coordinates $`x`$ and $`y`$), assuming boundedness in the $`x`$ direction and periodicity in the $`y`$ direction. The domain represents a superconducting core surrounded by a blanket of insulating material (air) or a normal metal. The order parameter vanishes outside the superconductor, and no superconducting charge carriers leave the superconductor. The whole system is driven by a time-independent externally applied magnetic field $`𝐇`$ that is parallel to the $`z`$ axis, $`𝐇=(0,0,H)`$. The vector potential and the supercurrent have two nonzero components, $`𝐀=(A_x,A_y,0)`$ and $`𝐉_s=(J_x,J_y,0)`$, while the magnetic field has only one nonzero component, $`𝐁=(0,0,B)`$, where $`B=_xA_y_yA_x`$.
### 2.4 Spatial Discretization
The physical configuration to be modeled (superconductor embedded in blanket material) is periodic in $`y`$ and bounded in $`x`$. In the $`x`$ direction, we distinguish three subdomains: an interior subdomain occupied by the superconducting material and two subdomains, one on either side, occupied by the blanket material. We take the two blanket layers to be equally thick, but do not assume that the problem is symmetric around the midplane. (Possible sources of asymmetry are material defects in the system, surface currents, and different field strengths on the two outer surfaces.)
We impose a regular grid with mesh widths $`h_x`$ and $`h_y`$,
$$\mathrm{\Omega }_{i,j}=(x_i,x_{i+1})\times (y_j,y_{j+1}),x_i=x_0+ih_x;y_j=y_0+jh_y,$$
(15)
assuming the following correspondences:
| Left outer surface: | $`x=x_0+\frac{1}{2}h_x`$, | $`i=0`$, |
| --- | --- | --- |
| Left interface: | $`x=x_{n_{sx}1}+\frac{1}{2}h_x`$, | $`i=n_{sx}1`$, |
| Right interface: | $`x=x_{n_{ex}}+\frac{1}{2}h_x`$, | $`i=n_{ex}`$, |
| Right outer surface: | $`x=x_{n_x}+\frac{1}{2}h_x`$, | $`i=n_x`$. |
One period in the $`y`$ direction is covered by the points $`j=1,\mathrm{},n_y`$. We use the symbols Sc and Bl to denote the index sets for the superconducting and blanket region, respectively,
Sc $`=`$ $`\{(i,j):(i,j)[n_{sx},n_{ex}]\times [1,n_y]\},`$ (16)
Bl $`=`$ $`\{(i,j):(i,j)[1,n_{sx}1][n_{ex}+1,n_x]\times [1,n_y]\}.`$ (17)
The order parameter $`\psi `$ is evaluated at the grid vertices,
$$\psi _{i,j}=\psi (x_i,y_j),(i,j)\text{Sc},$$
(18)
the components $`A_x`$ and $`A_y`$ of the vector potential at the midpoints of the respective edges,
$$A_{x;i,j}=A_x(x_i+\frac{1}{2}h_x,y_j),A_{y;i,j}=A_y(x_i,y_j+\frac{1}{2}h_y),(i,j)\text{Sc}\text{Bl},$$
(19)
and the induced magnetic field $`B`$ at the center of a grid cell,
$`B_{i,j}`$ $`=`$ $`B(x_i+\frac{1}{2}h_x,y_j+\frac{1}{2}h_y)`$
$`=`$ $`{\displaystyle \frac{A_{y;i+1,j}A_{y;i,j}}{h_x}}{\displaystyle \frac{A_{x;i,j+1}A_{x;i,j}}{h_y}},(i,j)\text{Sc}\text{Bl},`$
see Fig. 1.
The values of the link variables and the supercurrent are computed from the expressions
$`U_{x;i,j}=\mathrm{e}^{i\kappa ^1h_xA_{x;i,j}}`$ , $`U_{y;i,j}=\mathrm{e}^{i\kappa ^1h_yA_{y;i,j}},`$ (21)
$`J_{x;i,j}={\displaystyle \frac{1}{\kappa h_x}}\mathrm{Im}\left[\psi _{i,j}^{}U_{x;i,j}\psi _{i+1,j}\right]`$ , $`J_{y;i,j}={\displaystyle \frac{1}{\kappa h_y}}\mathrm{Im}\left[\psi _{i,j}^{}U_{y;i,j}\psi _{i,j+1}\right].`$ (22)
The discretized TDGL equations are
$`{\displaystyle \frac{\mathrm{d}\psi _{i,j}}{\mathrm{d}t}}`$ $`=`$ $`\left(L_{xx}(U_{x;,j})\psi _{,j}\right)_i+\left(L_{yy}(U_{y;i,})\psi _{i,}\right)_j+N\left(\psi _{i,j}\right),(i,j)\text{Sc},`$ (23)
$`\sigma {\displaystyle \frac{\mathrm{d}A_{x;i,j}}{\mathrm{d}t}}`$ $`=`$ $`\left(D_{yy}A_{x;i,}\right)_j\left(D_{yx}A_{y;,}\right)_{i,j}+J_{x;i,j},(i,j)\text{Sc}\text{Bl},`$ (24)
$`\sigma {\displaystyle \frac{\mathrm{d}A_{y;i,j}}{\mathrm{d}t}}`$ $`=`$ $`\left(D_{xx}A_{y;,j}\right)_i\left(D_{xy}A_{x;,}\right)_{i,j}+J_{y;i,j}.(i,j)\text{Sc}\text{Bl},`$ (25)
where
$`\left(L_{xx}(U_{x;,j})\psi _{,j}\right)_i`$ $`=`$ $`h_x^2\left[U_{x;i,j}\psi _{i+1,j}2\psi _{i,j}+U_{x;i1,j}^{}\psi _{i1,j}\right],`$ (26)
$`\left(L_{yy}(U_{y;i,})\psi _{i,}\right)_j`$ $`=`$ $`h_y^2\left[U_{y;i,j}\psi _{i,j+1}2\psi _{i,j}+U_{y;i,j1}^{}\psi _{i,j1}\right],`$ (27)
$`N\left(\psi _{i,j}\right)`$ $`=`$ $`\tau _{i,j}\psi _{i,j}|\psi _{i,j}|^2\psi _{i,j},`$ (28)
$`\left(D_{yy}A_{x;i,}\right)_j`$ $`=`$ $`h_y^2\left[A_{x;i,j+1}2A_{x;i,j}+A_{x;i,j1}\right],`$ (29)
$`\left(D_{xx}A_{y;,j}\right)_i`$ $`=`$ $`h_x^2\left[A_{y;i+1,j}2A_{y;i,j}+A_{y;i1,j}\right],`$ (30)
$`\left(D_{yx}A_{y;,}\right)_{i,j}`$ $`=`$ $`h_x^1h_y^1\left[\left(A_{y;i+1,j}A_{y;i,j}\right)\left(A_{y;i+1,j1}A_{y;i,j1}\right)\right],`$ (31)
$`\left(D_{xy}A_{x;,}\right)_{i,j}`$ $`=`$ $`h_x^1h_y^1\left[\left(A_{x;i,j+1}A_{x;i,j}\right)\left(A_{x;i1,j+1}A_{x;i1,j}\right)\right].`$ (32)
The interface conditions are
$$\psi _{n_{sx}1,j}=U_{x;n_{sx}1,j}\psi _{n_{sx},j},\psi _{n_{ex}+1,j}=U_{x;n_{ex},j}^{}\psi _{n_{ex},j},j=1,\mathrm{},n_y.$$
(33)
At the outer boundary, $`B`$ is given,
$$B_{0,j}=H_{L_j},B_{n_x,j}=H_{R_j},j=1,\mathrm{},n_y.$$
(34)
The resulting approximation is second-order accurate .
## 3 Time Integration
We now address the integration of Eqs. (23)–(25). The first equation, which controls the evolution of $`\psi `$, involves the second-order linear finite-difference operators $`L_{xx}`$ and $`L_{yy}`$, whose coefficients depend on $`A_x`$ and $`A_y`$, and the local nonlinear operator $`N`$, which involves neither $`A_x`$ nor $`A_y`$. Each of the other two equations, which control the evolution of $`A_x`$ and $`A_y`$ respectively, involves likewise a second-order linear finite-difference operator, but with constant coefficients, and the nonlinear supercurrent operator, which involves $`\psi `$, $`A_x`$, and $`A_y`$. The following algorithms are distinguished by whether the various operators are treated explicitly or implicitly.
### 3.1 Fully Explicit Integration
Algorithm I uses a fully explicit forward Euler time-marching procedure for $`\psi `$, $`A_x`$, and $`A_y`$. Starting from an initial triple $`(\psi ^0,A_x^0,A_y^0)`$, we solve for $`n=0,1,\mathrm{}`$,
$`{\displaystyle \frac{\psi _{i,j}^{n+1}\psi _{i,j}^n}{\mathrm{\Delta }t}}`$ $`=`$ $`\left(L_{xx}(U_{x;,j}^n)\psi _{,j}^n\right)_i+\left(L_{yy}(U_{y;i,}^n)\psi _{i,}^n\right)_j+N\left(\psi _{i,j}^n\right),(i,j)\text{Sc},`$ (35)
$`\sigma {\displaystyle \frac{A_{x;i,j}^{n+1}A_{x;i,j}^n}{\mathrm{\Delta }t}}`$ $`=`$ $`\left(D_{yy}A_{x;i,}^n\right)_j\left(D_{yx}A_{y;,}^n\right)_{i,j}+J_{x;i,j}^n,(i,j)\text{Sc}\text{Bl},`$ (36)
$`\sigma {\displaystyle \frac{A_{y;i,j}^{n+1}A_{y;i,j}^n}{\mathrm{\Delta }t}}`$ $`=`$ $`\left(D_{xx}A_{y;,j}^n\right)_i\left(D_{xy}A_{x;,}^n\right)_{i,j}+J_{y;i,j}^n.(i,j)\text{Sc}\text{Bl},`$ (37)
where $`J^n`$ is defined in terms of $`\psi ^n`$, $`A_x^n`$, and $`A_y^n`$ in the obvious way. The initial triple is usually chosen so the superconductor is in the Meissner state, with a seed present to trigger the transition to the vortex state.
Algorithm I has been described in . It has been implemented in a distributed-memory multiprocessor environment (IBM SP2); the transformations necessary to achieve the parallelism have been described in . The code uses the Message Passing Interface (MPI) standard as implemented in the MPICH software library for domain decomposition, interprocessor communication, and file I/O. The code has been used extensively to study vortex dynamics in superconducting media . The underlying algorithm provides highly accurate solutions but requires a significant number of time steps for equilibration. For stability reasons, the time step $`\mathrm{\Delta }t`$ cannot exceed 0.0025.
### 3.2 Semi-Implicit Integration
Algorithm II is generated by an implicit treatment of the second-order linear finite-difference operators $`D_{yy}`$ and $`D_{xx}`$ in the equations for $`A_x`$ and $`A_y`$, respectively,
$`{\displaystyle \frac{\psi _{i,j}^{n+1}\psi _{i,j}^n}{\mathrm{\Delta }t}}`$ $`=`$ $`\left(L_{xx}(U_{x;,j}^n)\psi _{,j}^n\right)_i+\left(L_{yy}(U_{y;i,}^n)\psi _{i,}^n\right)_j+N\left(\psi _{i,j}^n\right),(i,j)\text{Sc},`$ (38)
$`\sigma {\displaystyle \frac{A_{x;i,j}^{n+1}A_{x;i,j}^n}{\mathrm{\Delta }t}}`$ $`=`$ $`\left(D_{yy}A_{x;i,}^{n+1}\right)_j\left(D_{yx}A_{y;,}^n\right)_{i,j}+J_{x;i,j}^n,(i,j)\text{Sc}\text{Bl},`$ (39)
$`\sigma {\displaystyle \frac{A_{y;i,j}^{n+1}A_{y;i,j}^n}{\mathrm{\Delta }t}}`$ $`=`$ $`\left(D_{xx}A_{y;,j}^{n+1}\right)_i\left(D_{xy}A_{x;,}^n\right)_{i,j}+J_{y;i,j}^n.(i,j)\text{Sc}\text{Bl}.`$ (40)
Equations (39) and (40) lead to two linear systems of equations,
$`\left(I{\displaystyle \frac{\mathrm{\Delta }t}{\sigma }}D_{yy}\right)A_{x;i}^{n+1}`$ $`=`$ $`F_i(\psi ^n,A_x^n,A_y^n),i=1,\mathrm{},n_x1,`$ (41)
$`\left(I{\displaystyle \frac{\mathrm{\Delta }t}{\sigma }}D_{xx}\right)A_{y;j}^{n+1}`$ $`=`$ $`G_j(\psi ^n,A_x^n,A_y^n),j=1,\mathrm{},n_y,`$ (42)
for the vectors of unknowns $`A_{x;i}=\{A_{x;i,j}:j=1,\mathrm{},n_y\}`$ and $`A_{y;j}=\{A_{y;i,j}:i=1,\mathrm{},n_x1\}`$. The matrix $`D_{yy}`$ has dimension $`n_y\times n_y`$ and is periodic tridiagonal with elements $`h_y^2,2h_y^2,h_y^2`$; the matrix $`D_{xx}`$ has dimension $`(n_x1)\times (n_x1)`$ and is tridiagonal with elements $`h_x^2,2h_x^2,h_x^2`$, (except along the edges, because of the boundary conditions). Both matrices are independent of $`i`$ and $`j`$. Furthermore, if the boundary conditions are time independent, they are constant throughout the time-stepping process. Hence, the coefficient matrices in Eqs. (41) and (42) need to be factored only once; in fact, the factorization can be done in the preprocessing stage and the factors can be stored.
In a parallel processing environment, the coefficient matrices extend over several processors, so Eqs. (41) and (42) are broken up in blocks corresponding to the manner in which the computational mesh is distributed among the processor set. We first solve the equations within each processor (inner iterations) and then couple the solutions across processor boundaries (outer iterations). Hence, we deal with interprocessor coupling in an iterative fashion. Two to three inner iterations usually suffice to reach a desired tolerance for convergence. After each inner iteration, each processor shares boundary data with its neighbors through MPI calls.
### 3.3 Implicit Integration
Algorithm III combines the semi-implicit treatment of $`A_x`$ and $`A_y`$ with an implicit treatment of the order parameter,
$`{\displaystyle \frac{\psi _{i,j}^{n+1}\psi _{i,j}^n}{\mathrm{\Delta }t}}`$ $`=`$ $`\left(L_{xx}(U_{x;,j}^n)\psi _{,j}^{n+1}\right)_i+\left(L_{yy}(U_{y;i,}^n)\psi _{i,}^{n+1}\right)_j+N\left(\psi _{i,j}^n\right),(i,j)\text{Sc},`$ (43)
$`\sigma {\displaystyle \frac{A_{x;i,j}^{n+1}A_{x;i,j}^n}{\mathrm{\Delta }t}}`$ $`=`$ $`\left(D_{yy}A_{x;i,}^{n+1}\right)_j\left(D_{yx}A_{y;,}^n\right)_{i,j}+J_{x;i,j}^n,(i,j)\text{Sc}\text{Bl},`$ (44)
$`\sigma {\displaystyle \frac{A_{y;i,j}^{n+1}A_{y;i,j}^n}{\mathrm{\Delta }t}}`$ $`=`$ $`\left(D_{xx}A_{y;,j}^{n+1}\right)_i\left(D_{xy}A_{x;,}^n\right)_{i,j}+J_{y;i,j}^n.(i,j)\text{Sc}\text{Bl}.`$ (45)
The second and third equation are solved as in the semi-implicit algorithm of the preceding section. The first equation is solved by a method similar to the method of Douglas and Gunn for the Laplacian.
We begin by transforming Eq. (43) into an equation for the correction matrix $`\varphi ^{n+1}=\psi ^{n+1}\psi ^n`$. The equation has the general form
$$\left(I\mathrm{\Delta }t(L_{xx}+L_{yy})\right)\varphi ^{n+1}=F(\psi ^n,A_x^n,A_y^n).$$
(46)
If $`\mathrm{\Delta }t`$ is sufficiently small, we may replace the operator in the left member by an approximate factorization,
$$\left(I\mathrm{\Delta }t(L_{xx}+L_{yy})\right)\left(I\mathrm{\Delta }tL_{xx}\right)\left(I\mathrm{\Delta }tL_{yy}\right),$$
(47)
and consider, instead of Eq. (46),
$$\left(I\mathrm{\Delta }tL_{xx}\right)\left(I\mathrm{\Delta }tL_{yy}\right)\varphi ^{n+1}=F(\psi ^n,A_x^n,A_y^n).$$
(48)
This equation can be solved in two steps,
$`\left(I\mathrm{\Delta }tL_{xx}\right)\phi `$ $`=`$ $`F,`$ (49)
$`\left(I\mathrm{\Delta }tL_{yy}\right)\varphi ^{n+1}`$ $`=`$ $`\phi .`$ (50)
The conditions (33), which must be satisfied at the interface between the superconductor and the blanket material, require some care. If we impose the conditions at every time step, then
$`\varphi _{n_{sx}1,j}^{n+1}`$ $`=`$ $`U_{x;n_{sx}1,j}^{n+1}\varphi _{n_{sx},j}^{n+1}+\left[U_{x;n_{sx}1,j}^{n+1}U_{x;n_{sx}1,j}^n\right]\psi _{n_{sx},j}^n,`$
$`\varphi _{n_{ex}+1,j}^{n+1}`$ $`=`$ $`\left(U_{x;n_{ex},j}^{n+1}\right)^{}\varphi _{n_{ex},j}^{n+1}+\left[\left(U_{x;n_{ex},j}^{n+1}\right)^{}\left(U_{x;n_{sx}1,j}^n\right)^{}\right]\psi _{n_{sx},j}^n,`$
for $`j=1,\mathrm{},n_y`$. These conditions couple the correction $`\varphi `$ to the update of $`A_x`$. To eliminate this coupling, we solve Eq. (46) subject to the reduced interface conditions
$`\varphi _{n_{sx}1,j}^{n+1}`$ $`=`$ $`U_{x;n_{sx}1,j}^{n+1}\varphi _{n_{sx},j}^{n+1},j=1,\mathrm{},n_y,`$ (51)
$`\varphi _{n_{ex}+1,j}^{n+1}`$ $`=`$ $`\left(U_{x;n_{ex},j}^{n+1}\right)^{}\varphi _{n_{ex},j}^{n+1},j=1,\mathrm{},n_y.`$ (52)
When Eq. (46) is replaced by Eq. (48), these conditions are inherited by the system (49).
### 3.4 Fully Implicit Integration
Algorithm IV uses a fully implicit integration procedure for the order parameter,
$`{\displaystyle \frac{\psi _{i,j}^{n+1}\psi _{i,j}^n}{\mathrm{\Delta }t}}`$ $`=`$ $`\left(L_{xx}(U_{x;,j}^n)\psi _{,j}^{n+1}\right)_i+\left(L_{yy}(U_{y;i,}^n)\psi _{i,}^{n+1}\right)_j+N\left(\psi _{i,j}^{n+1}\right),(i,j)\text{Sc},`$ (53)
$`\sigma {\displaystyle \frac{A_{x;i,j}^{n+1}A_{x;i,j}^n}{\mathrm{\Delta }t}}`$ $`=`$ $`\left(D_{yy}A_{x;i,}^{n+1}\right)_j\left(D_{yx}A_{y;,}^n\right)_{i,j}+J_{x;i,j}^n,(i,j)\text{Sc}\text{Bl},`$ (54)
$`\sigma {\displaystyle \frac{A_{y;i,j}^{n+1}A_{y;i,j}^n}{\mathrm{\Delta }t}}`$ $`=`$ $`\left(D_{xx}A_{y;,j}^{n+1}\right)_i\left(D_{xy}A_{x;,}^n\right)_{i,j}+J_{y;i,j}^n.(i,j)\text{Sc}\text{Bl}.`$ (55)
The new element here is the term $`N\left(\psi _{i,j}^{n+1}\right)`$ in the first equation.
The second and third equations are solved again as in the semi-implicit algorithm. The first equation is solved by a slight modification of the method used in the implicit algorithm of the preceding section, The modification is brought about by the approximation
$$N\left(\psi ^{n+1}\right)=\tau \psi ^{n+1}|\psi ^{n+1}|^2\psi ^{n+1}\frac{1}{\mathrm{\Delta }t}\left(S\left(\psi ^n\right)\psi ^n\right),$$
(56)
where $`S`$ is a nonlinear map,
$$S(\psi )=\frac{\tau ^{1/2}\psi }{\left[|\psi |^2+(\tau |\psi |^2)\mathrm{exp}(2\tau \mathrm{\Delta }t)\right]^{1/2}}.$$
(57)
(This approximation is explained in the remark below.) Equation (53) is again of the form (46), but with a different right-hand side,
$$(I\mathrm{\Delta }t(L_{xx}+L_{yy}))\varphi ^{n+1}=G(\psi ^n,A_x^n,A_y^n).$$
(58)
The difference is that, where $`F`$ in Eq. (46) contains a term $`(\mathrm{\Delta }t)N\left(\psi ^n\right)`$, $`G`$ in Eq. (58) contains the more complicated term $`S\left(\psi ^n\right)\psi ^n`$.
#### Remark
The approximation (56) is suggested by semigroup theory. Symbolically,
$$N(\psi )=\underset{\mathrm{\Delta }t0}{lim}\frac{S(\mathrm{\Delta }t)\psi \psi }{\mathrm{\Delta }t}.$$
(59)
To find an expression for the “semigroup” $`S`$, we start from the continuous TDGL equations (6)–(8) (zero-electric potential gauge, $`\mathrm{\Phi }=0`$), using the polar representation $`\psi =|\psi |\mathrm{e}^{i\varphi }`$,
$`_t|\psi |`$ $`=`$ $`\mathrm{\Delta }|\psi ||\psi ||\varphi \kappa ^1𝐀|^2+\tau |\psi ||\psi |^3,`$ (60)
$`|\psi |_t\varphi `$ $`=`$ $`2(|\psi |)(\varphi \kappa ^1𝐀)+|\psi |(\varphi \kappa ^1𝐀),`$ (61)
$`\sigma _t𝐀`$ $`=`$ $`\times \times 𝐀+\kappa ^1|\psi |^2(\varphi \kappa ^1𝐀).`$ (62)
At this point, we are interested in the effect of the nonlinear term $`|\psi |^3`$ on the dynamics. To highlight this effect, we concentrate on the time evolution of the scalar $`u=|\psi |`$ and the vector $`v=\varphi \kappa ^1𝐀`$. (In physical terms, $`u^2`$ is the density of superconducting charge carriers, while $`u^2v`$ is $`\kappa `$ times the supercurrent density.) Ignoring their spatial variations, we have a dynamical system,
$`u^{}`$ $`=`$ $`u|v|^2+\tau uu^3,`$ (63)
$`v^{}`$ $`=`$ $`\epsilon u^2v,`$ (64)
where denotes differentiation with respect to $`t`$, and $`\epsilon =(\kappa ^2\sigma )^1`$. This system yields a pair of ordinary differential equations for the scalars $`x=u^2`$ and $`y=|v|^2`$,
$`x^{}`$ $`=`$ $`2x(\tau xy),`$ (65)
$`y^{}`$ $`=`$ $`2\epsilon xy.`$ (66)
If $`\kappa `$ is large, $`\epsilon `$ is small, and the dynamics are readily analyzed. To leading order, $`y`$ is constant; $`y=0`$ is the only meaningful choice. (Recall that $`xy^{1/2}`$ is $`\kappa `$ times the magnitude of the supercurrent density.) Then the dynamics of $`x`$ are given by
$$x^{}=2x(\tau x).$$
(67)
We integrate this equation from $`t=t_n`$ to $`t`$,
$$x(t)=\frac{\tau x(t_n)}{x(t_n)+(\tau x(t_n))\mathrm{exp}(2\tau (tt_n))}.$$
(68)
In particular,
$$x(t_{n+1})=\frac{\tau x(t_n)}{x(t_n)+(\tau x(t_n))\mathrm{exp}(2\tau \mathrm{\Delta }t)},$$
(69)
where $`\mathrm{\Delta }t=t_{n+1}t_n`$. Since $`x(t_n)=|\psi ^n|^{1/2}`$ and $`x(t_{n+1})=|\psi ^{n+1}|^{1/2}`$, it follows that
$$|\psi ^{n+1}|=\frac{\tau ^{1/2}|\psi ^n|}{[|\psi ^n|^2+(\tau |\psi ^n|^2)\mathrm{exp}(2\tau \mathrm{\Delta }t)]^{1/2}}.$$
(70)
The phase $`\varphi `$ of $`\psi `$ is constant in time. If we multiply both sides by $`\mathrm{e}^{i\varphi }`$, we obtain the expression (57) for the “semigroup” $`S`$.
## 4 Evaluation
We now present the results of several experiments, where the algorithms described in the preceding section were applied to a benchmark problem.
### 4.1 Benchmark Problem
The benchmark problem adopted for this investigation is the equilibration of a vortex configuration in a homogeneous superconductor without defects ($`\kappa =16`$, $`\sigma =1`$, $`\tau =1`$) embedded in a thin insulator (air), where the entire system is periodic in the direction of the free surfaces ($`y`$).
The superconductor measures $`128\xi `$ in the transverse ($`x`$) direction. The thickness of the insulating layer on either side is taken to be $`2\xi `$, so the width of the entire system is $`132\xi `$. The period in the $`y`$ direction is taken to be $`192\xi `$, so the entire system measures $`132\xi \times 192\xi `$.
The computational grid is uniform, with a mesh width $`h_x=h_y=\frac{1}{2}\xi `$. The periodic boundary conditions in the $`y`$ direction are handled through ghost points, so the computational grid has $`264\times 386`$ vertices. The index sets for the superconductor and blanket (see Eqs. (16) and (17)) are
Sc $`=`$ $`\{(i,j):i=5,\mathrm{},260,j=1,\mathrm{},386\},`$ (71)
Bl $`=`$ $`\{(i,j):i=1,\mathrm{},4,261,\mathrm{},264,j=1,\mathrm{},386\}.`$ (72)
The applied field is uniform,
$$H_L=H_R=H=0.5.$$
(73)
(Units of $`H`$ are $`H_c\sqrt{2}`$, so $`H0.707\mathrm{}H_c`$). As there is no transport current in the system, the solution of the TDGL equations tends to an equilibrium state.
### 4.2 Benchmark Solution
First, preliminary runs were made to determine, for each algorithm, the optimal number of processors in a multiprocessing environment. Figure 2 shows the wall clock time for 50 time steps against the number of processors on the IBM SP2.
Each algorithm shows a saturation around 16 processors, beyond which any improvement becomes marginal. All problems were subsequently run on 16 processors.
Next, we used the fully explicit Algorithm I to establish a benchmark equilibrium configuration. We integrated Eqs. (35)–(37) with a time step $`\mathrm{\Delta }t=0.0025`$ (units of $`\xi ^2/D`$), the maximal value for which the algorithm remained stable, and followed the evolution of the vortex configuration by monitoring the number of vortices and their positions. Equilibrium was reached after 10,000,000 time steps, when the number of vortices remained constant and the vortex positions varied less than $`1.0\times 10^6`$ (units of $`\xi `$). The equilibrium vortex configuration had 116 vortices arranged in a hexagonal pattern; see Fig. 3.
The wall clock time for the entire computation was approximately 3,000 minutes. The elapsed time per time step (0.018 seconds) is a measure for the computational cost of Algorithm I.
### 4.3 Evaluation of Algorithms II–IV
With the benchmark solution in place, we evaluated each of the remaining algorithms (II–IV) for stability, accuracy, and computational cost.
We found the stability limit in the obvious way, gradually increasing the time step and integrating to equilibrium until arithmetic divergences caused the algorithm to fail. Equilibrium was defined by the same criteria as for the benchmark solution: no change in the number of vortices and a variation in the vortex positions of less than $`1.0\times 10^6`$.
Because each algorithm defines its own path through phase space, one cannot expect to find identical equilibrium configurations nor equilibrium configurations that are exactly the same as the benchmark. The equilibrium vortex configurations for the four algorithms were indeed different, albeit slightly. To measure the differences quantitatively, we computed the following three parameters: (i) the number of vortices in the superconducting region, (ii) the mean bond length joining neighboring pairs of vortices, and (iii) the mean bond angle subtended by neighboring bonds throughout the vortex lattice. In all cases, the number of vortices was the same (116); the mean bond length varied less than $`1.0\times 10^3\xi `$, and the mean bond angle varied by less than $`1.0\times 10^3`$ radians. Within these tolerances, the equilibrium vortex configurations were the same.
The results are given in Table 1;
$`\mathrm{\Delta }t`$ is the time step at the stability limit (units of $`\xi ^2/D`$), $`N`$ the number of time steps needed to reach equilibrium, and $`C`$ the cost of the algorithm (seconds per time step). From these data we obtain the wall clock time needed to compute the equilibrium configuration, $`T=NC/60`$ (minutes).
Note that the existence of a stability limit for Algorithm IV is a consequence of the implementation in a multiprocessing environment. Since we restrict interprocessor communication to the end of each time step, the fully implicit character of the algorithm is lost. On a single processor, Algorithm IV is fully implicit, and the stability limit is infinite.
## 5 Further Evaluation of Algorithm IV
We evaluated the fully implicit Algorithm IV in more detail by considering its speedup in a multiprocessing environment and its performance under a multi-timestepping procedure.
### 5.1 Parallelism
First, we investigated the speedup of Algorithm IV in a multiprocessing environment, using the the benchmark problem and two other problems, obtained form the benchmark problem by twice doubling the size of the system in each direction. The mesh width was kept constant at $`\frac{1}{2}\xi `$, so the resulting computational grid had $`264\times 386`$ vertices for the benchmark problem, $`528\times 772`$ vertices for the intermediate problem, and $`1056\times 1544`$ vertices for the largest problem. Speedup was defined as the ratio of the wall clock time (exclusive of I/O) to reach equilibrium on $`p`$ processors divided by the time to reach equilibrium on a single processor for the benchmark and intermediate problem, or twice the time to reach equilibrium on two processors for the largest problem. (The largest problem did not fit on a single processor.) The results are given in Fig. 4.
The curve for the benchmark problem was obtained as an average over many runs; the data for the intermediate and largest problem were obtained from single runs, hence they are less smooth. The speedup is clearly linear when the number of processors is small; it becomes sublinear at about 12 processors for the smallest problem, 14 processors for the intermediate problem, and 18 processors for the largest problem.
### 5.2 Multi-timestepping
The final set of experiments shows that the performance of Algorithm IV is enhanced by a multi-timestepping procedure, where $`𝐀`$ is updated less frequently than $`\psi `$,
$`{\displaystyle \frac{\psi _{i,j}^{n+1}\psi _{i,j}^n}{\mathrm{\Delta }t}}`$ $`=`$ $`\left(L_{xx}(U_{x;,j}^n)\psi _{,j}^{n+1}\right)_i+\left(L_{yy}(U_{y;i,}^n)\psi _{i,}^{n+1}\right)_j+N\left(\psi _{i,j}^{n+1}\right),(i,j)\text{Sc},`$ (74)
$`\sigma {\displaystyle \frac{A_{x;i,j}^{n+m}A_{x;i,j}^n}{m\mathrm{\Delta }t}}`$ $`=`$ $`\left(D_{yy}A_{x;i,}^{n+m}\right)_j\left(D_{yx}A_{y;,}^n\right)_{i,j}+J_{x;i,j}^n,(i,j)\text{Sc}\text{Bl},`$ (75)
$`\sigma {\displaystyle \frac{A_{y;i,j}^{n+m}A_{y;i,j}^n}{m\mathrm{\Delta }t}}`$ $`=`$ $`\left(D_{xx}A_{y;,j}^{n+m}\right)_i\left(D_{xy}A_{x;,}^n\right)_{i,j}+J_{y;i,j}^n.(i,j)\text{Sc}\text{Bl}.`$ (76)
When $`m=1`$, both $`\psi `$ and $`𝐀`$ are updated at every time step, but when $`m`$ is greater than 1, $`𝐀`$ is updated only every $`m`$th time step. In the limit as $`m\mathrm{}`$, this procedure yields the frozen-field approximation, which is a good approximation of the Ginzburg–Landau model near the upper critical field when the charge of the superconducting charge carriers is small .
We applied this modification of Algorithm IV with $`m=10,15`$ to the benchmark problem of §4. The results are given in Table 2.
All computations were done with $`\mathrm{\Delta }t=0.19`$ (units of $`\xi ^2/D`$). The data for the computation with $`m=1`$ are taken from Table 1. We observe that the cost of the algorithm ($`C`$, seconds per time step) decreases with increasing $`m`$, while the number of time steps to equilibrium ($`N`$) increases. If $`m=10`$, the overall wall-clock time ($`T=NC/60`$, minutes) is approximately one-third less than the wall-clock time for $`m=1`$. For larger values of $`m`$, the increase in the number of steps needed to reach equilibrium offsets any gain from the decrease in cost. These data suggest an optimal strategy, where $`𝐀`$ is updated every 10 time steps. Figure 5 shows the effect of the updating frequency on the evolution of the free-energy functional (9). Note the dramatic increase of the wall-clock time for $`m=15`$. Eventually, the curve for $`m=15`$ merges with the other curves, but this happens well beyond the range of the figure.
## 6 Conclusions
The results of the investigation lead to the following conclusions.
(i) One can increase the time step $`\mathrm{\Delta }t`$ nearly 80-fold, without losing stability, by going from the fully explicit Algorithm I to the fully implicit Algorithm IV.
(ii) As one goes to the fully implicit Algorithm IV, the complexity of the matrix calculations and, hence, the cost $`C`$ of a single time step increase.
(iii) The increase in the cost $`C`$ per time step is more than offset by the increase in the size of the time step $`\mathrm{\Delta }t`$. In fact, the wall clock time needed to compute the same equilibrium state with the fully implicit Algorithm IV is one-sixth of the wall clock time for the fully explicit Algorithm I.
(iv) The (physical) time to reach equilibrium—that is, $`N\mathrm{\Delta }t`$, the number of time steps needed to reach equilibrium times the step size—is (approximately) the same for all algorithms, namely, 25,000 (units of $`\xi ^2/D`$).
(v) The fully implicit Algorithm IV displays linear speedup in a multiprocessing environment. The speedup curves show sublinear behavior when the number of processors is large.
(vi) The performance of the fully implicit Algorithm IV can be improved further by a multi-timestepping procedure, where the vector potential $`𝐀`$ is updated less frequently than the order parameter $`\psi `$.
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# 1 Introduction
## 1 Introduction
The massive binary systems are often considered as sites of high energy processes in which $`\gamma `$-ray production is expected. $`\gamma `$-rays in these systems may be produced in interactions of relativistic particles injected by: a neutron star or a black hole (e.g. Bignami, Maraschi & Treves 1977, Kirk, Ball & Skjaeraasen 1999); a shock wave created in collision of the pulsar and stellar winds (e.g. Harding & Gaisser 1990) or two stellar winds (e.g. Eichler & Usov 1993). Observations of these systems in GeV and TeV energies suggest that in fact it may be the case. For example in last years the Compton GRO detectors reported point sources with flat spectra coincident with some massive binaries e.g. Cyg X-3, LSI 303<sup>o</sup>+61, or Cen X-3. Some positive detections of massive binaries by Cherenkov telescopes has been also claimed, although they were frequently accompanied by many negative reports (see for a review Weekes 1992, Moskalenko 1995). In the recent review on TeV observations only Cen X-3 has been mentioned as possible massive binary active at these energies (Weekes 1999).
These theoretical predictions and positive observations has stimulated analysis of propagation of VHE $`\gamma `$-rays in the anisotropic radiation fields of accretion disks around compact objects in massive binaries (Carraminana 1992, Bednarek 1993), and in the radiation fields of massive stars (e.g. Protheroe & Stanev 1987, Moskalenko, Karakuła & Tkaczyk 1993, Moskalenko & Karakuła 1994). Recently we have performed Monte Carlo simulations of cascades initiated by monoenergetic $`\gamma `$-ray beams injected by a discrete source (probably a compact object) in the radiation field of a massive companion. Two massive binaries, Cyg X-3 (Lamb et al. 1977, Merck et al. 1995) and LSI 303<sup>o</sup>+61 (van Dijk et al. 1993, Hermsen et al. 1977, Thompson et al. 1995), have been discussed in the context of this problem (Bednarek 1997).
Another massive X-ray binary, Cen X-3, containing a neutron star with the period 4.8 s on a 2.09 day orbit around an O-type supergiant, has been detected above $`100`$ MeV by the EGRET detector on the Compton Observatory (Vestrand, Sreekumar & Mori 1997). There are evidences that $`\gamma `$-ray emission at these energies has a form of outbursts and is modulated with a 4.8 s period of the pulsar. Based on the observations in late 80’s, two Cherenkov groups (Raubenheimer et al. 1989, Brazier et al. 1990, North et al. 1990) reported detection of positive signal at TeV energies from Cen X-3 at an orbital phase of $`0.75`$, which is modulated with a period of the pulsar. This emission has been localized by North et al. (1991) and Raubenheimer & Smith (1997) to a relatively small region between the pulsar orbit and the surface of a massive companion which may be the accretion wake or the limb of the star. More recently the Durham group has detected a persistent flux of $`\gamma `$-rays above 400 GeV on a lower level than previous reports (Chadwick et al. 1998,1999). No evidence of correlation with the pulsar or orbital periods has been found and no evidence of correlation with the X-ray flux has been detected (Chadwick et al. 1999a,b).
The purpose of this paper is to compute the $`\gamma `$-ray light curves and $`\gamma `$-ray spectra escaping from Cen X-3 system, assuming different geometries and energy distributions for primary electrons or photons injected into the radiation field of the massive star by a compact object. The $`\gamma `$-ray spectrum escaping from the system is a consequence of anisotropic cascades which initiate primary electrons or $`\gamma `$-ray photons in the thermal radiation of a massive star in Cen X-3. The results of computations are discussed in the context of observations of Cen X-3 in GeV and TeV energy ranges.
We assumed the following parameters for the massive binary Cen X-3: the radius of the O star $`r_\mathrm{s}=8.6\times 10^{11}`$ cm, its surface temperature $`T_\mathrm{s}=3\times 10^4`$ K (Krzemiński 1974), the binary star separation $`a=1.4r_\mathrm{s}=1.2\times 10^{12}`$ cm, the inclination angle $`70^o`$ and the circular orbit since the known eccentricity of the binary is $`<0.0008`$ (Fabbiano & Schreier 1977).
## 2 Propagation of VHE $`\gamma `$-rays in the radiation field of a massive star
In the previous paper (Bednarek 1997), we considered the propagation of monoenergetic $`\gamma `$-ray beams in the radiation filed of a massive star in the case of two massive binaries: Cyg X-3 and LSI 303<sup>o</sup>+61. We determined general conditions for which VHE $`\gamma `$-rays may initiate such kind of inverse Compton cascade (ICS) (see Sect. 2 in Bednarek 1997). The massive star in Cen X-3 system has also high enough luminosity that VHE photons injected inside a few stellar radii should initiate the cascade. This expectation has been checked by computing the optical depth for $`\gamma `$-ray photons as a function of distance from the surface of the massive star, the photon energies, and the angle of injection $`\alpha `$ measured with respect to the direction determined by the place of photon injection and the center of the massive star. The optical depth has been computed as shown in Appendix A of Bednarek (1997). We have found that for the parameters of the massive star in Cen X-3 system, the optical depth for photons injected at angles $`\alpha >30^o`$ and at the distance $`x=1.4r_\mathrm{s}`$, can be greater than one. If photons are injected closer to the massive star, then the optical depth can be greater than one for all injection angles (see curve for $`x=1.1r_\mathrm{s}`$ in Fig. 1a). The optical depth reaches maximum (characteristic cusps visible in Fig. 1a), for photons propagating in directions tangent to the massive star surface, i.e. at an angle $`\alpha _\delta =\pi \delta `$, where $`\delta `$ is the angle intercepted by the massive star. For higher angles $`\alpha `$, the optical depth is computed up to the moment of photon collision with the stellar surface and therefore it is lower than for $`\alpha _\delta `$. As expected the optical depth for $`\gamma `$-ray photon reaches maximum at photon energies determined by characteristic energies of soft star photons (see Fig. 1b). Photons with energies below a few tens of GeV escape without significant absorption. Therefore it is expected that cascade $`\gamma `$-ray spectra should show a break close to these energies.
## 3 Gamma-ray spectra from ICS cascade
Let us assume that the $`\gamma `$-ray photon is injected by the compact object which is on an orbit around the star with the parameters characteristic for Cen X-3 system. As we showed above these photons may create $`e^\pm `$ pair in collision with the soft star photon. The secondary pairs can next produce ICS $`\gamma `$-rays, initiating in this way the ICS cascade in the massive star radiation which is seen anisotropic in respect to the location of the injection place of the primary photons or electrons. In the subsections we describe the cascade scenario and discuss the results of calculations for different initial injection conditions, i.e. type, distribution, and spectra of primary particles.
### 3.1 ICS cascade with isotropisation of secondary electrons
The $`\gamma `$-ray photon with energy $`E_\gamma `$ interacts with the star photon creating $`e^\pm `$ pair at the propagation distance $`l_\gamma `$. The place of the creation can be simulated randomly from
$`P_1`$ $`=`$ $`\mathrm{exp}\left({\displaystyle _0^{\mathrm{l}_\gamma }}\lambda _{\gamma \gamma }^1(E_\gamma ,x_\gamma ,\alpha _\gamma )𝑑l\right)`$ (1)
where $`l`$ is the distance measured along the photon propagation, $`\lambda _{\gamma \gamma }(E_\gamma ,x_\gamma ,\alpha _\gamma )`$ is the mean free path for $`\gamma `$-ray photon with energy $`E_\gamma `$, injected at the angle $`\alpha _\gamma `$ and at the distance $`x_\gamma `$ from the center of the massive star in Cen X-3 system, and $`P_1`$ is the random number. If condition (1) can not be fulfilled for chosen random number $`P_1`$, we accept that the $`\gamma `$-ray: escapes from the star radiation field, or collides with the star surface.
Simple trigonometry allows us to determine the distance of produced $`e^\pm `$ pair, $`x_{e^\pm }`$, from the center of the massive star. The energy of created electron (positron) $`E_\mathrm{e}`$ is chosen by sampling from the differential spectrum of pairs which can be produced by the $`\gamma `$-ray photon at the propagation distance $`l_\gamma `$. It is obtained from
$`P_2`$ $`=`$ $`\left({\displaystyle _{0.5}^{E_e}}{\displaystyle \frac{dW}{dEdx}}𝑑E\right)\left({\displaystyle _{0.5}^{E_{e,max}}}{\displaystyle \frac{dW}{dEdx}}𝑑E\right)^1`$ (2)
where $`dW/dEdx`$ is defined in Bednarek (1997, Appendix B), $`E_{e,max}`$ is the maximal possible energy of the electron produced in $`\gamma `$-ray – soft photon collision, and $`P_2`$ is the random number. The energy of positron is then $`E_p=E_\gamma E_e`$.
We assume that secondary $`e^\pm `$ pairs are locally isotropised by the random component of the magnetic field. This will happen if the random component of the magnetic field $`B_\mathrm{r}`$ inside the binary system is strong enough (Bednarek 1997), i.e.
$`B_r4\times 10^{29}\gamma _\mathrm{e}^2T^4D,`$ (3)
where $`\gamma _\mathrm{e}`$ is the Lorentz factor of electron (or positron), $`T`$ is the surface temperature of the massive star, and $`D`$ is the dilution factor of the star radiation defined as the part of the sphere intersepted by the star. $`D0.1`$ for the separation of the companions $`x=1.4r_\mathrm{s}`$. The $`e^\pm `$ pairs cool locally and produce next generation of $`\gamma `$-rays in the inverse Compton scattering of soft photons coming from the massive star. In order to select the direction of the secondary $`\gamma `$-ray we compute at first the energy loss rate on ICS for pairs with energy $`E_{e^\pm }`$ assuming that the electron (positron) is located at the distance $`x_{e^\pm }`$ and propagates at an angle $`\alpha _{e^\pm }`$ to the direction defined by $`x_{e^\pm }`$ (see Appendix C in Bednarek 1997). The direction of motion of a secondary $`\gamma `$-ray, which in fact covers with the direction of motion of relativistic electron at the moment of $`\gamma `$-ray production, is obtained by sampling from the distribution of the energy loss rates computed as a function of the angle $`\alpha _{e^\pm }`$ after its normalization to the total energy loss rate for electrons moving isotropically at $`x_{e^\pm }`$. The cosine of the angle of $`\gamma `$-ray emission ($`\mathrm{cos}\alpha _\gamma `$) is obtained from
$`P_3`$ $`=`$ $`\left({\displaystyle _{\mathrm{cos}\alpha _\gamma }^1}{\displaystyle \frac{dL}{dtd\nu }}𝑑\nu \right)\left({\displaystyle _1^1}{\displaystyle \frac{dL}{dtd\nu }}𝑑\nu \right)^1,`$ (4)
where $`P_3`$ is the random number, and
$`{\displaystyle \frac{dL}{dtd\nu }}`$ $`=`$ $`{\displaystyle _0^{E_{\gamma ,max}}}{\displaystyle \frac{dN}{dtdE_\gamma }}{\displaystyle \frac{dN_{e^\pm }}{d\nu }}E_\gamma 𝑑E_\gamma ,`$ (5)
$`dN/dtdE_\gamma `$ is the photon spectrum produced by electrons in ICS process (see Appendix C in Bednarek 1997), $`dN_{e^\pm }/d\nu `$ describes the number of electrons moving inside the unit cosine angle $`d\nu =d(\mathrm{cos}\alpha )`$, and $`E_{\gamma ,max}`$ is the maximum energy of the $`\gamma `$-ray photon produced in the ICS process by an electron with energy $`E_{e^\pm }`$, located at the distance $`x_{e^\pm }`$, and propagating at the angle $`\alpha _{e^\pm }=\alpha _\gamma `$. In this way the directions of motion of secondary ICS $`\gamma `$-rays are determined. Now we need to determine their energies. The mean energy of secondary $`\gamma `$-ray photon, $`E_\gamma ^{}`$, produced by secondary electron with energy $`E_{e^\pm }`$ at the distance $`x_{e^\pm }`$ from the star and at an angle $`\alpha _\gamma `$, is obtained from the formula
$`<E_\gamma >`$ $`=\left({\displaystyle \frac{dL}{dtd\nu }}\right)\left({\displaystyle \frac{dN}{dtd\nu }}\right)^1,`$ (6)
where
$`{\displaystyle \frac{dN}{dtd\nu }}`$ $`=`$ $`{\displaystyle _0^{E_{\gamma ,max}}}{\displaystyle \frac{dN}{dtdE_\gamma }}{\displaystyle \frac{dN_{e^\pm }}{d\nu }}𝑑E_\gamma `$ (7)
is the rate of photon production by an electron in the ICS process. Above described procedure is repeated for all cascade $`e^\pm `$ pairs up to the moment of their ’complete cooling’, i.e. to the moment at which they are not able to produce $`\gamma `$-ray photons in ICS process with energies above certain applied value which in our simulations is chosen as equal to 100 MeV. We follow all secondary ICS $`\gamma `$-rays with energies above the threshold for $`e^\pm `$ pair production in radiation of a massive star.
### 3.2 Injection of monodirectional and monoenergetic photons or electrons
It is likely that high energy photons and/or electrons are injected by the neutron star (a compact object in Cen X-3 system) in the form of highly collimated beams. Such beams may be formed: in the outer gaps of pulsar magnetospheres (e.g. Cheng, Ho & Rudermann 1986); in collisions of collimated protons escaping from the pulsar magnetospheres with the matter of an accretion disk (Cheng & Ruderman 1989); or they may emerge from regions of the magnetic poles of the neutron star (Kiraly & Meszaros 1988). In our simulations it is assumed that highly collimated beam of $`\gamma `$-rays with energy 10 TeV emerge from the distance $`x=1.4r_\mathrm{s}`$ at a certain angle $`\alpha `$, measured with respect to the direction determined by the center of the massive star and the compact object. In fact such situation corresponds also to the case of the highly collimated electron beams with comparable energy since at these energies electrons transfer almost all its primary energy to a single $`\gamma `$-ray photon in ICS process because the scattering occurs in the Klein-Nishina domain.
We are interested in the spectra of secondary photons which escape from the above described type of cascade at a range of angles $`\mathrm{\Delta }\mathrm{cos}\theta =0.2`$. In Fig. 2, we show such secondary spectra for different injection angles of primary monoenergetic photons with energy $`E_\gamma =10^7`$ MeV: $`\mathrm{cos}\alpha =1`$ (a), 0.5 (b), 0. (c), $`0.5`$ (d), and $`1`$ (e), observed at the range of angles: $`\mathrm{cos}\theta =0.8÷1.`$ (dotted histograms), $`0.4÷0.6`$ (dashed), $`0.2÷0`$ (dot-dashed), $`0.6÷0.4`$ (dot-dot-dashed), and $`1÷0.9`$ (long-dashed). The cascading effects are the strongest (the highest numbers of secondary $`\gamma `$-rays) if the monoenergetic $`\gamma `$-ray beam is injected at the intermediate angles $`\alpha `$ (see thick solid histograms in Fig. 2, which show the secondary photon spectra intergated over all solid angle). This is a consequence of the highest optical depth for photons propagating at directions tangent to the massive star limb (curve for $`x=1.4r_\mathrm{s}`$ in Fig. 1a). All secondary spectra are well described by a power law with the index $`1.5`$ at energies below $`10`$ GeV (as expected in ICS cascades with complete cooling of electrons). The photon intensities in these spectra are the highest at direction tangent to the limb of the massive star (dot-dot-dot-dashed histograms in Fig. 2). At higher photon energies ($`>10`$ GeV), the spectra show characteristic cut-offs due to the absorption of secondary photons in the massive star radiation. The spectra recover at TeV energies with the intensities depending on the observation angle. The highest intensities are observed at small angles $`\alpha `$ for which the optical depth for TeV $`\gamma `$-rays is the lowest (see Fig. 1). In Figure 2f, we compare the spectra escaping at fixed range of angles $`\mathrm{cos}\theta =00.2`$ but for different angles of injection of primary photons $`\mathrm{cos}\alpha =1.`$ (dotted histogram), 0.5 (dashed), 0. (dot-dashed), $`0.5`$ (dot-dot-dot-dashed), and $`1.`$ (long dashed). It is clear that primary photons injected at directions of the highest optical depth ($`\mathrm{cos}\alpha =0.5÷0.`$) produce the secondary photons with the highest intensities than thes onces injected at directions of the lowest optical depth ($`\mathrm{cos}\alpha =1.`$).
Note that in the type of cascade discussed here, the intensity of secondary photons observed at the direction of injection of primary photons is not completely dominated by these primary photons because of the large optical depths for the considered injection place ($`x=1.4r_\mathrm{s}`$), and the isotropisation of secondary cascade $`e^\pm `$ pairs. Our assumptions on the propagation of photons in massive binaries are different than these ones applied in e.g. Kirk, Ball & Skjaeraasen (1999).
### 3.3 Isotropic injection of monoenergetic photons or electrons
The monoenergetic photons and electrons can be isotropically injected into the binary system by the young pulsars with strong pulsar winds. In fact, it is believed the pulsar winds are composed with relativistic electrons (positrons) which have typical Lorentz factors of the order $`10^7`$ (Rees & Gunn 1974). These electrons, if accelerated by the electric fields generated during magnetic reconnection in the pulsar wind zone, can be injected approximately isotropically. Therefore we consider the case of isotropic injection of photons or electrons with energies 10 TeV. If only 10 TeV electrons are injected, they should produce photons with comparable energies in a single scattering as mentioned above. The secondary photon spectra, produced in cascade under such initial conditions, are shown in Fig. 3 for different range of observation angles defined by: $`\mathrm{cos}\theta =0.8÷1.`$ (dotted), $`0.4÷0.6`$ (dashed), $`0÷0.2`$ (dot-dashed), $`0.6÷0.4`$ (dot-dot-dot-dashed), $`1÷0.8`$ (long-dashed). The photon intensities observed at directions defined by very small angles $`\theta `$ and directions tangent to the massive star limb behaves differently in different energy ranges. The intensities of GeV photons are the highest for directions tangent to the stellar limb and the lowest for small angles. This is in contrary what is observed at TeV energies. Note, that significant amount of secondary photons emerges also at the range of angles $`\mathrm{cos}\theta =1÷0.8`$. This is on the opposite side of the massive star than the location of the compact source of primary photons and/or electrons (directions obscured by the star!). Therefore, the soft radiation of a luminous star may work as a kind of lens for high energy $`\gamma `$-ray photons causing the effects of focusing of very high energy photons.
### 3.4 Isotropic injection of electrons with the power law spectrum
We consider also the case of isotropic injection of electrons with the power law spectrum. Such electrons can be accelerated at the shock front created in collision of the pulsar wind with the surrounding matter as proposed in the model by Kennel & Coroniti (1984). It is assumed that electrons have the power law spectrum with index $`2`$ (as expected from the theory of acceleration in strong shocks). They are injected from the discrete source orbiting the massive star in Cen X-3. In the binary system the shock may form relatively close to the compact object because the pulsar in Cen X-3 is relatively slow, and the density of surrounding plasma inside the binary system is high.
The results of calculations of the secondary photon spectra escaping from ICS cascade, after integration over all solid angle, are shown in Fig. 4a for different distances of the discrete source of primary electrons from the massive star. The results show that the integrated spectra in the energy region below $`10`$ GeV do not depend significantly on the location of the injection distance (in the range from the surface up to $`5r_\mathrm{s}`$). This effect may result from our assumption on the local capturing of secondary $`e^\pm `$ pairs by the random magnetic field. If the dominant magnetic field has ordered structure inside the binary then the propagation effects of pairs may be important. Such more complicated cascade scenario is out of the scope of this paper and will be discussed in the future work. The photon intensity decreases drastically at TeV energies for injection distance closer to the massive star surface.
We investigate also the dependence of the shape of the escaping spectrum of secondary $`\gamma `$-rays on the observation angle $`\theta `$, for distance of injection $`x=1.4r_\mathrm{s}`$ (see Fig. 4b). The features of these angular dependent spectra are similar to these ones described above for monoenergetic injection of electrons. The highest intensities at TeV energies are observed at small angles $`\theta `$ and the lowest intensities at directions behind the massive star. The highest intensities at GeV energies are for directions tangent to the massive star limb and the lowest for small angles and at directions obscured by the massive star.
### 3.5 Isotropic injection of photons with the power law spectrum
Finally we discuss the case of isotropic injection of photons with the power law spectrum and spectral index $`2`$. We show in Fig. 5 the spectra of escaping $`\gamma `$-rays for: different distances of the injection place from the massive star (Fig. 5a); separately, the escaping spectra of primary $`\gamma `$-rays and secondary $`\gamma `$-rays for the injection place at $`x=1.4r_\mathrm{s}`$ (Fig. 5b); and the angular dependence of secondary $`\gamma `$-ray spectra on the observation angle $`\theta `$ (Fig. 5c). General features of these spectra are very similar to the escaping spectra produced by electrons with the power law spectrum. The significant differences appear at low energies (below a few GeV) and at high energies (above $`1`$ TeV), and results from the contribution to the escaping spectrum from the primary $`\gamma `$-rays which do not cascade in the radiation of the massive star in Cen X-3. The escaping primary $`\gamma `$-rays flattens the spectrum at energies below a few $`10^4`$ Gev (spectral index close to 1.9), and dominates the angle integrated spectrum at TeV energies.
## 4 Consequences for gamma-ray escape from Cen X-3
In the previous section we investigated general features of the $`\gamma `$-ray spectra escaping from the radiation field of a massive star. In order to have results of calculations which can be directly compared with the observations we compute the $`\gamma `$-ray light curves expected from Cen X-3 system assuming that the compact object in this binary system injects $`\gamma `$-ray photons or electrons with the power law spectrum. The parameters of the Cen X-3 system, used in computations, are mentioned in the Introduction. Note that the orbit of the compact object is almost circular, so then the expected light curve should be symmetrical. Therefore we compute only the photon fluxies for the phases from 0 to 0.5, where the phase is measured from the side of the observer.
The $`\gamma `$-ray light curves of photons escaping from the system in the case of isotropic injection of electrons (or positrons) with the power law spectrum and spectral index $`2`$ are shown in Figs. 6a,b for photons with energies: above 100 MeV, i.e. the EGRET energy range (a), and above 300 GeV, i.e. Cherenkov technique energy range (b). The results are shown for the cut-offs in the spectrum of electrons at $`10^7`$ MeV (full histogram) and at $`10^8`$ MeV (dashed histogram). Note however, that the case with cut-off at $`10^8`$ MeV we show only for comparison since it may not be completely right. Our assumption on the isotropization of secondary $`e^\pm `$ pairs with the Lorentz factors $`10^7`$ may not be justified in this case(see Eq. 3). The light curves show that the $`\gamma `$-ray flux should change drastically during $`2.09`$ day the orbital period of the system by at least an order of magnitude. However the $`\gamma `$-ray light curves observed at different energy ranges behaves completely different. When the photon flux above 100 MeV increases from the phase equal 0. up to the eclipse of the compact object by the massive star, which occurs for the phase $`0.38`$, the photon flux above 300 GeV decreases. This is the result of propagation of photons in the anisotropic radiation of a massive star as discussed in details in section 3. In Fig. 7, we show the spectra of $`\gamma `$-rays which should be seen by the observer for different phases of the compact object: 0. (full histogram), 0.15 (dotted), 0.35 (dashed), and 0.5 (dot-dashed). The photon spectra above 100 MeV have similar shapes but different intensities. The spectra above 300 GeV differs significantly not only in the intensities but also by the shape.
We have also computed the $`\gamma `$-ray light curves in the case of isotropic injection of primary photons with the power law spectrum and index $`2`$ (see Figs. 8a,b). Specific histograms in these figures show the light curves for all escaping $`\gamma `$-ray photons (full histograms), primary photons which escape without cascading (dotted), and secondary photons produced in cascades (dashed). As expected the light curves for secondary photons in the case of injection of primary photons and electrons are very similar. However the contribution of escaping primary photons to the $`\gamma `$-ray light curves with energies above 100 MeV dominates the secondary photons. Altogether, the $`\gamma `$-ray light curves at energies above 100 MeV are very flat with the strong decrease for phases between $`0.38÷0.62`$ resulting from the eclipse condition. During the eclipse, the observer may only detect secondary photons produced in cascades (dashed histogram in Fig. 8b), but on the level of about an order of magnitude lower. The $`\gamma `$-ray light curves above 300 GeV do not differ significantly for the case of injection of primary photons or electrons (compare Fig 6b and 8b). The contribution of primary non-cascading photons dominates only for small values of the phase (dotted histogram in Fig. 8b). From these computations it becomes clear that investigation of the $`\gamma `$-ray light curves at photon energies above 100 MeV (but not above 300 GeV) should allow to distinguish what kind of primary particles is dominantly produced by the compact object, photons or electrons (positrons), provided that these particles are injected isotropically with the power law spectrum.
We also show in Fig. 9a,b the spectra of escaping photons for different phases of the compact object, separately for secondary cascade photons and for primary photons which escape without interaction. The photon spectral index below $`10`$ GeV for all escaping photons (primary plus secondary) vary with phase only in relatively small range, from $`1.8`$ to $`2`$. The observed photon fluxies are almost constant. At TeV energies the spectra change drastically with the phase of the compact object, similarly to the above discussed case of injection of primary electrons. Note, that for phase 0.5 (corresponding to the total eclipse of the compact object by the massive star) only secondary photons at energies below $`10`$ GeV can be observed (Fig. 9a and b).
As we have discussed in the Introduction, Cen X-3 has been detected in GeV and TeV energy range. The emission in GeV energy range can be fitted by the power law with the spectral index $`1.81\pm 0.37`$ (Vestrand et al. 1997). This index is consistent with our results for both discussed models of isotropic injection of primary photons or electrons with the power law spectrum and spectral index $`2`$. However, the EGRET observations indicate modulation of GeV emission with the pulsar’s spin period, which should not be observed in the case of injection of primary electrons since the escaping photons at these energies were produced in the cascade process and the information on the pulsar period should disappear. Therefore the model with injection of primary electrons by the compact object seems not work. The modulation with the pulsar period might be observed in the case of injection of primary photons from the compact object, provided that the secondary photons do not completely dominate the primary escaping photons. In fact, this is evident from our simulations (see Figs. 9a,b). However, as it is seen in Fig. 8a, the photon flux, although constant though most of the phase range, should drop drastically during the eclipse of the compact object by the massive star. This feature has not been observed but also can not be rejected by the EGRET observations (Vestrand et al. 1997).
Cen X-3 has been also reported as a source of TeV photons modulated with the orbital period of the binary system by earlier, less sensitive Cherenkov observations (Brazier et al. 1990, North et al. 1990). Recent observations report that Cen X-3 is a source of steady emission above $`400`$ GeV (Chadwick et al. 1998). However modulation with the pulsar and orbital periods has not been found (Chadwick et al. 1999b). Our calculations show that in both models, injection of primary photons and injection of primary electrons by the compact object, the modulation of the signal with the orbital period should be very clear. In contrary, the modulation with the pulsar period should not be observed because the secondary cascade photons determines the light curve in the case of injection of primary electrons and dominate or give similar contribution to the light curve in the case of injection of primary photons (see dotted histogram in Fig. 8b).
## 5 Conclusion
We considered the cascade initiated by photons or electrons injected from the compact object in the radiation field of a massive companion in Cen X-3 system, assuming that secondary electrons are isotropised by the magnetic field in the binary. It is found that the features of the escaping photons from such massive binaries (the light curves, photon spectra as a function of the phase of the compact object) may allow to distinguish which particles are injected by the compact object, i.e. photons or electrons. If the cascades are initiated by electrons then the escaping secondary cascade photons should not show features of modulation with the pulsar period. This seems to be in contradiction with the observations of Cen X-3 at energies above 100 MeV. Therefore we reject such model. If primary photons are injected isotropically with the power law spectrum and spectral index $`2`$, then the spectrum of escaping photons at energies below $`10`$ GeV is dominated, for most of the phases, by the primary non-cascading photons. Then, the modulation with the pulsar period can be observed. At TeV energies, the modulation of the photon flux with the 2.09 day binary period should be strong. The observation of modulation of TeV signal with the orbital period of the system have been reported by earlier observations of Cen X-3 system (mensioned above) but not confirmed recently (Chadwick et al. 1999b). We think that this problem needs further investigation since the sensitivity of present observations is still rather poor (see the TeV $`\gamma `$-ray light curve in Fig. 3 in Chadwick et al. 1999b).
If the lack of modulation of the TeV photon flux with the Cen X-3 binary period is real, then a more complicated model has to be investigated. An extended source, e.g. a shock inside the binary system, injecting primary photons or electrons which initiate cascades in the soft radiation of a massive star should be developed. However such computations will require much more computing time in order to get satisfactory statistics, and therefore are left for the future work.
In the present calculations we neglected the X-ray radiation field produced by the compact object. Its energy density $`L_\mathrm{X}10^{38}`$ erg s<sup>-1</sup> (Giacconi et al. 1971) is comparable to the energy density of thermal photons from the massive star, so their photon density is a few orders of magnitude lower inside the volume of the binary system. We neglect also the heating effects of the massive star by the X-rays coming form the compact object since the power emitted by the stellar surface is higher than the X-ray power falling on the massive star from the orbital distance of the compact object equal to $`1.4`$ stellar radii. However note that X-rays, produced in the accretion disk around a compact object, i.e. close to the production site of primary photons, may absorb $`\gamma `$-rays if they are also produced close to the inner disk radius (see e.g. Bednarek 1993).
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# A Possible X-ray and Radio Counterpart of the High-Energy Gamma-ray Source 3EG J2227+6122
## 1 Introduction
The nature of the persistent high-energy ($`>100`$ MeV) $`\gamma `$-ray sources in the Galaxy remains an enigma almost three decades after their discovery. The third installment of the EGRET catalog (Hartman et al. 1999) lists a total of 270 sources, of which 93 are likely or possibly identified with blazars, five with rotation-powered pulsars, one with Cen A, and one with the LMC. This leaves 170 unidentified sources, of which 57, or one third, lie within $`|b|10^{}`$ along the Galactic plane. This excess of low-latitude sources must comprise a Galactic population that is either similar to the already identified $`\gamma `$-ray pulsars, or representative of an entirely new class of $`\gamma `$-ray emitter associated with the disk population. These Galactic sources have proven extremely difficult to identify.
Rotation-powered pulsars are likely to dominate the Galactic $`\gamma `$-ray source population, but their detectability at both radio and $`\gamma `$-ray wavelengths depends on their beam patterns. The shapes of radio pulsar beams as determined by the rotating vector model (Radhakrishnan & Cooke 1969) demands that $`5070\%`$ of young radio pulsars are not visible from Earth (Frail & Moffett 1993; Tauris & Manchester 1998). The clear differences between the broad observed $`\gamma `$-ray beam patterns and the narrow radio pulses implies that $`\gamma `$-ray emission is probably visible from a more complete range of directions than are the radio beams. Indeed, the identification of the high-energy $`\gamma `$-ray source Geminga as the first radio quiet, but otherwise ordinary pulsar (Halpern & Holt 1992; Bertsch et al. 1992), provides a likely prototype for the remaining unidentified Galactic sources. That Geminga was the brightest unidentified EGRET source, at a distance of only $`250`$ pc, argues that there should be others. The possible existence of intrinsically radio-quiet $`\gamma `$-ray pulsars also cannot be ruled out.
Several authors have considered the pulsar hypothesis from statistical or theoretical points of view (Halpern & Ruderman 1993; Helfand 1994; Kaaret & Cottam 1996). The most detailed theoretical treatment of the pulsar model for the Galactic $`\gamma `$-ray sources is that of Romani & Yadigaroglu (1995) and Yadigaroglu & Romani (1995;1997). They developed a numerical calculation of $`\gamma `$-ray production and beaming in the outer-gap model that successfully reproduces the basic observed features of the pulse profiles and the $`\gamma `$-ray efficiency as a function of age. By combining this model with a Monte Carlo simulation of the Galactic pulsar population, they estimated that a total of 22 pulsars should be detected by EGRET at the threshold of the first EGRET catalog. This number is very close to the actual number of unidentified sources at $`|b|<10^{}`$ in that catalog.
Candidate neutron-star identifications for additional EGRET sources have been found in the form of point-like X-ray sources with no obvious optical counterparts. One such candidate is present in the supernova remnant CTA1 which is close to 3EG J0010+7309 (Lamb & Macomb 1997; Brazier et al. 1998). Another one is in the $`\gamma `$-Cygni supernova remnant which is coincident with 3EG J2020+4017 (Brazier et al. 1996). Both of these so far lack detected of pulsations. Most notable among the probable identifications is the radio star and Be/X-ray binary LSI $`+61^{}303`$ (Strickman et al. 1998) which has long been associated with the $`\gamma `$-ray source 2CG 135+01. Another possible candidate is a 34 ms X-ray pulsar with a Be star companion in the error circle of 3EG J0634+0521 (Kaaret et al. 1999; Cusumano et al. 2000). Additional EGRET sources have recently been tentatively identified with known radio pulsars based on the detection of a corresponding pulsed $`\gamma `$-ray signal, namely 3EG J1048–5840 with PSR B1046–58 (Kaspi et al. 2000), and the millisecond pulsar PSR J0218+4232 with 3EG J0222+4253 (Kuiper et al. 2000). X-ray synchrotron nebulae that are inferred to be powered by pulsars have been detected in the error circles of the EGRET sources 2EG J1811–2339 (Oka et al. 1999) and 3EG J1420–6038 (Roberts & Romani 1998; Roberts et al. 1999). PSR B1951+32 in the supernova remnant CTB 80 has been detected by EGRET (Ramanamurthy et al. 1995) even though it does not exceed the threshold for inclusion in the EGRET catalogs. It has also been suggested that accreting neutron-star binaries might be EGRET sources, although so far only one example has been found in the possible detection of intermittent $`\gamma `$-rays from Cen X-3 (Vestrand, Sreekumar, & Mori 1997).
Informed by the pulsar scenario, we are studying several EGRET sources that are at “intermediate” Galactic latitudes, $`3^{}<|b|<8^{}`$, and that are not apparently variable. This strategy increases the likelihood that (a) the EGRET source is real, (b) its position is not affected by errors in the diffuse emission model or neighboring weak sources, (c) it is relatively nearby, (d) the intervening column density is not too large for soft X-ray observations, and (e) the corresponding optical fields are not too crowded. The absence of variability is important, since the known $`\gamma `$-ray pulsars show little if any long-term variability, while the blazars often flare dramatically.
One of our targets is 3EG J2227+6122, a source at Galactic coordinates $`(\mathrm{},b)=(106.^{}5,3.^{}2)`$ with a 95% error circle of radius $`0.^{}46`$ (Hartman et al. 1999). Its average flux is $`4.1\times 10^7`$ photon cm<sup>-2</sup> s<sup>-1</sup> ($`>100`$ MeV) and its photon spectral index is $`2.24\pm 0.14`$. It is not obviously variable (McLaughlin et al. 1996). The total Galactic 21 cm column density in this direction is $`8.2\times 10^{21}`$ cm<sup>-2</sup> (Stark et al. 1992). Prior to this work, there were no known pulsars or blazar-like radio sources in this field, and no prominent X-ray sources. In this paper we present the results of X-ray, radio, and optical observations of the region of 3EG J2227+6122 which together point to a possible identification.
## 2 Observations
### 2.1 ROSAT HRI Survey
The ROSAT HRI was used to cover the 95% error circle of 3EG J2227+6122 in four overlapping pointings, as shown in Figure 1. These observations were made during 1996 August 7–13. Exposure times ranged between 14,000 and 19,000 s. A total of six compact X-ray sources were detected in this field to a limiting count rate of approximately $`1\times 10^3`$ counts s<sup>-1</sup>. This limit corresponds to an intrinsic flux of $`4\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup> for a thermal plasma of $`T=3\times 10^6`$ K and $`N_\mathrm{H}=2\times 10^{20}`$ cm<sup>-2</sup>, spectral parameters that are expected of stellar coronal X-ray sources which should be the dominant population in this field. The X-ray source positions, count rates, and optical counterpart data are listed in Table 1. Optical magnitudes and positions are taken from the USNO–A2.0 catalog (Monet et al. 1996).
### 2.2 Optical Identifications of ROSAT Sources
Optical spectroscopic identifications were made for five of the six ROSAT HRI sources using the KPNO 2.1m telescope and Goldcam spectrograph. Four of these sources are bright K and M type stars. The fifth is a bright emission-line star which is classified as a Herbig Ae/Be type by Hang, Liu, & Xia (1999); it is also an IRAS source. Our optical spectrum of this star is essentially identical to that of Hang et al., and it confirms their classification in detail. This star also illuminates a prominent nebulosity which is visible on the Palomar Observatory Sky Survey (POSS) plates. Our spectroscopy and subsequent H$`\alpha `$ imaging show that this nebula is bright in H$`\alpha `$. Herbig Ae/Be stars are commonly detected as X-ray sources (e.g., Zinnecker & Preibisch 1994), but we have no reason to suspect that this one is the origin of 3EG J2227+6122. We also note that the X-ray positions of all five identified sources are coincident with their optical positions to the expected statistical accuracy of $`1^{\prime \prime }3^{\prime \prime }`$ (see Figure 2). Furthermore, any systematic offset between the average X-ray and optical positions is $`0.^{\prime \prime }6`$, which is not significantly different from zero. Therefore, we do not find it necessary to make any a posteriori adjustments to the X-ray astrometry.
The sixth HRI source, RX J2229.0+6114, remains unidentified optically. The remainder of this paper is largely concerned with the properties of this source and the evaluation of its credentials as a possible identification of 3EG J2227+6122. Within a conservative $`5^{\prime \prime }`$ error radius, which is justified empirically by the data in Figure 2, there is nothing at this location on the POSS plates. Figure 3 shows images of this field to a limit of $`R=24.5`$ and $`B=24.0`$ that we obtained using the MDM 2.4m telescope. The nearest bright object is a star of $`R16.8`$ which lies $`5.^{\prime \prime }6`$ southeast of the X-ray position; it is labeled star A in Figure 3. An optical spectrum which we obtained on the KPNO 2.1m telescope shows that it is in fact a highly reddened A star, an unlikely counterpart considering the hard X-ray properties of this source discussed below. The brightest object in the X-ray error circle, just north of star A, has $`R=21.3`$. A Keck spectrum of this object kindly obtained by R. H. Becker shows no emission lines or other interesting features. Although we cannot firmly classify the spectrum, it is possibly an ordinary star or galaxy. On the other hand, if it is a truly featureless spectrum, then it could be a Crab-like pulsar or a BL Lac object. Nonetheless, we regard this X-ray source as unidentified to a limit of $`R21.3`$ and $`B>24`$, and possibly to $`R23`$ if the $`R=21.3`$ object is not its counterpart.
### 2.3 Radio Observations
We created a 20 cm image of the entire EGRET error circle for 3EG J2227+6122 by constructing a mosaic from 13 snapshots obtained on 31 March 1996 using the Very Large Array (VLA<sup>1</sup><sup>1</sup>1The VLA, part of the National Radio Astronomy Observatory, is operated by Associated Universities, Inc., under cooperative agreement with the National Science Foundation.) in its C configuration. The synthesized beam yields a resolution of $`15^{\prime \prime }`$ FWHM; the median rms in the image is 0.12 mJy. Subsequently, the NRAO/VLA Sky Survey (NVSS – Condon et al. 1998) covered this field to a somewhat lower sensitivity (rms $`0.5`$ mJy) at an angular resolution of $`45^{\prime \prime }`$. In addition, we have examined images from the 92 cm Westerbork Northern Sky Survey (WENNS – Rengelink et al. 1997), the Greenbank 20 cm single-dish image (Condon & Broderick 1991), and the Greenbank 6 cm catalog (Becker, White, & Edwards 1991), as well as observing a portion of the field again at 6 cm on 12 January 1997 with the VLA in its D configuration.
Only one radio source in the EGRET error circle is coincident with an X-ray source: the single unidentified source RX J2229.0+6114. An image of this radio source constructed from the $`I`$, $`Q`$, and $`U`$ maps of the NVSS is displayed in Figure 4$`a`$ where it is seen that RX J2229.0+6114 lies at the center of an incomplete circular radio shell with a diameter of $`3.^{}5`$. The NVSS catalog lists two components for this source with a combined total flux density of $`73\pm 5`$ mJy. More remarkable is the high degree of linear polarization present throughout the shell. The NVSS catalog lists a polarized flux density of $`21`$ mJy, albeit with large reported errors. However, we have constructed a polarization map from the archived $`Q`$ and $`U`$ data, and the signal is unmistakable. We find a peak polarized flux density of 3.9 mJy per beam (the map rms is 0.23 mJy) and an integrated polarized flux density of 17 mJy, yielding a polarized fraction of $`25\%`$. In contrast, the 20 mJy source $`4^{}`$ south of the shell has a maximum polarized signal of 0.8 mJy, yielding an upper limit to its polarization of $`4\%`$. Other sources in the field are also unpolarized at the 5% level. The polarization vectors are approximately radial throughout the shell, implying a tangential magnetic field. However, it is important to note that an unknown amount of Faraday rotation along the line of sight to the source could alter this interpretation.
Our additional observations of this source at 20 and 6 cm confirm its high degree of polarization. While our higher resolution data clearly over-resolve the source, leading to substantial missing flux density, both images also yield polarized fractions of $`30\%`$; the 6 cm image is displayed in Figure 4$`b`$. The orientation of the polarization vectors at 6 cm are similar to those in the NVSS image, although they are not reliable since no polarization calibration was carried out during this observation. Note that the source $`4^{}`$ south of the shell is resolved in this map into an elongated structure $`20^{\prime \prime }`$ in extent, reminiscent of an extragalactic double radio source; its steep spectral index of $`\alpha =0.9`$ (where $`S_\nu \nu ^\alpha `$) derived from these data in conjunction with the WENNS catalog entry is consistent with this interpretation.
No source is listed in the WENNS catalog at the position of radio shell. However, a map extracted from the archive shows a clear excess coincident with the shell; an image smoothed with a $`60^{\prime \prime }`$ Gaussian yields a crude estimated flux density of $`35`$ mJy, which is probably uncertain by a factor of 2. It is clear from this image, however, that while the southern source has peak and integrated flux densities of 51 mJy and 60 mJy, respectively, the shell source is not significantly brighter at 92 cm than it is at 20 cm.
The shell source is not detected in the Greenbank 20 cm images because of confusion with bright diffuse emission nearby, but is clearly seen in the Greenbank 6 cm maps; its flux density is listed in Becker et al. (1991) as 80 mJy. Since the $`3.^{}5`$ beam of the Greenbank telescope at this wavelength is well-matched to the source size, this is probably the most reliable measure of the source flux density, as all of the interferometric measurements resolve out some fraction of the flux. Our scaled-array observations with $`15^{\prime \prime }`$ beams yield the same flux densities at the two frequencies to within the relatively large errors. Thus, we conclude that all available measurements are consistent with a flat spectral index for the radio shell from 92 cm through 6 cm, with an integrated flux density of $`80`$ mJy and a $`25\%`$ polarized fraction.
There are no other notable radio sources in the error circle of 3EG J2227+6122. The brightest source, with a flux density at 6 cm of 494 mJy (Becker et al. 1991), is the H II region Sharpless 141. Most significant for the classification of the $`\gamma `$-ray source, however, is the fact that there is no compact, flat-spectrum radio source in this field comparable to the well-identified EGRET blazars, all of which have 6 cm flux densities in excess of 400 mJy (Mattox et al. 1997). The upper limit on such a source in the field of 3EG J2227+6122 is $`20`$ mJy, the flux density limit of the Becker et al. catalog.
### 2.4 H$`\alpha `$ Imaging
To search for further evidence concerning the nature of the radio nebula and the compact X-ray source, we obtained H$`\alpha `$ images of an $`7.^{}3\times 7.^{}3`$ region around VLA J2229.0+6114 using the MDM 2.4m telescope and a 39 Å wide filter centered at 6563 Å. Figure 5 shows the combined image. Diffuse H$`\alpha `$ structures are present throughout, with a peak surface brightness of $`1.7\times 10^{16}`$ erg cm<sup>-2</sup> s$`{}_{}{}^{1}\mathrm{arcsec}_{}^{2}`$ above the background level in the northwest and southeast corners of the image. This value is comparable to the average of the diffuse ionized gas at $`b=0^{}`$ near this location ($`6.3\times 10^{17}`$ erg cm<sup>-2</sup> s$`{}_{}{}^{1}\mathrm{arcsec}_{}^{2}`$; Reynolds 1985). However, there is no structure that appears correlated with either the radio shell or the location of the X-ray source. The $`1\sigma `$ noise level in this image is $`1.2\times 10^{17}`$ erg cm<sup>-2</sup> s$`{}_{}{}^{1}\mathrm{arcsec}_{}^{2}`$; the implications of the lack of H$`\alpha `$ emission at this level from VLA J2229.0+6114 or RX J2229.0+6114 will be discussed in §3.
### 2.5 ASCA Observation of the Unidentified Source RX J2229.0+6114
To investigate further the nature of RX J2229.0+6114, we obtained an ASCA observation beginning on 1999 August 4. A total of 114,500 s of exposure was obtained with each of the two GIS detectors, and 97,600 s with each of the two SIS detectors operated in 1-CCD mode. A prominent hard X-ray source was detected at the position (J2000) $`22^h29^m05.^s9,+61^{}14^{}16^{\prime \prime }`$ (corrected for the ASCA temperature variation by the method of Gotthelf et al. 2000). This position is consistent with that of the ROSAT HRI source to within $`8^{\prime \prime }`$, and is well within the $`12^{\prime \prime }`$ radius ASCA error circle (at 90% confidence). A contour map of the combined ASCA GIS images is superposed on the ROSAT HRI images in Figure 1, and a more detailed view of the GIS image is shown in Figure 6. Analysis of the SIS and GIS photons from RX/AX J2229.0+6114, extracted using standard methods, shows that the spectrum (Figure 7) is best fitted by a power law of photon index $`\mathrm{\Gamma }=1.51\pm 0.14`$ at 90% confidence, and $`N_\mathrm{H}=(6.3\pm 1.3)\times 10^{21}`$ cm<sup>-2</sup>; the confidence contours of these spectral parameters are shown in Figure 8. The intrinsic 2–10 keV flux is $`1.56\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>.
Some care was required in the extraction of source and background regions for this analysis since there is also diffuse X-ray flux in this region. Weak, diffuse emission surrounding the compact source as well as to the northwest of it appears to be much softer than the compact source, and its uncertain distribution is a significant source of systematic error in the spectral parameters of RX/AX J2229.0+6114. We chose a radius of $`4^{}`$ for the source region in all detectors, and background regions which are as close as possible, but which exclude a region of radius $`5^{}.25`$ around the source. As we change the source and background extraction regions, the spectral index of RX/AX J2229.0+6114 can become as steep as 1.8, which we consider the maximum systematic deviation from the best fitted value of $`1.51\pm 0.14`$. There is also a weak, soft ASCA source coincident with the Herbig Ae/Be star RX J2226+6113. No other significant sources are seen in the ASCA images to a flux limit of $`6\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup>.
The fitted $`N_\mathrm{H}=(6.3\pm 1.3)\times 10^{21}`$ cm<sup>-2</sup> of RX/AX J2229.0+6114 is almost equal to the total Galactic 21 cm column in this direction, $`N_\mathrm{H}=8.4\times 10^{21}`$ cm<sup>-2</sup> (Stark et al. 1992), indicating that the X-ray source is at least 2 kpc distant, and possibly much farther. In the following sections we adopt a fiducial distance of 3 kpc as typical for the X-ray measured $`N_\mathrm{H}`$ if RX/AX J2229.0+6114 is Galactic. Absorption accounts for the relative weakness of the source in the HRI. The fitted ASCA spectral parameters would extrapolate to an HRI count rate of $`(5.5\pm 1.5)\times 10^3`$ counts s<sup>-1</sup>, approximately twice that observed. This difference can be taken either as marginal evidence of long-term variability of this X-ray source, or as an indication that some extended emission is escaping detection by the HRI.
We also searched the ASCA spectrum of RX/AX J2229.0+6114 for an emission line of Fe K$`\alpha `$. No such line is detected, and 95% confidence upper limits on its EW range from 300 eV for a narrow line at $`z=0`$, to 50 eV for a narrow line at $`z=0.12`$. For a broad line (Gaussian $`\sigma =0.5`$ keV), the corresponding limits are 440 eV at $`z=0`$ and 150 eV at $`z=0.12`$. Beyond $`z=0.12`$, the X-ray luminosity of the source would exceed $`10^{44}`$ erg s<sup>-1</sup>, and such luminous AGNs do not usually exhibit an Fe K$`\alpha `$ line.
The elapsed time of the ASCA observation was 67 h. During this time there was no evidence for short-term variability on time scales of hours. The ASCA GIS bit assignments were configured for high time resolution in order to search for pulsations. Approximately 54,600 s of data were obtained with 3.9 ms resolution, and 59,100 s with 0.5 ms resolution. These were searched, without success, for periodic pulsations for all periods as short as 10 ms using the Rayleigh test. In addition to a coherent search of the entire photon list (extracted from a region of radius $`4^{}`$ around RX/AX J2229.0+6114 ), the period search was performed on individual segments of 30,000 s in length in case there is a short-period pulsar with a large $`\dot{P}`$ or acceleration. For example, the Crab pulsar’s $`\dot{P}`$ term contributes –0.17 cycles over an elapsed time of 30,000 s. Further searches of the data for a high $`B`$ pulsar are in progress, and the results will be reported in a subsequent paper.
## 3 Possible Interpretations of RX/AX J2229.0+6114 and VLA J2229.0+6114
### 3.1 A Chance Coincidence?
We considered the hypothesis that RX/AX J2229.0+6114 is a chance superposition of a Galactic or extragalactic X-ray source unrelated to the radio shell VLA J2229.0+6114. The X-ray properties of RX/AX J2229.0+6114 are by themselves inconclusive as to its Galactic or extragalactic origin. Power-law fits to the ASCA spectrum require a column density slightly less than the maximum Galactic 21 cm value, which allows either a background AGN or a distant Galactic source. The power-law photon index $`\mathrm{\Gamma }=1.5`$ is harder than that of most radio-quiet QSOs, but is typical for a $`\gamma `$-ray pulsar (Wang et al. 1998). Although there is no prominent Fe K$`\alpha `$ emission line as would be expected from a low-redshift Seyfert galaxy, we cannot rule out that this is a luminous QSO for which Fe K$`\alpha `$ is generally not seen. The absence of variability on short and long time scales is also ambiguous, since AGNs are often but not always variable during an observation of this length. The spatial resolution and sensitivity of the ASCA SIS are not adequate to determine if RX/AX J2229.0+6114 is spatially extended, which if true, would be strong evidence of a Galactic pulsar-powered nebula. Similarly, it is not possible to conclude whether or not the diffuse soft X-ray emission surrounding RX/AX J2229.0+6114 is associated with it.
The absence of an emission-line optical counterpart to a limit of $`R>23`$ somewhat favors a neutron star, although the Galactic absorption $`A_R=4.6`$ mag (Schlegel et al. 1998) allows for the possibility that a faint object in the error circle in Figure 3 is a QSO of dereddened $`R18.4`$. Such a QSO would have $`\alpha _{\mathrm{ox}}1.0`$, which would be near the extreme of X-ray loudness among QSOs (Wilkes et al. 1994). Thus, we regard the absence of a suitable optical counterpart for RX/AX J2229.0+6114 as somewhat discouraging an AGN classification unless it is a new type of extreme $`\gamma `$-ray quasar (see §4).
### 3.2 An H II Region?
We also considered the hypothesis that, if the X-ray source RX J2229.0+6114 is unrelated to the radio shell, then the latter might be an H II region. The principle reason for entertaining this possibility is the flat radio spectrum of VLA J2229.0+6114, which is reminiscent of thermal bremsstrahlung for which $`S_\nu \nu ^{0.1}`$ in the radio. While the high degree of polarization strongly favors nonthermal emission, we also have independent evidence against the H II region hypothesis from our optical imaging. The argument proceeds as follows. We can use the observed radio flux to estimate the required electron density and ionizing flux. The thermal bremsstrahlung emissivity is
$$ϵ_\nu =6.8\times 10^{38}\frac{Z^2n_en_i}{T^{1/2}}e^{h\nu /kT}g_{ff}(\nu ,T)\mathrm{erg}\mathrm{cm}^3\mathrm{s}^1\mathrm{Hz}^1.$$
$`(1)`$
Assuming $`T=8,000`$ K and an observing frequency $`\nu =4.86\times 10^9`$ Hz, the Gaunt factor $`g_{ff}4.91`$. Approximating the nebula as a sphere of ionized hydrogen of radius $`r_s=\theta d/2`$ where $`\theta 3^{}`$ is its angular diameter, the flux density received at Earth is
$$S_\nu =0.095\left(\frac{d}{3\mathrm{kpc}}\right)n_e^2\mathrm{mJy}.$$
$`(2)`$
In this case, since we observe $`S_\nu 80`$ mJy, the mean electron density in the nebula must be $`n_e30`$ cm<sup>-3</sup>.
The rate of ionizations $`Q`$ is related to the Strömgren radius $`r_s`$ and the case B recombination coefficient $`\alpha _B(T)3.1\times 10^{13}`$ cm<sup>3</sup> s<sup>-1</sup> via
$$Q=\frac{4\pi }{3}r_s^3n_en_+\alpha _B(T)=8.4\times 10^{43}\left(\frac{d}{3\mathrm{kpc}}\right)^3n_e^2\mathrm{s}^1.$$
$`(3)`$
Combining Equations (2) and (3), we can eliminate $`n_e`$ and solve for $`Q`$ as function of the observed radio flux and the unknown distance $`d`$,
$$Q=7.1\times 10^{46}\left(\frac{d}{3\mathrm{kpc}}\right)^2\left(\frac{S_\nu }{80\mathrm{mJy}}\right)\mathrm{s}^1.$$
$`(4)`$
For a range of plausible distances between 2 and 8 kpc, $`Q`$ is in the range $`(0.325.0)\times 10^{47}`$ s<sup>-1</sup>, which corresponds to stars of spectral type B0. The absolute visual magnitudes of such stars are in the range –4.1 to –4.4. Taking into account an estimated extinction $`A_V=1.4`$ mag kpc<sup>-1</sup> (up to a maximum of $`A_V=5.7`$), the apparent magnitude of the exciting star should fall in the range $`m_V=10.315.8`$.
The main difficulties with the H II region hypothesis are the absence of H$`\alpha `$ recombination radiation and the lack of a suitable exciting star. We describe each of these failed predictions in turn. Corresponding to the radio bremsstrahlung emission there should be H$`\alpha `$ with emission coefficient $`ϵ(\mathrm{H}\alpha )=3.2\times 10^{25}`$ erg cm<sup>3</sup> s<sup>-1</sup>, giving a total observed flux of
$$F(\mathrm{H}\alpha )=8.2\times 10^{14}\left(\frac{d}{3\mathrm{kpc}}\right)n_e^2\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1.$$
$`(5)`$
It is more useful to relate the expected surface brightness in H$`\alpha `$ directly to the surface brightness in the radio image by combining Equations (2) and (5). Independent of any other variables, the proportionality is
$$\mathrm{\Phi }(\mathrm{H}\alpha )=8.6\times 10^{16}\mathrm{\Sigma }(4.86\mathrm{GHz}),$$
$`(6)`$
where $`\mathrm{\Phi }(\mathrm{H}\alpha )`$ is the surface brightness of H$`\alpha `$ in erg cm<sup>-2</sup> s<sup>-1</sup> arcsec<sup>-2</sup> and $`\mathrm{\Sigma }(4.86\mathrm{GHz})`$ is the radio surface flux density in $`\mu `$Jy arcsec<sup>-2</sup>. The peak observed $`\mathrm{\Sigma }(4.86\mathrm{GHz})`$ is $`11\mu `$Jy arcsec<sup>-2</sup> in the northwest sector of the nebula, which predicts $`\mathrm{\Phi }(\mathrm{H}\alpha )=9.6\times 10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup> arcsec<sup>-2</sup>. Even if we allow for the maximum extinction along the line of sight (4.6 mag at H$`\alpha `$), we would still expect to see a surface flux of $`1.4\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup> arcsec<sup>-2</sup> at this location, which should be a prominent feature in Figure 5. Instead, there is no sign of such emission to a limit at least one order of magnitude smaller. The absence of H$`\alpha `$ emission at this level is difficult to understand if the radio emission is truly thermal bremsstrahlung at $`T8,000`$ K.
The second failing of the H II region hypothesis is the absence of an exciting star within the radio shell. The radio flux requires a B0 star of $`10.3m_V15.8`$. We have ruled out spectroscopically the bright star of $`R13`$ at the bottom of Figure 3; it is an early K star. The A star of $`R=16.8`$ closest to the X-ray source is also inadequate, as are half a dozen red stars of $`R1617`$ which lie within the area of the radio shell. Thus, even after taking into account the maximum extinction along this line of sight, it is not likely that we have missed a star which could excite such an H II region at any distance less than 10 kpc. And at a larger distance than this, one would be suspicious of the presence of an H II region more than 500 pc from the Galactic plane. We consider that an H II region classification of VLA J2229.0+6114 is therefore ruled out by its high polarization and lack of corresponding optical evidence. Since we also lack an optical counterpart of RX/AX J2229.0+6114, we proceed to consider scenarios in which these sources could be Galactic objects associated with each other and powered by a neutron star.
### 3.3 A Young Supernova Remnant?
If VLA J2229.0+6114 and RX/AX J2229.0+6114 are associated, then it is highly likely that they reside in the Galaxy. We consider here whether the radio shell can be a small-diameter supernova remnant (SNR) or a pulsar wind nebula. The SNR hypothesis encounters an immediate difficulty, since shell components generally have steeper radio spectra characterized by $`\alpha =0.40.7`$, where $`S_\nu \nu ^\alpha `$. The compact cores of plerionic (Crab-like) remnants have flatter radio spectra ($`\alpha =0.00.3`$), but they have center-filled morphologies.
The observed $`25\%`$ linear polarization is high, but not unprecedented in Galactic SNRs. The integrated polarization of the Crab Nebula, for example, is $`10\%`$, but some features in the nebula show up to 50% linearly polarized flux (Wilson, Samarasinha, and Hogg 1985). The Crab-like remnant G21.5–0.9 exhibits 20–30% polarization in a circumferential ring, although the integrated polarized fraction for the whole remnant is somewhat lower (Becker & Szymkowiak 1981). The typical polarized fraction for shell-like remnants is 5–10%. Extragalactic radio sources also exhibit polarization, but typical values are a few percent, and no extended source with a polarized fraction $`>20`$% at 20 cm has ever been detected. Irrespective of these comparisons, the high degree of linear polarization observed leads to the inescapable conclusion that the radio source associated with RX J2229.0+6114 is nonthermal.
There is no clear manifestation of the shell at any energy above the radio. The X-rays could be dominated by a point source, and our search for optical emission-line filaments, which was sensitive in the velocity range $`\pm 1,000`$ km s<sup>-1</sup>, yielded null results as described above. Nevertheless, it is clear that both the radio and the X-ray sources are nonthermal. The X-rays could be in part magnetospheric synchrotron emission from a young neutron star, and in part a compact plerion that is unresolved by the ASCA SIS. One way that a flat radio spectrum can be accommodated is if the characteristic synchrotron frequency of the lowest energy electrons is comparable to or greater than the observing frequency, with corresponding constraints on the magnetic field strength and electron energies as described in §3.4 below. $`S_\nu \nu ^{+1/3}`$ obtains in the low-frequency limit. Thus, there is enough evidence to warrant a SNR interpretation for VLA J2229.0+6114, albeit one with extreme properties.
### 3.4 A Pulsar Bow-Shock Nebula?
An alternative interpretation of the radio shell is nonthermal emission from a shock between a relativistic pulsar wind and the surrounding interstellar medium (ISM). Two scenarios can be considered for the required power in the pulsar wind. The first assumes that the shock is expanding with velocity $`v_s`$, and is driven by the difference in the interior and exterior pressure. The second method assumes that the radio shell is a bow shock which travels at the velocity of the pulsar $`v_p`$, and whose apex is located where the pulsar wind pressure equals $`\rho _0v_p^2`$. Since the incomplete radio shell in fact resembles a bow shock, we chose the latter scenario for consideration. Thus, the spin-down power $`\dot{E}`$ is related to the ambient ISM density $`\rho _0(1.4n_\mathrm{H}m_p)`$ and the apex distance $`r_0`$ via
$$\dot{E}=4\pi r_0^2c\rho _0v_p^2=1.9\times 10^{38}\left(\frac{n_\mathrm{H}}{0.1}\right)\left(\frac{d}{3\mathrm{kpc}}\right)^2\left(\frac{v_p}{100\mathrm{km}\mathrm{s}^1}\right)^2\mathrm{erg}\mathrm{s}^1.$$
$`(7)`$
Here we have measured $`1.^{}7`$ for the apex distance, and we assume a relatively low-density ISM. The rather large size of the radio shell means that the required $`\dot{E}`$ is of order $`10^{38}`$ erg s<sup>-1</sup> for any reasonable distance, pulsar velocity, or ISM density.
It is worth noting that the direction of the pulsar’s velocity vector under this hypothesis is nearly perpendicular to and away from the Galactic plane, consistent with its birth as a young star in the disk. While not in itself strong evidence that this is a young pulsar, a velocity in the opposite direction would have been difficult to accommodate, as this source is already at just about the maximum $`z`$-height expected for a Population I object ($`z=160`$ pc for $`d=3`$ kpc).
Under the bow shock hypothesis we assume that the radio flux is nonthermal emission from the shocked pulsar wind itself – i.e., the reverse shock rather than the shocked ISM – as theorized by Hester & Kulkarni (1988) in their model of the core of the supernova remnant CTB 80. A Crab-like pulsar wind is thought to contain relativistic electrons with $`\gamma 1\times 10^6`$ (Kennel & Coroniti 1984a,b). We hypothesize an electron-positron wind, and we suppose that the shock produces a relativistic Maxwellian distribution of particle energies rather than a power law, and an equipartition magnetic field of $`B4\times 10^5`$. Then the synchrotron power (in $`\nu F_\nu `$) will be emitted mostly near $`\nu 25\nu _c4\times 10^{15}`$ Hz (Tavani 1996) where $`\nu _c`$ is the characteristic frequency $`\nu _c=(3\gamma ^2eB)/(4\pi m_ec)`$. If a suprathermal power-law tail develops, then the frequency of the synchrotron peak can move down a factor of 10 to $`\nu 2.5\nu _c4\times 10^{14}`$ Hz. Thus, we may expect little X-ray emission coincident with the radio shell. The low-energy limit of the synchrotron emission from such a thermal distribution has an asymptotic spectrum of the form $`F_\nu \nu ^{+1/3}`$, just as in the non-thermal case.
Assuming that the observed radio flux extrapolates with $`\alpha 0.0`$ up to $`10^{16}`$ Hz, the total synchrotron power from the bow shock is $`8.6\times 10^{36}`$ erg s<sup>-1</sup>, which is a significant fraction of the $`\dot{E}`$ inferred from the radius of the shell. We would not easily detect the optical synchrotron emission from the reverse shock because its surface brightness would be at most $`11\mu `$Jy arcsec<sup>-2</sup> as it is in the radio, which corresponds to 21.1 mag arcsec<sup>-2</sup> in the $`R`$ band, or 25.7 mag arcsec<sup>-2</sup> after correcting for the maximum extinction. Any steepening of the spectral index above the radio band would further decrease this optical estimate.
The majority of $`\dot{E}`$ may go into accelerating and heating the shocked ISM. Although the forward shock from such pulsar winds is sometimes detected in H$`\alpha `$, the low density of the ISM in this case, and the relative inefficiency of the non-radiative shock which is usually inferred, would predict much less H$`\alpha `$ emission than was required in the H II region scenario for VLA J2229.0+6114. A simple scaling relation can be derived to approximate the average surface brightness in H$`\alpha `$ of a nonradiative shock. According to Raymond (1991), it is the neutral fraction of hydrogen passing freely through the shock that emits on average 0.2 H$`\alpha `$ photons per atom via collisional excitation before it is ionized by the hot, shocked ISM. This number is relatively independent of shock velocity because the collisional excitation and collisional ionization cross sections scale similarly with temperature. If the bulk of the H$`\alpha `$ is emitted in a hemispherical region of radius $`r_0`$, then the average H$`\alpha `$ surface brightness is
$$\mathrm{\Phi }(\mathrm{H}\alpha )=1.1\times 10^{18}\left(\frac{n_\mathrm{H}}{0.1}\right)\left(\frac{v_p}{100\mathrm{km}\mathrm{s}^1}\right)\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1\mathrm{arcsec}^2.$$
$`(8)`$
After including the effects of extinction, this value is 2–3 orders of magnitude smaller than the limiting sensitivity of our H$`\alpha `$ image.
In this scenario, neither the reverse shock nor the forward shock are energetic enough to emit the hard X-rays which are seen from RX/AX J2229.0+6114. Since the X-rays are apparently more compact in size than the radio shell, we would assume that they are entirely magnetospheric synchrotron emission from a young neutron star and/or a compact plerion that is unresolved by the ASCA SIS. Its power-law spectral index $`\mathrm{\Gamma }=1.5`$ is compatible with this interpretation, as is its lack of variability. The 2–10 keV luminosity of this source is $`1.7\times 10^{33}(d/3\mathrm{kpc})^2`$ erg s<sup>-1</sup>, which is $`10^5`$ of the pulsar spin-down power as inferred from the bow-shock interpretation of the radio shell. This is perhaps one of the weaknesses of this model, since the ratio $`L_x/\dot{E}`$ for rotation-powered pulsars is typically observed to be $`10^410^3`$. While $`\dot{E}`$ can perhaps be reduced by lowering the ambient density $`n_\mathrm{H}`$ to 0.01, the unusual flat radio spectrum of VLA J2229.0+6114 remains a significant mystery. Thus, none of the scenarios that we have considered is without difficulty.
### 3.5 A Pulsar Wind Bubble?
Some of the problems associated with the large spin-down power inferred from the bow-shock interpretation of the radio source VLA J2229.0+6114 could be ameliorated if it is a bubble confined by the static pressure $`nkT`$ of the ISM instead of the dynamic pressure $`\rho _0v_p^2`$. This scenario would require the pulsar to have fortuitously a smaller velocity than the thermal speed of the ISM, but if it did, the inferred spin-down power could be much smaller,
$$\dot{E}=4\pi r_0^2c(n_e+n_i)kT=2.4\times 10^{36}\left(\frac{n_\mathrm{H}}{0.1}\right)\left(\frac{d}{3\mathrm{kpc}}\right)^2\left(\frac{T}{10^4\mathrm{K}}\right)\mathrm{erg}\mathrm{s}^1.$$
$`(9)`$
We assume an ambient temperature of $`10^4`$ K as appropriate if the ISM were pre-ionized by the recent UV flash from the supernova which produced the neutron star, or if an older pulsar is now surrounded by the cooling remnant of the past supernova shock. In this case, the X-ray luminosity could indeed be $`10^410^3`$ of $`\dot{E}`$ as is observed from most other pulsars. We would still not expect to detect diffuse optical continuum or H$`\alpha `$ emission associated with the radio shell, although the extrapolated radio spectrum of the shell would have to turn down below $`10^{15}`$ Hz in order not to exceed $`\dot{E}`$. The pulsar would then be required to have a very high efficiency of conversion of $`\dot{E}`$ into $`\gamma `$-rays if it were to be the counterpart of the EGRET source 3EG J2227+6122. The observed photon flux of 3EG J2227+6122 above 100 MeV, $`4.1\times 10^7`$ photon cm<sup>-2</sup> s<sup>-1</sup> translates, for the photon spectral index of 2.14, into a luminosity of $`3.7\times 10^{35}(d/3\mathrm{kpc})^2`$ erg s<sup>-1</sup>, approximately 15% of $`\dot{E}`$ from Equation (9). Such high $`\gamma `$-ray efficiencies are also inferred for the adolescent pulsars Geminga and PSR B1055–52; they could be reduced by a modest amount of anisotropy in the $`\gamma `$-ray beam pattern. It could be argued that the shape of the radio source VLA J2229.0+6114 resembles an incomplete spherical bubble as much as it does a cometary bow-shock surface.
## 4 Implications of the Identification (or Not) of 3EG J2227+6122
Despite the lack of a completely satisfactory theory linking the X-ray source RX/AX J2229.0+6114 and the radio shell VLA J2229.0+6114, we are encouraged that their properties are not inconsistent with what we would expect for a pulsar counterpart of 3EG J2227+6122. Of all the pulsars detected by EGRET, the Crab is the only one whose $`\gamma `$-ray spectral index ($`\mathrm{\Gamma }=2.19\pm 0.02`$) is identical to that of 3EG J2227+6122 ($`\mathrm{\Gamma }=2.24\pm 0.14`$). All of the less energetic pulsars have flatter spectra, and there is a clear trend of spectral flattening (and increasing $`\gamma `$-ray efficiency) with increasing age or decreasing spin-down power. Furthermore, since the observed $`\gamma `$-ray flux of the Crab is approximately 5 times that of 3EG J2227+6122, the estimated distance to RX/AX J2229.0+6114 of 3 kpc, 1.5 times that of the Crab, is basically what one would expect for a pulsar of half the spin-down power and similar $`\gamma `$-ray efficiency. At this intermediate Galactic latitude of $`3^{}`$, a distance of 3 kpc is just about the maximum that one would expect for a young Population I object which should be found within the ISM layer of the disk, and it is just about right to accommodate the observed X-ray absorption column of $`6.3\times 10^{21}`$ cm<sup>-2</sup>. Essentially, by selecting an EGRET source at this Galactic latitude, we are predisposed to finding a pulsar that is similar to the Crab but more distant. Accordingly, it is perhaps not surprising that the bow-shock interpretation of the radio shell VLA J2229.0+6114 requires $`\dot{E}2\times 10^{38}`$ erg s<sup>-1</sup>, about half that of the Crab. If this identification for 3EG J2227+6122 is correct, the main difference between it and the Crab would be the presence of a prominent bow shock, and the relative weakness of X-ray emission from the pulsar magnetosphere and/or plerion. The failure to detect X-ray pulsations from RX/AX J2229.0+6114 is mysterious, however, unless a compact plerion dominates over the pulsar emission.
We presented an alternative scenario in which the spin-down power of the putative pulsar could be much lower, $`2.4\times 10^{36}`$ erg s<sup>-1</sup>. Although this would allow an older and therefore more “common” class of pulsar, the 3 kpc distance would strain its $`\gamma `$-ray production efficiency if it behaves similarly to $`\gamma `$-ray pulsars in this regime of $`\dot{E}`$, such as Vela, PSR B1951+32, and PSR B1706–44. Those three pulsars have (isotropic) $`\gamma `$-ray efficiencies between 0.4% and 7%, while a $`\gamma `$-ray luminosity of $`3.7\times 10^{35}(d/3\mathrm{kpc})^2`$ erg s<sup>-1</sup> for 3EG J2227+6122 would require an efficiency of 15%. Perhaps this is not implausible, however, because the older pulsars Geminga and PSR B1055–52 have even higher apparent efficiency. Thus, it is possible that we are dealing with a highly efficient example of an intermediate age $`\gamma `$-ray pulsar.
The observed X-ray spectral index of RX/AX J2229.0+6114, $`\mathrm{\Gamma }=1.5`$, was in fact predicted for strong $`\gamma `$-ray pulsars by Wang et al. (1998) in the context of the outer-gap model. In this model, an outer-gap accelerator sends $`e^\pm `$ pairs flowing inward and outward along open magnetic field lines. These particles continuously radiate $`\gamma `$-rays by the curvature mechanism. When the inward flowing particles approach the surface of the star, the $`>100`$ MeV $`\gamma `$-rays that they emit convert into secondary $`e^\pm `$ pairs in the inner magnetosphere wherever $`B\mathrm{sin}\varphi >2\times 10^{10}`$ G, where $`\varphi `$ is the angle between the photon and the B field. Those secondary pairs must radiate away their energy instantaneously in the strong local B field. Such a synchrotron decay spectrum has $`\mathrm{\Gamma }=1.5`$ between $`E_{\mathrm{min}}=0.1`$ keV and $`E_{\mathrm{max}}=5`$ MeV. We emphasize that in this theory, the X-ray power-law is not supposed to be a simple continuation of the EGRET spectrum. Rather, it is a separate component radiated by the secondary $`e^\pm `$ pairs that are created when some of the primary $`\gamma `$-rays convert in the strong B field. The X-ray luminosity generated by this mechanism is
$$L_x(210\mathrm{keV})1.5\times 10^{31}f\left(\frac{0.1}{P}\right)^2\left(\frac{B_\mathrm{p}}{10^{12}}\right)^{1/2}\mathrm{erg}\mathrm{s}^1,$$
$`(10)`$
where $`P`$ is the rotation period, $`B_\mathrm{p}`$ is the surface polar magnetic field, and $`f`$ is the fraction of the Goldreich-Julian current that flows through the outer-gap accelerator back to the polar cap. This mechanism could make a significant contribution to the observed X-ray luminosity of RX/AX J2229.0+6114 only if $`P`$ is small, $`0.01`$ s, and if $`f`$ is of order unity. Otherwise, most of the X-rays must be generated by another mechanism, including perhaps an extended synchrotron nebula.
An alternative hypothesis, that RX/AX J2229.0+6114 is a twin of the intermediate-age pulsar Geminga, would be difficult to accommodate since Geminga’s $`\dot{E}=3.3\times 10^{34}`$ erg s<sup>-1</sup> is inadequate to account for the $`\gamma `$-ray luminosity of 3EG J2227+6122 at a distance of 3 kpc unless a beaming factor $`>10`$ is supposed. Also, the nonthermal X-ray luminosity of RX/AX J2229.0+6114, $`1.7\times 10^{33}(d/3\mathrm{kpc})^2`$ erg s<sup>-1</sup>, is 3 orders of magnitude more than that of Geminga.
If RX/AX J2229.0+6114 proves not to be the identification of 3EG J2227+6122, then the absence of any other X-ray candidate at the level of $`4\times 10^{14}`$ erg s<sup>-1</sup> is difficult to reconcile with any of the established classes of $`\gamma `$-ray emitters. If 3EG J2227+6122 is a pulsar but RX/AX J2229.0+6114 is not its counterpart, it implies that highly efficient (or highly beamed) $`\gamma `$-ray pulsars can avoid producing soft or hard X-rays at a level below $`10^4`$ of their apparent $`\gamma `$-ray luminosity. At least two mechanisms of X-ray emission have been observed to accompany all $`\gamma `$-ray pulsars at such levels or higher, and they were explained in the context of the outer-gap model by Wang et al. (1998). One is the synchrotron emission from secondary pairs as described above. The second is thermal emission from the heated polar caps that are impacted by the inward-going accelerated particles from the outer-gap accelerator, as well as from the original heat of formation emitted over the whole surface. If 3EG J2227+6122 were a low-luminosity pulsar like Geminga, then it must be less than 1 kpc distant in order to be detected by EGRET, yet the intervening column density must be greater than $`10^{21}`$ cm<sup>-2</sup> in order for its thermal soft X-rays not to reach us. Many of the unidentified EGRET sources could be similarly situated radio-quiet pulsars, made exceedingly difficult to identify in X-rays because of interstellar absorption.
One might still hypothesize that 3EG J2227+6122 is an extragalactic object. The one tempting piece of evidence along those lines is the steep $`\gamma `$-ray spectral index, which if it is not due to a Crab-like pulsar, is compatible with many of the known EGRET blazars. But as discussed above, there is no compact, flat-spectrum radio source in this field to an upper limit $`20`$ mJy, or 50–100 times fainter than the typical identified EGRET blazars. Furthermore, there is no extragalactic X-ray source in this field other than possibly RX/AX J2229.0+6114 itself, and it is radio quiet (as a compact source) at the 1 mJy level. Thus, there is no escaping the fact that 3EG J2227+6122 would be a member of a new class of $`\gamma `$-ray source if extragalactic. However, it would not necessarily be unique in this regard. Similar considerations concerning the high-latitude source 3EG J1835+5918 were recently presented by Mirabal et al. (2000), where even more stringent upper limits to its potential pulsar or blazar counterparts were obtained. Since many EGRET sources remain unidentified, more of them may prove, upon detailed study and especially after more precise localizations are obtained by GLAST, to be of a previously unrecognized type, for example, the frequently imagined “radio-quiet blazar”. It is of interest that several $`\gamma `$-ray loud quasars are seen to have flat X-ray spectra in the 2–10 keV band, with $`\mathrm{\Gamma }=1.31.5`$ (Tavecchio et al. 2000). Their broad-band spectra can be fitted with Inverse Compton jet models only if a significant proton component or Poynting flux outside the emission region is the main carrier of power. If RX/AX J2229.0+6114 were a radio-quiet blazar of this type, then its flat X-ray spectrum and lack of an optical counterpart might be just what is needed for the prototype of a new or extreme class of $`\gamma `$-ray loud AGN.
## 5 Conclusions
ROSAT and ASCA observations of the error circle of 3EG J2227+6122 reveal only one candidate X-ray counterpart with an intrinsic 2–10 keV flux of $`1.56\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, fitted by a power-law spectrum of photon index $`\mathrm{\Gamma }=1.51\pm 0.14`$ and $`N_\mathrm{H}=(6.3\pm 1.3)\times 10^{21}`$ cm<sup>-2</sup>. There is no other candidate to a flux limit of $`6\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. This X-ray source happens to coincide with a highly polarized, incomplete radio shell which resembles a bow-shock nebula or a wind-blown bubble. The X-ray measured column density implies a distance of $`3`$ kpc, although an extragalactic source has not been ruled out definitively. The flat radio spectrum is different from that of all other shell SNRs, and our optical searches have turned up no evidence for associated emission-line gas. On balance, we favor the hypothesis that RX/AX J2229.0+6114 is indeed an energetic pulsar counterpart of 3EG J2227+6122, and that VLA J2229.0+6114 is an associated nebula powered by a large fraction of the pulsar spin-down power. The best method for proving this and quantifying the age and energetics of this unique source would be to obtain more sensitive X-ray imaging and timing observations. We have planned both Chandra and XMM observations. Chandra has the ability to resolve a point source from the hypothesized compact synchrotron nebula, and to obtain a precise position for deeper optical follow-up. XMM with its high throughput may permit a pulsar discovery and a precise ephemeris to be developed, which could then lead to detection of pulses in the archival ASCA and EGRET data, and in future observations by GLAST. Deeper optical spectroscopy, and possibly infrared imaging, would constitute a more definitive test for an extragalactic counterpart of RX/AX J2229.0+6114 such as a Seyfert galaxy or a quasar. The absence of such a counterpart would itself be strong evidence that RX/AX J2229.0+6114 is a neutron star.
This work was supported by NASA grants NAG 5-3229 and NAG 5-7814. We thank R. H. Becker for obtaining an optical spectrum with the Keck telescope, and M. Eracleous for assistance with other spectroscopic observations at Kitt Peak. We also thank R. H. Becker for calling our attention to the important polarization information contained in the NVSS images.
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# Time of avalanche mixing of granular materials in a half filled rotated drum
## Abstract
The avalanche mixing of granular solids in a slowly rotated 2D upright drum is studied. We demonstrate that the account of the difference $`\delta `$ between the angle of marginal stability and the angle of repose of the granular material leads to a restricted value of the mixing time $`\tau `$ for a half filled drum. The process of mixing is described by a linear discrete difference equation. We show that the mixing looks like linear diffusion of fractions with the diffusion coefficient vanishing when $`\delta `$ is an integer part of $`\pi `$. Introduction of fluctuations of $`\delta `$ supresses the singularities of $`\tau (\delta )`$ and smoothes the dependence $`\tau (\delta )`$.
The evolution of granular materials in rotated drums is one of the most interesting topics of soft matter physics . Granular flows induced in such a way demonstrate a lot of intriguing phenomena including, e.g., segregation of different fractions, mixing, self-organization, self-organized criticality, etc . In principle, the system is three-dimensional, and the process features depend strongly on the particular characteristics of the grains involved. A theoretical description of these processes is still unsolved problem, and only the simplest cases were considered analytically.
Decrease of the dimensionality of the system by one simplifies the consideration, although does not make the arising problems to be less interesting. The avalanche mixing phenomenon observed in a slowly rotated 2D upright drum, i.e., in a flat disc, turned to be so impressive that one of figures of the paper has made a front page of the Nature.
Let us introduce the avalanche mixing problem . A partially filled 2D upright drum rotates very slowly around its axis (see Fig. 1). Small grains of the granular material move as a whole with the drum. There is, however, one exception: periodically, when the grains appear on the free surface, they fall down along it. If, initially, there is several separated fractions of the material, the grains of the different fractions may mix only in the avalanches which fall down along the free surface. One has to describe the evolution of the mixing, i.e., to describe how the process develops. The main quantity under study is the time of mixing $`\tau `$, that is the characteristic time after that the material will be mixed homogeneously. It may be introduced, since the relaxation to the final homogeneous state in such a process proceeds by an exponentional law.
There are two particular values of the angle of the free surface slope of a granular material – the angle of marginal stability and the angle of repose. (We assume that the free surface is straight. That is reasonable for a slowly rotated drum.) When the free surface slope achieves the angle of marginal stability, the grains from the sector of the angle $`\delta `$ near the upper half of a free surface fall down (see Fig. 1 – grains fall down from sector u to sector d). Here, $`\delta `$ is the difference between the angle of marginal stability and the angle of repose. The avalanche stops when the free surface slope approaches the repose angle, and the drum has to turn by the angle $`\delta `$ to launche the next avalanche. Therefore, the only reason for mixing in the problem under consideration is the granular flow (avalanches) between these wedges.
We assume that the fractions become to be mixed homogeneously after the fall of the avalanche – the mixture in sector d in Fig. 1. This case, in which the avalanche mixing proceeds most quickly, appears to be near the experimentaly realized situation . At least, using this assumption, one may obtain the minimal possible time of the avalanche mixing.
When $`\delta 0`$, the mixing was described in the frames of a simple geometrical approach , and an analytical theory describing the experimental observations surprisingly well has been proposed. Nevertheless, there is a very interesting particular case of the half filling, in which one certainly can not assume that $`\delta `$ is zero. (Note, that in the experiment of Metcalfe et al , $`\delta 8^{}`$.) Indeed, this assumption leads immediately to an infinite value of the mixing time for a half filled drum: there is no any mixing between the grains of the adjacent wedges if $`\delta =0`$, and the material will never mix.
The finite difference between the angle of marginal stability and the repose angle makes $`\tau `$ to be restricted even if the drum is half filled. Hence, the finiteness of $`\delta `$ determines the avalanche mixing process in this particular case that is considered here.
Let us assume that there are two fractions – black and white, and that the small grains of the different fractions are distinguished from each other only by their color. The material in sector d of the angle $`\delta `$ is mixed homogeneously after the avalanches falls. Then, after the first half of turn, the mixing dynamics can be described by a set of concentration values, $`\{c_m\}`$, where $`c_m`$ is the concentration of the black fraction in sector $`𝐝`$ (Fig. 1) after the $`m`$-th avalanche falls, that happens at time $`t=m\delta `$. Here, time is an angle of drum rotation. Note that a drum radius is not essential for the consideration and does not appear in our expressions. In the case of a half filled drum, the equation for $`c_m`$ turns to be very simple:
$$c_m=\left\{1\left(\frac{\pi }{\delta }\left[\frac{\pi }{\delta }\right]\right)\right\}c_{m[\pi /\delta ]}+\left\{\frac{\pi }{\delta }\left[\frac{\pi }{\delta }\right]\right\}c_{m[\pi /\delta ]1}.$$
(1)
It can be understood easily from Fig. 1. The grains of only two adjacent wedges mix. In each of these wedges the material is homogeneous after the previous avalanches, so two terms in the right part of Eq. (1) are the contributions of the black grains from one of these two wedges, taken with the corresponding weight – their relative angles in sector u in Fig. 1. Here, $`[]`$ denotes an integer part of a number. The values $`c_m,m=[\pi /\delta ],[\pi /\delta ]+1,\mathrm{},1,0`$ may be taken as an initial condition if one does not consider the first half of turn. Obviously, the initial conditions are not essential for the mixing time value.
For brevity, we use the notation $`n[\pi /\delta ]`$ for the integer part and $`a\pi /\delta [\pi /\delta ]`$ for the fractional part. With these notations, Eq. (1) may be written in a more compact form:
$$c_{m+n}=c_ma(c_mc_{m1}).$$
(2)
According to usual prescriptions , one may search for the solution of Eq. (1) or Eq. (2) in the form of a linear combination of the terms $`\lambda _j^m`$, where $`\lambda _j`$ are the roots of the characteristic polynomial
$$\lambda ^{n+1}(1a)\lambda a=0.$$
(3)
To obtain the long time relaxation we should find, among the roots $`\lambda _j`$ of Eq. (3), the maximal value $`|\lambda _j|\lambda _{max}`$ that is lower then the root $`\lambda =1`$ (note that there is no multiple roots; $`\lambda =1`$ corresponds to a stationary situation). The mixing time $`\tau `$ is expressed in the terms of $`\lambda _{max}`$: $`\tau ^1=(1/\delta )\mathrm{log}(1/\lambda _{max})`$, so it may be calculated directly for any value of $`a`$ and a finite $`n`$. The result is shown by the solid line in Fig. 2. One can see that the mixing time is infinite for only integer values of $`\pi /\delta `$. It has local minima near the points $`\pi /\delta =k+1/2`$, where $`k=2,3,\mathrm{}`$.
The process of the mixing may be interpreted in the following way. It follows from Eq. (2) that the long scale and the short one are separated if $`n1`$. Therefore, in such a situation, it is reasonable to use different variables for these two scales. Let us introduce a new time variable $`\stackrel{~}{t}`$. We observe periodically the system under study with the interval of periodicity $`\pi `$ (e.g., using flashes like in the stroboscopic effect) and describe the result of the observation by the quantity $`c(\phi ,\stackrel{~}{t})`$, the concentration of the black fraction at the angle $`\phi `$ at the “time” $`\stackrel{~}{t}`$. Here, $`0<\phi <\pi `$.
At $`n1`$, in the continuous limit, $`c_mc(m\delta )`$. Subtracting $`c((ma)\delta )`$ from both sides of Eq. (2) and using the obvious relation, $`(n+a)\delta =\pi `$, we obtain immediately the equation for $`c(\phi ,\stackrel{~}{t})`$:
$$\frac{c(\phi ,\stackrel{~}{t})}{\stackrel{~}{t}}=\frac{\delta ^2a(1a)}{2\pi }\frac{^2c(\phi ,\stackrel{~}{t})}{\phi ^2}.$$
(4)
Hence, in such an approach, the mixing looks like the diffusion governed by Eq. (4). We do not write out its solution since only the mixing time, i.e., the time of the exponentional relaxation to a homogeneous state, is interesting for us. It follows from the diffusion coefficient and is of the form:
$$\tau ^1=\frac{2}{\pi }\delta ^2\left(\frac{\pi }{\delta }\left[\frac{\pi }{\delta }\right]\right)\left\{1\left(\frac{\pi }{\delta }\left[\frac{\pi }{\delta }\right]\right)\right\}.$$
(5)
This result (which is obtained in the limit $`\delta \pi `$) is shown by the dashed curve in Fig. 2. One may see that the deviations from the results obtained directly from the characteristic polinomial equation are small even for large $`\delta `$. The considerable difference is visible only in the unphysical range, $`\pi /3<\delta <\pi /2`$.
Above, we considered the fixed $`\delta `$. Nevertheless, really, in an experiment, the angle $`\delta `$ is is not fixed but fluctuates with time. A particular form of the distribution function of $`\delta `$ is dependent on characteristics of the mixed grains and is not discussed here. What is the effect of these fluctuations?
Let the difference $`\delta _m`$ between the marginal stability angle and the repose angle for the $`m`$-th avalanche be fluctuating and, hence, depending on $`m`$. Let the distribution function for it be $`P(\delta )`$. One introduces the total angle of the drum turn, $`\theta _m=_{i=1}^m\delta _i`$. Then, after the $`m`$-th avalanche, the wedge between $`\theta _{m1}`$ and $`\theta _m`$ of width $`\delta _m`$ is homogeneously mixed with the concentration of the black fraction $`c_m`$.
The material for the mixing is taken from the wedge $`\theta _{m1}\pi <t<\theta _m\pi `$. Now, unlike the previously considered case without the fluctuations, it may be more than two wedges with different concentration of the black grains in this sector. Therefore, if we assume $`\theta _{m1}\pi <\theta _j<\mathrm{}<\theta _{j+k}<\theta _m\pi `$, where the value of $`k`$ depends on the particular set of $`\{\delta _m\}`$, one obtaines the following equation:
$`c_m=`$ $`{\displaystyle \frac{1}{\delta _m}}\{[\theta _j(\theta _{m1}\pi )]c_j+\delta _{j+1}c_{j+1}+\mathrm{}+`$ (7)
$`\delta _{j+k}c_{j+k}+[(\theta _m\pi )\theta _{j+k}]c_{j+k+1}\}`$
that is the strict generalization of Eq. (1). We study this equation numerically with account of the fluctuations of $`\delta _m`$ to find the mixing time. The results for the homogeneous distributions $`P(\delta )=\mathrm{\Theta }(\delta (1\mathrm{\Delta })\overline{\delta })\mathrm{\Theta }((1+\mathrm{\Delta })\overline{\delta }\delta )/(2\mathrm{\Delta })`$ with $`\mathrm{\Delta }=1/8`$ and $`\mathrm{\Delta }=1/4`$ are presented in Fig. 2 (here, $`\mathrm{\Theta }()`$ is the theta-function). (Note, that, in Fig. 2, we do not connect the points obtained from the numerics.) One sees that the singularities of $`\tau `$ are effectively suppressed by the fluctuations. The oscillations of the dependence decrease with decrease of $`\delta `$, and $`(\pi /(2\delta ^2))\tau ^1`$ approaches a constant value at small $`\delta `$. For small values of the parameter $`\mathrm{\Delta }`$, this value is approached more quickly than for large ones but it is easy to check that, for all $`C>0`$, it equals approximately $`0.175`$. That is close to $`1/6=_0^1𝑑aa(1a)`$, i.e., to the average value of $`(\pi /(2\delta ^2))\tau ^1`$ at small $`\delta `$.
In the experiment , $`\delta =8^{}\pm 2^{}`$. From our estimation, that leads to $`\tau 70`$ turns of the drum. There are no available experimental data yet to be compared with this value. The measuring of $`\tau `$ in such a situation would be really interesting. One can expect that our estimation is quite reasonable because the agreement of the results of the geometrical approach for $`\delta =0`$ with the experiment is excellent.
In conclusion, we have described analytically the dynamics of the avalanche mixing in the case of a half filled slowly rotated drum when the difference $`\delta `$ between the angle of the marginal stability and the angle of repose plays the principal role determining the value of the mixing time. The mixing looks like diffusion of grains between wedges with the diffusion coefficient vanishing when the angle $`\delta `$ is an integer part of $`\pi `$. In these points, the dependence $`\tau (\delta )`$ has singularities which may be effectively suppressed by fluctuations of $`\delta `$. The proposed approach is, in fact, geometrical, and does not contain any other parameters apart of $`\delta `$ or its distribution. Two main assumptions were made: (i) We considered the granular material consisting of very small grains (much smaller than the drum size). (ii) The mixing of the fractions after each elementary avalanche were assumed to be homogeneous. These assumptions reserve space for future study of the avalanche mixing problem.
SND thanks PRAXIS XXI (Portugal) for a research grant PRAXIS XXI/BCC/16418/98. JFFM was partially supported by the projects PRAXIS/2/2.1/FIS/299/94. We also thank E. Lage for reading the manuscript and A.N. Samukhin and H. Watanabe for many useful discussions.
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# RECURSION OPERATORS OF SOME EQUATIONS OF HYDRODYNAMIC TYPE
## I. Introduction
Most of the integrable nonlinear partial differential equations admit Lax representations
$$L_t=[A,L],$$
(1)
where $`L`$ is a pseudo-differential operator of order $`m`$ and $`A`$ is a pseudo differential operator. Recently we established a new method for such integrable equations to construct their recursion operators. This method uses the hierarchy of equations
$$L_{t_n}=[A_n,L]$$
(2)
and the Gel’fand-Dikkii construction of the $`A_n`$-operators. Defining an operator $`R_n`$ in the form
$$A_n=LA_{nm}+R_n,$$
(3)
one then obtains relations among the hierarchies
$$L_{t_n}=LL_{t_{nm}}+[R_n;L]$$
(4)
This equation allows to find $`L_{t_n}`$ in terms of $`L_{t_{nm}}`$. It is important to note that one does not need to know the exact form of $`A_n`$. For further details of the method see .
Here we extend this method to equations of hydrodynamic type . These equations and their Hamiltonian formulation (sometimes called the dispersion-less KdV system) were studied by Dubrovin and Novikov . See for more details on this subject. It is known that these equations admit a nonstandard Lax representation
$$\frac{L}{t}=\{A,L\}_k,$$
(5)
where $`A,L`$ are differentiable functions of $`t,x,p`$ on a Poisson manifold $`M`$ with local coordinates $`(x,p)`$ and $`\{,\}_k`$ is the Poisson bracket. On $`M`$ we take this Poisson bracket $`\{,\}_k=p^k\{,\}`$, where $`\{,\}`$ is the canonical Poisson bracket and k is an integer. For more information on Poisson manifolds see ,. Equations of hydrodynamic type with the above Lax representations were studied in -. Having such Lax representation, we can consider a whole hierarchy of equations
$$\frac{L}{t_n}=\{A_n,L\}_k.$$
(6)
We can also represent function $`A_{n+m}`$ in the form given in (3) and apply our method for constructing of a recursion operator for the equation (6). There are some other works \- which give also recursion operators of some equations of hydrodynamic type. The form of these operators are different then the recursion operators presented in this work. Our method produces recursion operators for hydrodynamic type of equations in the form $`=A+BD^1`$ where $`A`$ and $`B`$ are functions of dynamical variables and their derivatives. All higher symmetries obtained by the repeated application of this recursion operator to translational symmetries belong also to the hydrodynamic type of equations. The recursion operators obtained in references - are of the form $`=CD+A+BD^1E`$, where $`A,B,C`$, and $`E`$ are functions of dynamical variables and their derivatives.
In the next section we discuss the Lax representation with Poisson brackets for polynomial Lax functions. In Sec.III we give the method of construction of the recursion operators following . In Sec.IV we give several examples for $`k=0`$ and $`k=1`$. In Sec.V we consider the Poisson bracket for general $`k`$ and let
$$L=p+S+Pp^1,$$
(7)
and find the Lax equations and the corresponding recursion operator for $`N=2`$. In Sec.VI we consider the Lax function
$$L=p^{\gamma 1}+u+\frac{v^{\gamma 1}}{(\gamma 1)^2}p^{\gamma +1},$$
(8)
and take $`k=0`$. We obtain the equations corresponding to the polytropic gas dynamics and its recursion operators , . It is interesting to note that the systems of equations and their recursion operators obtained in these sections V and VI are transformable into each other. In Sec.VII we give a method reduction from $`N+1`$ system to an $`N`$ system by letting one of the symmetrical variables (defined in the text) to zero. Reduced systems are shown to be also integrable, i.e., they admit recursion operators.
## II. Lax formulation with Poisson bracket
We start with the definition of the standard Poisson bracket. Let $`f(x,p)`$ and $`g(x,p)`$ be differentiable functions of their arguments. Then the standard Poisson bracket is defined by (see and for more details)
$$\{f,g\}=\frac{f}{p}\frac{g}{x}\frac{f}{x}\frac{g}{p}.$$
(9)
We give a slight modification of this bracket as
$$\{f,h\}_k=p^k\{f,g\},$$
(10)
where $`k`$ is an integer. Then we have
Lemma 1. For any $`k`$ the bracket $`\{f,g\}_k`$ defines a Poisson brackets.
Proof. We should check only the Jacobi identity. Other properties of Poisson brackets are evidently true. The standard bracket $`\{f,g\}`$ satisfies the Jacobi identity. For all $`k`$ we have to show that
$$\{\{f,g\}_k,h\}_k+\{\{h,f\}_k,g\}_k+\{\{g,h\}_k,f\}_k=0.$$
First, note that
$$\{\{f,g\}_k,h\}_k=p^k\{p^k\{f,g\},h\}=p^{2k}\{\{f,g\},h\}+kp^{k1}\{f,g\}h_x.$$
Thus, we have
$$\{\{f,g\}_k,h\}_k+\{\{h,f\}_k,g\}_k+\{\{g,h\}_k,f\}_k=p^{2k}(\{\{f,g\},h\}+\{\{h,f\},g\}+$$
$$\{\{g,h\},f\})+kp^{k1}(\{f,g\}h_x+\{h,f\}g_x+\{g,h\}f_x).$$
Equality
$$\{\{f,g\},h\}+\{\{h,f\},g\}+\{\{g,h\},f\}=0$$
holds and it is easy to check that
$$\{f,g\}h_x+\{h,f\}g_x+\{g,h\}f_x=0.$$
Hence, $`\{,\}_k`$, for all $`k`$ defines a Poisson bracket. $`\mathrm{}`$
For each $`k𝐙`$ we can consider hierarchies of equations of hydrodynamic type, defined in terms of the Lax function
$$L=p^{N1}+\underset{i=1}{\overset{N2}{}}p^iS_i(x,t)$$
(11)
by the Lax equation
$$\frac{L}{t_n}=\{(L^{\frac{n}{N1}})_{k+1};L\}_k,$$
(12)
where $`n=j+l(N1)`$ and $`j=1,2,\mathrm{},(N1),l𝐍`$ . So we have a hierarchy for each $`k`$ and $`j=1,\mathrm{},(N1)`$. Also, we require $`nk+1`$ to ensure that $`(L^{\frac{n}{N1}})_{k+1}`$ is not zero. With the choice of Poisson brackets $`\{,\}_k`$, we must take a certain part of the series expansion of $`L^{\frac{n}{N1}}`$ to get the consistent equation (12). This part is $`(L^{\frac{n}{N1}})_{k+1}`$.
The Lax function (11) can also be written in terms of symmetric variables $`u_1,\mathrm{},u_N`$
$$L=\frac{1}{p}\underset{j=1}{\overset{N}{}}(pu_j)$$
(13)
that is $`u_1,\mathrm{},u_N`$ are roots of the polynomial
$$p^{N1}+S_{N2}p^{N2}+\mathrm{}+S_1.$$
In new variables the equation (12) is invariant under transposition of variables.
## III. Recursion Operators
For each hierarchy of the equations (12), depending on the pair $`(N,k)`$, we can find a recursion operator.
Lemma 2. For any $`n`$
$$L_n=LL_{n(N1)}+\{R_n;L\}_k,$$
(14)
where function $`R_n`$ has a form
$$R_n=\underset{i=0}{\overset{N2}{}}p^{ik}A_i(S_1\mathrm{}S_{N2},S_{1,n1}\mathrm{}S_{N2,n1}).$$
(15)
Proof.
$$(L^{\frac{n}{N1}})_{k+1}=[L(L^{\frac{n}{N1}1})_{k+1}+L(L^{\frac{n}{N1}1})_{<k+1}]_{k+1}$$
So,
$`(L^{\frac{n}{N1}})_{k+1}=L(L^{\frac{n}{N1}1})_{k+1}+(L(L^{\frac{n}{N1}1})_{<k+1})_{k+1}`$
$`(L(L^{\frac{n}{N1}1})_{k+1})_{<k+1}.`$ (16)
If we put
$$R_n=(L(L^{\frac{n}{N1}1})_{<k+1})_{k+1}(L(L^{\frac{n}{N1}1})_{k+1})_{<k+1},$$
then
$$(L^{\frac{n}{N1}})_{k+1}=L(L^{\frac{n}{N1}1})_{k+1}+R_n.$$
Hence,
$`L_n=\{(L^{\frac{n}{N1}})_{k+1};L\}_k=\{L(L^{\frac{n}{N1}1})_{k+1}+R_n;L\}_k`$
$`=LL_{n(N1)}+\{R_n;L\}_k,`$ (17)
and (14) is satisfied. Evaluating powers of $`(L(L^{\frac{n}{N1}1})_{<k+1})_{k+1}`$ and $`(L(L^{\frac{n}{N1}1})_{k+1})_{<k+1}`$ we get that $`R_n`$ has form (15). $`\mathrm{}`$
Lemma 3. A recursion operator for the hierarchy (12) is given by equalities, for $`m=N2,N3,\mathrm{},1`$,
$$\begin{array}{c}S_{m,n+(N1)}=_{j=1}^{m+1}S_jS_{mj,n}+_{j=1}^{m+1}(j+1k)A_{j+1}S_{mj,x}\\ _{j=1}^{m+1}(mj)A_{j+1,x}S_{mj},\end{array}$$
(18)
where to simplify the above formula we have defined that $`S_{N1}=1`$ and $`S_{N1,x}=0`$, $`S_{N1,n}=0`$. Coefficients $`A_{N2},A_{N3},\mathrm{},A_0`$ can be found from the recursion relations, for $`m=N2,\mathrm{},1`$
$`(N1)A_{m,x}={\displaystyle \underset{j=m}{\overset{N1}{}}}S_jS_{N2+mj}+{\displaystyle \underset{j=m}{\overset{N2}{}}}(j+1k)A_{j+1}S_{N2+mj,x}`$
$`{\displaystyle \underset{j=m}{\overset{N2}{}}}(N2+mj)A_{j+1,x}S_{N2+mj}`$ (19)
Proof. Let us write the equality (14), using (15) for $`R_n`$
$$\underset{i=1}{\overset{N2}{}}p^iS_{i,n+(N1)}=\left(p^{N1}+\underset{i=1}{\overset{N2}{}}p^iS_i\right)\left(\underset{i=1}{\overset{N2}{}}p^iS_{i,n}\right)+$$
$$p^k\left(\underset{j=0}{\overset{N1}{}}(jk)p^{jk1}A_j\right)\left(\underset{j=1}{\overset{N2}{}}p^jS_{j,x}\right)$$
$$p^k\left(\underset{j=0}{\overset{N1}{}}p^{jk}A_{j,x}\right)\left((N1)p^{N2}+\underset{j=1}{\overset{N2}{}}jp^{j1}S_j\right)$$
To have the equality, the coefficients of $`p^{2N3},\mathrm{},p^{N1}`$ and $`p^2`$ must be zero, it gives recursion relations to find $`A_{N2},\mathrm{},A_0`$. The coefficients of $`p^{N2},\mathrm{},p^1`$ give the expressions for $`S_{N2,n+(N1)},\mathrm{},S_{1,n+(N1)}`$. $`\mathrm{}`$
Although the recursion operator $``$, given by (18), is a pseudo-differential operator, but it gives a hierarchy of local symmetries starting from the equation itself. Indeed, equalities (18), (III. Recursion Operators) give expressions $`S_{N2,n+(N1)}`$ , $`\mathrm{}`$, $`S_{1,n+(N1)}`$ in terms of $`S_{N2}`$ ,$`\mathrm{}`$, $`S_1`$ and $`S_{N2,n}`$, $`\mathrm{}`$ , $`S_{1,n}`$. Hence, the recursion operator $``$ is constructed in such a way that
$$\{(L^{\frac{n}{N1}+1})_{k+1};L\}_k=\left(\{(L^{\frac{n}{N1}})_{k+1};L\}_k\right)$$
(20)
## IV. Some Integrable Systems
We shall consider first some examples for $`k=0`$, $`k=1`$ and the general case in the next section.
### A. Multicomponent hierarchy containing also the shallow water wave equations , $`k=0`$
This hierarchy corresponds to the case $`k=0`$. Let us give the first equation of hierarchy and a recursion operator for $`N=2,3`$.
Proposition 1. In the case $`N=2`$ one has the Lax function
$$L=p+S+Pp^1$$
and the Lax equation for $`n=2`$, given by (47), when $`k=0`$
$$\begin{array}{ccc}\frac{1}{2}S_t& =& SS_x+P_x,\\ \frac{1}{2}P_t& =& SP_x+PS_x.\end{array}$$
(21)
and the recursion operator, given by (48),
$$=\left(\begin{array}{cc}S+S_xD_x^1& 2\\ 2P+P_xD_x^1& S\end{array}\right).$$
(22)
These equations are known as the shallow water wave equations or as the equations of polytropic gas dynamics for $`\gamma =2`$ (See Sec.VI).
The first two symmetries of the system (21) are given by
$`S_{t_1}`$ $`=`$ $`(S^3+6SP)_x,`$
$`P_{t_1}`$ $`=`$ $`(3S^2P+3P^2)_x,`$ (23)
$`S_{t_2}`$ $`=`$ $`(S^4+12S^2P+6P^2)_x,`$
$`P_{t_2}`$ $`=`$ $`(4S^3P+12SP^2)_x.`$ (24)
These are all commuting symmetries.
Remark 1. In symmetric variables the system (21) is written as
$$\begin{array}{ccc}\frac{1}{2}u_t& =& (u+v)u_x+uv_x,\\ \frac{1}{2}v_t& =& vu_x+(u+v)v_x,\end{array}$$
(25)
and the recursion operator (22) takes the form
$$=\left(\begin{array}{cc}u+v+u_xD_x^1& 2u+u_xD_x^1\\ 2v+v_xD_x^1& u+v+v_xD_x^1\end{array}\right).$$
(26)
Proposition 2.In the case $`N=3`$ one has the Lax function
$$L=p^2+pS+P+p^1Q$$
and the Lax equation with $`n=3`$ is
$$\begin{array}{ccc}\frac{1}{3}S_t& =& (\frac{1}{2}P\frac{1}{8}S^2)S_x+\frac{1}{2}SP_x+Q_x,\\ \frac{1}{3}P_t& =& \frac{1}{2}QS_x+(\frac{1}{8}S^2+\frac{1}{2}P)P_x+SQ_x,\\ \frac{1}{3}Q_t& =& \frac{1}{4}SQS_x+\frac{1}{2}QP_x+(\frac{1}{8}S^2+\frac{1}{2}P)Q_x.\end{array}$$
(27)
The recursion operator, corresponding to this equation, is
$$=\left(\begin{array}{ccc}\frac{S^2}{4}+P+P_xD_x^1\frac{S_x}{4}D_x^1S& \frac{S}{2}+\frac{S_x}{2}D_x^1& 3\\ \frac{3Q}{2}+(Q_x+\frac{P_xS}{2})D_x^1\frac{P_x}{4}D_x^1S& P+\frac{P_x}{2}D_x^1& 2S\\ \frac{SQ}{4}+(\frac{SQ_x}{2}+\frac{S_xQ}{2})D_x^1\frac{Q_x}{4}D_x^1S& \frac{3Q}{2}+\frac{Q_x}{2}D_x^1& P\end{array}\right).$$
(28)
Proof. Using (III. Recursion Operators) we find the function $`R_n`$ and using (18) we find the recursion operator (28). $`\mathrm{}`$
Remark 2. In symmetric variables the equation (27) is written as
$$\begin{array}{ccc}\frac{1}{3}u_t& =& (\frac{1}{8}u^2+\frac{1}{2}(uv+uw+vw)+\frac{1}{8}(v+w)^2)u_x\\ & & +(\frac{1}{4}u^2+\frac{1}{4}uv+\frac{3}{4}uw)v_x+(\frac{1}{4}u^2+\frac{1}{4}uw+\frac{3}{4}uv)w_x,\\ \frac{1}{3}v_t& =& (\frac{1}{4}v^2+\frac{1}{4}uv+\frac{3}{4}vw)u_x+(\frac{1}{4}v^2+\frac{1}{4}vw+\frac{3}{4}uv)w_x\\ & & +(\frac{1}{8}v^2+\frac{1}{2}(uv+uw+vw)+\frac{1}{8}(u+w)^2)v_x,\\ \frac{1}{3}w_t& =& (\frac{1}{4}w^2+\frac{1}{4}uw+\frac{3}{4}wv)u_x+(\frac{1}{4}w^2+\frac{1}{4}wv+\frac{3}{4}uw)v_x\\ & & +(\frac{1}{8}w^2+\frac{1}{2}(uv+uw+vw)+\frac{1}{8}(v+u)^2)w_x,\end{array}$$
(29)
and the recursion operator takes the form (67) given in the Appendix.
### B. Toda hierarchy ($`k=1`$)
Toda hierarchy corresponds to the case $`k=1`$. Let us give the first equation of hierarchy and a recursion operator for $`N=2`$ and $`N=3`$.
Proposition 3. In the case $`N=2`$ and $`n=1`$ one has the Lax function
$$L=p+S+Pp^1$$
and the Lax equation for $`n=1`$ , given by (41),
$$\begin{array}{ccc}S_t& =& P_x,\\ P_t& =& PS_x,\end{array}$$
(30)
and the recursion operator, given by (42),
$$=\left(\begin{array}{cc}S& 2+P_xD_x^1P^1\\ 2P& S+S_xPD_x^1P^1\end{array}\right).$$
(31)
The first two symmetries of the equation (30) are given by
$`S_{t_1}`$ $`=`$ $`(2SP)_x,`$
$`P_{t_1}`$ $`=`$ $`P(2P+S^2)_x,`$ (32)
$`S_{t_2}`$ $`=`$ $`(3S^2P+3P^2)_x,`$
$`P_{t_2}`$ $`=`$ $`P(6PS+S^3)_x.`$ (33)
Remark 3. In symmetric variables the equation (30) is written as
$$\begin{array}{ccc}u_t& =& uv_x,\\ v_t& =& vu_x,\end{array}$$
(34)
and the recursion operator (31) takes the form
$$=\left(\begin{array}{cc}u+v+uv_xD_x^1u^1& 2u+uv_xD_x^1v^1\\ 2v+vu_xD_x^1u^1& u+v+vu_xD_x^1v^1\end{array}\right).$$
(35)
Proposition 4. In the case $`N=3`$ and $`n=1`$ one has the Lax function
$$L=p^2+pS_1+P+p^1Q$$
and the Lax equation with $`n=1`$ is
$$\begin{array}{ccc}S_t& =& P_x\frac{1}{2}SS_x,\\ P_t& =& Q_x,\\ Q_t& =& \frac{1}{2}QS_x.\end{array}$$
(36)
The recursion operator, corresponding to this equation, is
$$=\left(\begin{array}{ccc}P\frac{1}{4}S^2+(\frac{1}{2}P_x\frac{1}{4}SS_x)D_x^1& \frac{1}{2}S& 3+2Q_xD_x^1Q^1\\ \frac{3}{2}Q+\frac{1}{2}Q_xD_x^1& P& 2S+(SQ)_xD_x^1Q^1\\ \frac{1}{4}SQ+\frac{1}{4}S_xQD_x^1& \frac{3}{2}Q& P+P_xQD_x^1Q^1\end{array}\right).$$
(37)
Proof. Using equalities (III. Recursion Operators) we find the function $`R_n`$ and using (18) we find the recursion operator (37). $`\mathrm{}`$
Remark 4. In symmetric variables the equation (36) is written as
$$\begin{array}{ccc}u_t& =& \frac{1}{2}u(u_x+v_x+w_x),\\ v_t& =& \frac{1}{2}v(+u_xv_x+w_x),\\ w_t& =& \frac{1}{2}w(+u_x+v_xw_x),\end{array}$$
(38)
and the recursion operator takes the form (68) given in the Appendix.
## V. Lax equation for general $`k`$
We shall only consider the case where $`N=2`$. We have the Lax function
$$L=p+S+Pp^1$$
(39)
and the Lax equation
$$\frac{L}{t_n}=\{(L^n)_{k+1};L\}_k.$$
(40)
We consider two cases $`k1`$ and $`k0`$.
### A. The first case $`k1`$
Proposition 5. In the case $`N=2`$ and $`k1`$ one has the Lax equation
$$\begin{array}{ccc}S_t& =& kP^{k1}P_x,\\ P_t& =& kP^kS_x.\end{array}$$
(41)
and the recursion operator for this equation is
$$=\left(\begin{array}{cc}S+(1k)S_xD_x^1& 2+kP^{k1}P_xD_x^1P^k\\ 2P+(1k)P_xD_x^1& S+kS_xP^kD_x^1P^k\end{array}\right).$$
(42)
Proof. The smallest power of $`p`$ in $`L^n`$ is $`n`$. To have powers less than $`k+1`$ we must put $`n=k`$. If there are no such powers then Poisson brackets are $`\{(L^n);L\}_k=0`$.
Let us calculate the Lax equation
$$L_t=\{(L^k)_{k+1};L\}_k=\{(L^k)_k;L\}_k$$
We have $`(L^k)_k=[(p+S+Pp^1)^k]_k=P^kp^k`$, thus
$$L_t=\{P^kp^k;p+S+Pp^1\}_k.$$
And we get the equation (41). Using (18), (III. Recursion Operators) we find the recursion operator (42). $`\mathrm{}`$
First two symmetries are given as follows
$`S_{t_1}`$ $`=`$ $`(P^kS)_x,`$
$`P_{t_1}`$ $`=`$ $`P^k\left(P+{\displaystyle \frac{k}{2}}S^2\right)_x,`$ (43)
$`S_{t_2}`$ $`=`$ $`(k+1)(k+2)\left({\displaystyle \frac{1}{2}}P^kS^2+{\displaystyle \frac{1}{k+1}}P^{k+1}\right)_x,`$
$`P_{t_2}`$ $`=`$ $`(k+1)(k+2)P^k\left(PS+{\displaystyle \frac{k}{6}}S^3\right)_x.`$ (44)
Remark 5. In symmetric variables the equation (41) is written as
$$\begin{array}{ccc}u_t& =& ku^kv^{k1}v_x,\\ v_t& =& ku^{k1}v^ku_x,\end{array}$$
(45)
and the recursion operator (42) takes the form
$$=\left(\begin{array}{cc}u+v+(1k)u_xD_x^1+& 2u+(1k)u_xD_x^1+\\ ku^kv^{k1}v_xD_x^1u^kv^{k+1}& ku^kv^{k1}v_xD_x^1u^{k+1}v^k\\ & \\ 2v+(1k)v_xD_x^1+& u+v+(1k)v_xD_x^1+\\ ku^{k1}v^ku_xD_x^1u^kv^{k+1}& ku^{k1}v^ku_xD_x^1u^{k+1}v^k\end{array}\right).$$
(46)
### B. The second case $`k0`$
Proposition 6. In the case $`N=2`$ and $`k0`$ one has the Lax equation
$$\begin{array}{ccc}S_t& =& (k+2)(k+1)SS_x+(k+2)P_x,\\ P_t& =& (k+2)(k+1)SP_x+(k+2)S_xP.\end{array}$$
(47)
and the recursion operator for this equation is
$$=\left(\begin{array}{cc}S+(1k)S_xD_x^1& 2+kP^{k1}P_xD_x^1P^k\\ 2P+(1k)P_xD_x^1& S+kS_xP^kD_x^1P^k\end{array}\right).$$
(48)
Proof. The largest power of $`p`$ in $`L^n`$ is $`p^n`$. To have powers larger than $`k+1`$ we must put $`n=k+1`$. Then we have
$$(L^{k+1})_{k+1}=[(p+S+Pp^1)^{k+1}]_{k+1}=p^{k+1},$$
thus
$$L_t=\{p^{k+1};p+S+Pp^1\}_k.$$
Then the Lax equation becomes
$`S_t`$ $`=`$ $`S_x,`$
$`P_t`$ $`=`$ $`P_x.`$
This is a trivial equation, let us calculate the second symmetry. We have $`(L^{k+2})_{k+1}=[(p+S+Pp^1)^{k+1}]_{k+1}=p^{k+2}+(k+2)Sp^{k+1}`$, thus
$$L_t=\{p^{k+2}+(k+2)Sp^{k+1};p+S+Pp^1\}_k.$$
we get the equation (47). Using (18), (III. Recursion Operators) we find the recursion operator (48). $`\mathrm{}`$
First two symmetries are given as follows
$`S_{t_1}`$ $`=`$ $`(k2)(k3)\left(PS+{\displaystyle \frac{1}{6}}(1k)S^3\right)_x,`$
$`P_{t_1}`$ $`=`$ $`(k2)(k3)\left(SS_xP+{\displaystyle \frac{1}{2}}(1k)S^2P_x+PP_x\right),`$ (49)
$`S_{t_2}`$ $`=`$ $`(2k)(3k)(4k)\left({\displaystyle \frac{1}{2}}S^2P+{\displaystyle \frac{1}{6}}S^4+{\displaystyle \frac{1}{2(2k)}}P^2\right)_x,`$
$`P_{t_2}`$ $`=`$ $`(2k)(3k)(4k)({\displaystyle \frac{1}{2}}S^2S_xP+{\displaystyle \frac{1}{6}}(1k)S^3P_x`$ (50)
$`+SPP_x+{\displaystyle \frac{1}{(2k)}}P^2S_x).`$
Remark 6. In symmetric variables the equation (47) is written as
$$\begin{array}{ccc}u_t& =& (k+2)(1k)(u+v)u_x+(k+2)uv_x,\\ v_t& =& (k+2)vu_x+(k+2)(1k)(u+v)v_x,\end{array}$$
(51)
and the recursion operator (48) takes the form
$$=\left(\begin{array}{cc}u+v+(1k)u_xD_x^1+& 2u+(1k)u_xD_x^1+\\ ku^kv^{k1}v_xD_x^1u^kv^{k+1}& ku^kv^{k1}v_xD_x^1u^{k+1}v^k\\ & \\ 2v+(1k)v_xD_x^1+& u+v+(1k)v_xD_x^1+\\ ku^{k1}v^ku_xD_x^1u^kv^{k+1}& ku^{k1}v^ku_xD_x^1u^{k+1}v^k\end{array}\right).$$
(52)
In this Section, to obtain the recursion operators we have considered two different cases $`k0`$ and $`k1`$ to simplify some technical problems in the method. At the end we obtained recursion operators having the same forms (42) and (48). Hence any one of these represent the recursion operator for $`k`$. It seems , comparing the results, that the systems of equations in one case are symmetries of the other case. For instance, the system (47) is a symmetry of system (41). Hence we may consider only one case with recursion operator $`(\text{42})`$ for all integer values of $`k`$.
## VI. Lax function for polytropic gas dynamics
In this section we consider another Lax function, introduced in ,
$$L=p^{\gamma 1}+u+\frac{v^{\gamma 1}}{(\gamma 1)^2}p^{\gamma +1}$$
(53)
and the Lax equation
$$\frac{L}{t}=\frac{\gamma 1}{\gamma }\{(L^{\frac{\gamma }{\gamma 1}})_1,L\}_0,$$
(54)
gives the equations of the polytropic gas dynamics.
Proposition 7. The Lax equation corresponding to (54) is
$$\begin{array}{ccc}u_t+uu_x+v^{\gamma 2}v_x& =& 0,\\ v_t+(uv)_x& =& 0.\end{array}$$
(55)
Proof. Expanding the function (53) around the point $`p=\mathrm{}`$, we have
$$\left(p^{\gamma 1}+u+\frac{v^{\gamma 1}}{(\gamma 1)^2}p^{\gamma +1}\right)^{\frac{\gamma }{\gamma 1}}=p^\gamma +\frac{\gamma }{\gamma 1}pu+\mathrm{}$$
all other terms have negative powers of $`p`$. Therefore
$$\left(L^{\frac{\gamma }{\gamma 1}}\right)_1=p^\gamma +\frac{\gamma }{\gamma 1}pu$$
and the Lax equation (54) corresponds to (55). $`\mathrm{}`$
Proposition 8. The recursion operator for the equation (55) is
$$=\left(\begin{array}{cc}u+\frac{u_x}{\gamma 1}D_x^1& \frac{2v^{\gamma 2}}{\gamma 1}+\frac{(v^{\gamma 2})_x}{\gamma 1}D_x^1\\ \frac{2v}{\gamma 1}+\frac{v_x}{\gamma 1}D_x^1& u+\frac{\gamma 2}{\gamma 1}u_xD_x^1\end{array}\right).$$
(56)
Proof. Using the equation
$$\frac{L}{t_{n+1}}=L\frac{L}{t_n}+\{R_n,L\}.$$
in the same way as for polynomial Lax function one can find the recursion operator (56). $`\mathrm{}`$
It is interesting to note that the equation (47) and equations of polytropic gas dynamics (55) are related by the following change of variables
$$\begin{array}{cc}S=& \frac{u}{(k+2)(k+1)},\\ & \\ P=& \frac{v^{\frac{1}{k+1}}}{(k+2)^2},\end{array}$$
(57)
where $`\gamma ={\displaystyle \frac{k+2}{k+1}}`$. We note that under this change of variables recursion operator (48) is mapped to the recursion operator (56).
## VII. Reduction
We consider reduction of the equation (12), written in symmetric variables, by setting $`u_1=0`$. Let us write the equation (12) as
$$\mathrm{\Delta }(u_N,\mathrm{},u_1)=0,$$
(58)
where $`\mathrm{\Delta }`$ is a differential operator. Then
$$\mathrm{\Delta }(u_N,\mathrm{},u_1)|_{u_1=0}=\left(\begin{array}{c}\stackrel{~}{\mathrm{\Delta }}(u_N,\mathrm{},u_2),\\ 0\end{array}\right)$$
(59)
where $`\stackrel{~}{\mathrm{\Delta }}`$ is another differential operator. Indeed, following for the Lax function $`L=\frac{1}{p}_{j=1}^N(pu_j)`$ we have
$$\frac{L}{t}=L\underset{j=1}{\overset{N}{}}\frac{u_{j,t}}{p+u_j},$$
$$\frac{L}{x}=L\underset{j=1}{\overset{N}{}}\frac{u_{j,x}}{p+u_j}$$
and
$$\frac{L}{p}=L\left(\frac{1}{p}+\underset{j=1}{\overset{N}{}}\frac{1}{p+u_j}\right).$$
Thus $`u_{j,t}=Res_{p=u_j}\{M,L\}_k`$, where $`M=(L^{\frac{n}{N1}})_{k+1}`$. The Lax equation (12) can be written as
$$\underset{j=1}{\overset{N}{}}\frac{u_{j,t}}{p+u_j}=p^kM_p\underset{j=1}{\overset{N}{}}\frac{u_{j,x}}{p+u_j}p^kM_x\left(\frac{1}{p}+\underset{j=1}{\overset{N}{}}\frac{1}{p+u_j}\right).$$
Note, that $`p^kM_x`$ and $`p^kM_p`$ are polynomials. So, if we put $`u_1=0`$ and calculate residue of right hand side at $`p=0`$ we get (59). A new equation
$$\stackrel{~}{\mathrm{\Delta }}(u_N,\mathrm{},u_2)=0$$
(60)
is also integrable and a recursion operator of this equation can be obtained as reduction of the recursion operator of the equation (58). Let $``$ be the recursion operator of (58) given by lemma 3 , then
$$|_{u_1=0}=\left(\begin{array}{cc}\stackrel{~}{R}& \\ 0\mathrm{}0& 0\end{array}\right).$$
(61)
Indeed, we found the recursion operator using formula (14). This formula can be written as
$$\underset{j=1}{\overset{N}{}}\frac{u_{j,t_n}}{p+u_j}=LL_{n(N1)}+p^kR_{n,p}\underset{j=1}{\overset{N}{}}\frac{u_{j,x}}{p+u_j}p^kR_{n,x}\left(\frac{1}{p}+\underset{j=1}{\overset{N}{}}\frac{1}{p+u_j}\right)$$
and in the same way as for reduction of (58) we have (61), note, that $`p^kR_{n,x}`$ and $`p^kR_{n,p}`$ are also polynomials.
Lemma 4. The operator $`\stackrel{~}{R}`$ is a recursion operator of the equation (60).
Proof. Equation (60) is an evolution equation, so, to prove that $`\stackrel{~}{R}`$ is a recursion operator we must prove that for any solution $`(u_N,\mathrm{},u_2)`$ of (60) the following equality holds (see )
$$D_{\stackrel{~}{\mathrm{\Delta }}}\stackrel{~}{R}=\stackrel{~}{R}D_{\stackrel{~}{\mathrm{\Delta }}},$$
where $`D_{\stackrel{~}{\mathrm{\Delta }}}`$ is a Frechet derivative of $`\stackrel{~}{\mathrm{\Delta }}`$.
If $`(u_N,\mathrm{},u_2)`$ is a solution of (60) then $`(u_N,\mathrm{},u_2,u_1=0)`$ is a solution of (58) and for the solution $`(u_N,\mathrm{},u_2,u_1=0)`$ we have
$$D_\mathrm{\Delta }=D_\mathrm{\Delta }.$$
(62)
Next
$$D_\mathrm{\Delta }|_{u_1=0}=\left(\begin{array}{cc}\stackrel{~}{D}& \\ 0\mathrm{}0& \end{array}\right)$$
and
$$|_{u_1=0}=\left(\begin{array}{cc}\stackrel{~}{R}& \\ 0\mathrm{}0& 0\end{array}\right).$$
Hence by (62) we have $`\stackrel{~}{D}\stackrel{~}{R}=\stackrel{~}{R}\stackrel{~}{D}`$. Calculating Frechet derivative we take derivatives with respect to one variable, considering other variables as constants. Thus, to calculate $`\stackrel{~}{D}`$ we can put $`u_1=0`$ and differentiate with respect to other variables or we can first differentiate and then put $`u_1=0`$. It means that $`\stackrel{~}{D}=D_{\stackrel{~}{\mathrm{\Delta }}}`$ and
$$D_{\stackrel{~}{\mathrm{\Delta }}}\stackrel{~}{R}=\stackrel{~}{R}D_{\stackrel{~}{\mathrm{\Delta }}}.$$
$`\mathrm{}`$
Let us consider reduction of systems, given by remark 2 and remark 4 and their recursion operators.
Proposition 9. Putting $`w=0`$ in (38) and (68) we obtain a new system
$$\begin{array}{ccc}u_t& =& \frac{1}{2}u(u_x+v_x),\\ v_t& =& \frac{1}{2}v(+u_xv_x),\end{array}$$
(63)
and its recursion operator
$$=\left(\begin{array}{cc}uv+\frac{u}{4}(u+v)& \frac{u}{4}(u+v)\\ +\frac{u}{4}(u_xv_x)D_x^1& +\frac{u}{4}(u_xv_x)D_x^1\\ & \\ \frac{v}{4}(u+v)& uv+\frac{v}{4}(u+v)\\ +\frac{v}{4}(u_x+v_x)D_x^1& +\frac{v}{4}(u_x+v_x)D_x^1\end{array}\right)$$
(64)
respectively. $`\mathrm{}`$
Proposition 10. Putting $`w=0`$ in (29) and (67) we obtain a new system
$$\begin{array}{ccc}\frac{1}{3}u_t& =& (\frac{1}{8}u^2+\frac{1}{2}uv+\frac{1}{8}v^2)u_x+(\frac{1}{4}u^2+\frac{1}{4}uv)v_x,\\ \frac{1}{3}v_t& =& (\frac{1}{4}v^2+\frac{1}{4}uv)u_x+(\frac{1}{8}v^2+\frac{1}{2}uv+\frac{1}{8}u^2)v_x,\end{array}$$
(65)
and its recursion operator
$$=\left(\begin{array}{cc}\frac{u^2}{4}+\frac{3uv}{4}+(\frac{u_xv}{2}+\frac{uv_x}{2})D_x^1& \frac{u}{4}(u+v)+(\frac{u_xv}{2}+\frac{uv_x}{2})D_x^1\\ \frac{u_x}{4}D_x^1u+\frac{u_x}{4}D_x^1v& +\frac{u_x}{4}D_x^1u\frac{u_x}{4}D_x^1v\\ & \\ \frac{v}{4}(u+v)+(\frac{uv_x}{2}+\frac{u_xv}{2})D_x^1& \frac{v^2}{4}+\frac{3uv}{4}+(\frac{uv_x}{2}+\frac{u_xv}{2})D_x^1\\ \frac{v_x}{4}D_x^1u+\frac{v_x}{4}D_x^1v& +\frac{v_x}{4}D_x^1u\frac{v_x}{4}D_x^1v\end{array}\right),$$
(66)
respectively. $`\mathrm{}`$
It is worth to mention that by reduction we obtain a new equation. For example, consider the case $`k=0`$. The equation (25, corresponding to $`N=2`$, and reduction of the equation (29), corresponding to $`N=3`$, are not related by a linear transformation of variables. Indeed, in the equation (25)) coefficients of $`u_x,v_x`$ are linear in $`u,v`$ but in the equation (65) coefficients of $`u_x,v_x`$ contain quadratic terms. Hence they can not be related by a linear transformation.
## VII. Conclusion
We have constructed the recursion operators of some equations of hydrodynamic type. The form of the these operators fall into the class of pseudo differential operators $`A+BD^1`$ where $`A`$ and $`B`$ are functions of dynamical variables and their derivatives. The generalized symmetries of these equations are local and all belong to the same class (i.e., they are also equations of hydrodynamic type). We have introduced a method of reduction which leads also to integrable class. These properties, bi Hamiltonian structure of the equations we obtained and equations with rational Lax functions will be communicated elsewhere.
## Acknowledgments
We thank Burak Gürel and Atalay Karasu for several discussions. This work is partially supported by the Scientific and Technical Research Council of Turkey and by Turkish Academy of Sciences.
Appendix. Recursion operators of the systems (29) and (38) are respectively given by
$$=\left(\begin{array}{ccc}\frac{u^2}{4}+\frac{3}{4}(uv+uw)+wv& \frac{u}{4}(u+v+w)+\frac{3uw}{2}& \frac{u}{4}(u+v+w)+\frac{3uv}{2}\\ +\frac{u_x}{2}(v+w)D_x^1& +\frac{u_x}{2}(v+w)D_x^1& +\frac{u_x}{2}(v+w)D_x^1\\ +\frac{u}{2}(v_x+w_x)D_x^1& +\frac{u}{2}(v_x+w_x)D_x^1& +\frac{u}{2}(v_x+w_x)D_x^1\\ \frac{u_x}{4}D_x^1u+\frac{u_x}{4}D_x^1v& +\frac{u_x}{4}D_x^1u\frac{u_x}{4}D_x^1v& +\frac{u_x}{4}D_x^1u+\frac{u_x}{4}D_x^1v\\ +\frac{u_x}{4}D_x^1w& +\frac{u_x}{4}D_x^1w& \frac{u_x}{4}D_x^1w\\ & & \\ \frac{v}{4}(u+v+w)+\frac{3vw}{2}& \frac{v^2}{4}+\frac{3}{4}(uv+vw)+uw& \frac{v}{4}(u+v+w)+\frac{3uv}{2}\\ +\frac{v_x}{2}(u+w)D_x^1& +\frac{v_x}{2}(u+w)D_x^1& +\frac{v_x}{2}(u+w)D_x^1\\ +\frac{v}{2}(u_x+w_x)D_x^1& +\frac{v}{2}(u_x+w_x)D_x^1& +\frac{v}{2}(u_x+w_x)D_x^1\\ \frac{v_x}{4}D_x^1u+\frac{v_x}{4}D_x^1v& +\frac{v_x}{4}D_x^1u\frac{v_x}{4}D_x^1v& +\frac{v_x}{4}D_x^1u+\frac{v_x}{4}D_x^1v\\ +\frac{v_x}{4}D_x^1w& +\frac{v_x}{4}D_x^1w& \frac{v_x}{4}D_x^1w\\ & & \\ \frac{w}{4}(u+v+w)+\frac{3vw}{2}& \frac{w}{4}(u+v+w)+\frac{3uw}{2}& \frac{w^2}{4}+\frac{3}{4}(uw+vw)+uv\\ +\frac{w_x}{2}(u+v)D_x^1& +\frac{w_x}{2}(u+v)D_x^1& +\frac{w_x}{2}(u+v)D_x^1\\ +\frac{w}{2}(u_x+v_x)D_x^1& +\frac{w}{2}(u_x+v_x)D_x^1& +\frac{w}{2}(u_x+v_x)D_x^1\\ \frac{w_x}{4}D_x^1u+\frac{w_x}{4}D_x^1v& +\frac{w_x}{4}D_x^1u\frac{w_x}{4}D_x^1v& +\frac{w_x}{4}D_x^1u+\frac{w_x}{4}D_x^1v\\ +\frac{w_x}{4}D_x^1w& +\frac{w_x}{4}D_x^1w& \frac{w_x}{4}D_x^1w\end{array}\right).$$
(67)
$$=\left(\begin{array}{ccc}(uv+uw+vw)& \frac{u}{4}(u+v+w)& \frac{u}{4}(u+v+w)\\ +\frac{u}{4}(u+v+w)& \frac{3uw}{2}& \frac{3uv}{2}\\ +\frac{u}{4}(u_xv_xw_x)D_x^1& +\frac{u}{4}(u_xv_xw_x)D_x^1& +\frac{u}{4}(u_xv_xw_x)D_x^1\\ u(wv_x+vw_x)D_x^1u^1& u(wv_x+vw_x)D_x^1v^1& u(wv_x+vw_x)D_x^1w^1\\ & & \\ \frac{v}{4}(u+v+w)& (uv+uw+vw)& \frac{v}{4}(u+v+w)\\ \frac{3vw}{2}& +\frac{v}{4}(u+v+w)& \frac{3uv}{2}\\ +\frac{v}{4}(u_x+v_xw_x)D_x^1& +\frac{v}{4}(u_x+v_xw_x)D_x^1& +\frac{v}{4}(u_x+v_xw_x)D_x^1\\ v(wu_x+uw_x)D_x^1u^1& v(wu_x+uw_x)D_x^1v^1& v(wu_x+uw_x)D_x^1w^1\\ & & \\ \frac{w}{4}(u+v+w)& \frac{w}{4}(u+v+w)& (uv+uw+vw)\\ \frac{3uw}{2}& \frac{3vw}{2}& +\frac{w}{4}(u+v+w)\\ +\frac{w}{4}(u_xv_x+w_x)D_x^1& +\frac{w}{4}(u_xv_x+w_x)D_x^1& +\frac{w}{4}(u_xv_x+w_x)D_x^1\\ w(uv_x+vu_x)D_x^1u^1& w(uv_x+vu_x)D_x^1v^1& w(uv_x+vu_x)D_x^1w^1\end{array}\right).$$
(68)
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# Neutrino Masses at 𝑣^{3/2}
## I Introduction
Neutrinos are exceedingly light compared to other fermionic elementary particles. For instance, the atmospheric neutrino data suggest $`m_\nu 0.06`$ eV. To understand the smallness of neutrino masses, we usually invoke the presence of a scale $`M`$ that is much larger than $`v200`$ GeV, the scale of electroweak symmetry breaking. This scale is conventionally taken to be roughly the reduced Planck mass, $`M_{Pl}2\times 10^{18}\mathrm{GeV}`$, or the GUT scale, $`M_{GUT}2\times 10^{16}\mathrm{GeV}`$. Setting $`M=1`$, $`v`$ becomes a small parameter of the theory. Charged fermion masses arise at order $`v`$, while neutrino masses arise at $`v^2`$. Note that any particle of mass $`v^3`$ is massless for all particle physics practicalities, although such masses could be interesting for cosmology.
The see-saw mechanism is certainly the simplest explanation of why neutrino masses are so small - however, the result $`v^2/M`$ gives neutrino masses which are $`303000`$ times too light to explain the atmospheric neutrino anomaly, for $`M=M_{GUT}`$ or $`M_{Pl}`$. Frequently this is fixed by instead taking $`M10^{14}10^{15}\mathrm{GeV}`$. In this paper we propose an alternative approach.
In supersymmetric theories, a standard framework has supersymmetry broken in a hidden sector at the intermediate scale $`m_I10^{10}10^{11}\mathrm{GeV}`$ which allows supergravity mediation of supersymmetry breaking to the standard model to generate the weak scale $`vm_I^2/M_{Pl}`$. We call this the intermediate-scale see-saw. An important question is at what order neutrino masses arise within this framework. If the relevant operator is $`(LH)^2/M_{Pl}`$, then $`m_\nu v^2/M_{Pl}`$ as before. However, with $`M=M_{Pl}=1`$, the small dimensionless parameter of the expansion is now $`m_I`$ rather than $`v`$. There are now more possibilities, namely $`m_\nu `$ could occur at order $`m_I^3,m_I^4,m_I^5\mathrm{}`$ corresponding to $`v^{3/2},v^2,v^{5/2}\mathrm{}`$ We propose that scale for atmospheric neutrino oscillations arises at $`O(m_I^3)`$. The scale of solar oscillations can also arise at $`O(m_I^3)`$ with an extra suppression due to approximate flavor symmetries, or it can instead occur at $`O(m_I^4)`$. Alternatively, the hierarchy between the atmospheric and solar scales can be traced to an approximately rank one loop integral, as discussed in section III.
## II Neutrino masses at $`m_I^3`$
How might neutrino masses arise at $`O(m_I^3)`$? People have noted before that in a framework with large extra dimensions such that $`M_{Pl}`$ is lowered to $`m_I`$, the ordinary see-saw result yields $`m_\nu v^2/m_I`$. Such an approach is quite distinct from our perspective, where $`m_I`$ is the small parameter. Other authors have explored the possibility of generating neutrino masses from higher dimension operators involving a field with a vacuum expectation value (vev) at an intermediate scale, such as $`10^{11}\mathrm{GeV}`$. However, in these models, a supersymmetry-conserving vev generates the small couplings and masses, and the size of the vev is essentially a free parameter of the theory. In this paper, we will instead utilize the direct connection to supersymmetry breaking explored in . Adopting this approach can have significant phenomenological ramifications, as we will find for the model of section III.
Right handed states, singlet under the standard model, might be light if they are protected by some global symmetry $`G`$, analogous to a symmetry used to prevent a Planck-scale $`\mu H_uH_d`$ term in the superpotential. In , it was noted that if $`G`$ is broken in the supersymmetry breaking sector, then it is quite natural to have light neutrinos. In the models presented in , the neutrino masses arose at $`O(m_I^4)`$.
We can employ the same framework to generate $`O(m_I^3)`$ masses instead. Consider the superpotential
$$\left[XNN+LNH\right]_F$$
(1)
where $`N`$ is a standard model singlet and $`X`$ is a field in the supersymmetry breaking sector that takes on an $`A`$ component vev $`X=m_I`$ (for instance, the superpotential
$$\left[S(X\overline{X}m_I^2)+\overline{X}^2Y+X^2\overline{Y}\right]_F$$
(2)
generates $`F_SF_Y=F_{\overline{Y}}m_I^2`$ and $`A_X=A_{\overline{X}}m_I`$, but $`F_X=F_{\overline{X}}=0`$). The Lagrangian then contains
$$\left[m_INN\right]_F+\left[LNH\right]_F,$$
(3)
and when the Higgs takes on a vev, we have small Dirac masses for the neutrinos in addition to the Majorana mass for the right handed neutrino. Setting $`M_{Pl}=1`$, the neutrino mass matrix is
$$m_{LR}=\left(\begin{array}{cc}0& m_I^2\\ m_I^2& m_I\end{array}\right),$$
(4)
leading to a see-saw mass $`m_\nu m_I^3`$ for the left handed neutrino. If we further assume that the Dirac masses are suppressed by $`\lambda _\tau `$, for instance as might occur with flavor symmetries, the neutrino mass is $`\lambda _\tau ^2m_I^30.1\mathrm{eV}`$. (A similar estimate $`\lambda _\mu ^2m_I^3`$ gives an approximately correct mass scale for the LOW solution to the solar neutrino problem if there are more than one $`N`$).
This model is exceedingly simple, but it illustrates the fact that once we take the small parameter of the theory to be $`m_I`$ rather than $`v`$, there is no a priori reason why we should expect neutrino masses to occur at $`m_I^4`$.
The phenomenology of this model is identical to that of ordinary see-saw neutrino physics, except that we do not expect the signals that might accompany broken GUT symmetries at the scale $`10^{14}\mathrm{GeV}`$. In contrast, the model that we consider next features additional weak-scale states and predicts a much richer phenomenology.
## III Radiative neutrino masses at $`m_I^3`$
In , a model was proposed in which the right handed neutrinos get masses at $`O(m_I^2)`$, rather than at $`O(m_I)`$. If the Yukawa couplings are order one, then all neutrino masses are weak scale. As discussed in , this is remedied rather simply.
Suppose that a hidden sector field $`X`$ acquires $`A`$ and $`F`$ component vevs of $`O(m_I)`$ and $`O(m_I^2)`$, respectively. Consider the operators
$$\frac{1}{M_{Pl}}\left(\left[XLNH_u\right]_F+\left[X^{}NN\right]_D\right),$$
(5)
which yield
$$\frac{m_I^2}{M_{Pl}}[NN]_F+\frac{m_I}{M_{Pl}}[LNH]_F+\frac{m_I^2}{M_{Pl}}[LNH]_A$$
(6)
when $`X`$ acquires its vevs. We take there to be just one $`N`$ superfield. The first term in (6) generates a weak scale mass for $`N`$, while the second term generates Yukawas of order $`m_I/M_{Pl}`$. The neutrino mass matrix which we generate at tree level is (again setting $`M_{Pl}=1`$)
$$m_{\mathrm{tree}}=\left(\begin{array}{cc}0& m_I^3\\ m_I^3& m_I^2\end{array}\right).$$
(7)
Although this is not the canonical see-saw, as it involves a Dirac mass at $`O(v^{3/2})`$, it nevertheless yields the light neutrino mass of $`v^2/M_{Pl}`$ as usual. The unusual feature is that the right handed neutrino is at the weak scale. The operators of (5) are easily justified, for instance by assuming ordinary R parity (under which $`N`$ is odd), together with an R symmetry where $`N`$ has R charge $`2/3`$, $`X`$ has R charge $`4/3`$ and $`L`$ and $`H_u`$ have R charge $`0`$ .
Once we have included the operator $`\left[X^{}NN\right]_D`$, it is impossible to forbid $`\left[X^{}XX^{}NN\right]_D`$. This operator induces a lepton number violating scalar mass of the form $`\delta ^2\stackrel{~}{n}\stackrel{~}{n}+h.c.`$, where $`\delta m_I^{5/2}/M_{Pl}^{3/2}`$. The importance of this operator is that it allows a radiative contribution to the light neutrino mass at $`O(m_I^3)`$ from the diagram of figure 1.
If we take $`m_{\stackrel{~}{\nu }_L}m_{\stackrel{~}{n}}m_{\chi ^0}`$, and call them generically $`\stackrel{~}{m}`$, we generate a neutrino mass
$$m_\nu \frac{g^2}{384\pi ^2}\frac{A^2v^2\delta ^2}{\stackrel{~}{m}^5}.$$
(8)
Taking $`A\stackrel{~}{m}v`$, this becomes
$$m_\nu \frac{g^2}{384\pi ^2}\left(\frac{m_I^3}{M_{Pl}^2}\right).$$
(9)
Note that this is larger than the lighter eigenvalue of (7) since it occurs at $`O(m_I^3)`$ rather than $`O(m_I^4)`$. A tree level neutrino mass at $`O(m_I^3)`$ is $`O(\mathrm{keV})`$ \- too heavy to be interesting. However, our mechanism automatically leads to a loop factor of order $`10^4`$ giving masses of $`O(.1\mathrm{eV}1\mathrm{e}\mathrm{V})`$ \- within one order of magnitude of the scale necessary to explain the atmospheric neutrino anomaly!
Although the $`A`$ terms couple $`\stackrel{~}{n}`$ to only a single linear combination of $`\stackrel{~}{\nu }`$’s, the loop diagram of figure 1 can generate more than one neutrino mass eigenvalue. For incoming $`\nu _i`$ and $`\nu _j`$, the value of the loop integral, $`L_{ij}`$ has nontrivial dependence on the corresponding sneutrino masses $`\stackrel{~}{m}_i`$ and $`\stackrel{~}{m}_j`$. The resulting mass matrix, $`m_{ij}A_iA_jL_{ij}`$ is not necessarily rank one, and we can expect a second and third eigenvalue. Although $`L_{ij}`$ does depend on $`i`$ and $`j`$, to a large extent it factors into $`f_if_j`$, and is approximately rank one. Consequently, the second eigenvalue is suppressed greatly compared to first. We have investigated this numerically, and for a broad range of the parameters the second eigenvalue is typically a factor of $`10^2`$ or smaller down from the first. Consequently, it may either be that the mass scale to explain the solar neutrino anomaly arises from this additional suppression, or that it arises at $`O(m_I^4)`$ from (7). Of course, one could alternatively use more than one $`N`$ and a hierarchical $`A`$ matrix to generate a hierarchical $`m_{ij}`$.
## IV Neutrino Mass Anarchy
The possibility has been explored elsewhere that neutrino mass matrices have no ordering structure, such as a flavor symmetry . Absence of flavor symmetry is even more reasonable in our framework - all suppressions arise naturally via loop factors or factorization, or occur at different orders in $`m_I`$.
If the parameters in the model display no apparent structure, that is, the $`A_i`$ are all roughly equal, and likewise the $`m_{\stackrel{~}{\nu }_i}`$ are roughly - but not exactly - the same, then we have a natural justification for the large mixing observed between $`\nu _\mu `$ and $`\nu _\tau `$. We would then expect the solution to the solar neutrino problem to similarly involve a large angle: either large angle MSW, vacuum oscillations or the LOW solution. The only small parameter required is $`\theta _{13}<0.16`$, required by CHOOZ, but we can view this as an accident rather than a fine tuning.
Even if we have a flavor symmetry which explains the structure of the charged fermion masses, that would not necessarily preclude such a scenario. If the structure of the $`A`$-terms were determined by a supersymmetry-preserving spurion $`\lambda _i`$, i.e.,
$$A_i\stackrel{~}{n}\stackrel{~}{\nu }^ih_u=\frac{\left[X\right]_F}{M_{Pl}}\lambda _i\stackrel{~}{n}\stackrel{~}{\nu }^ih_u,$$
(10)
then we expect a hierarchy in $`A_i`$ related to that which we find in the charged leptons. However, if $`X`$ carries a flavor index itself, i.e.,
$$A_i\stackrel{~}{n}\stackrel{~}{\nu }^ih_u=\frac{\left[X_i\right]_F}{M_{Pl}}\stackrel{~}{n}\stackrel{~}{\nu }^ih_u,$$
(11)
then the situation is quite different. Because the structure of $`X_i`$ is determined in the supersymmetry breaking sector, it need not be related to the structure of the lepton masses. In this case, even with a flavor symmetry, we would expect large angles to arise in the neutrino sector.
## V Phenomenology and cosmology
The phenomenology of the model presented in section III is very interesting. It is essentially identical to that explored in for the case of a single $`N`$ superfield. The presence of a weak-scale $`\stackrel{~}{n}`$ that mixes through weak-scale $`A`$ terms to the left-handed sneutrinos can profoundly affect the sneutrino spectrum. For instance, a sneutrino mass eigenstate is not subject to the $`Z`$-width constraint if it is mostly composed of $`\stackrel{~}{n}`$, and its mass can be different from that of $`\stackrel{~}{l}_L`$ by far more than just the $`D`$-term splitting. The $`A`$ terms could potentially induce invisible Higgs decays into light sneutrino pairs, and $`\stackrel{~}{\nu }\stackrel{~}{l}`$ might be the dominant decay mode for the charged Higgs boson. Finally, cascade decays producing heavy sneutrinos that subsequently decay into a Higgs and light sneutrino could conceivably be the dominant source of Higgs production at the LHC. There are numerous other potential consequences, arising in particular scenarios, which we will not discuss here.
The presence of the additional $`\stackrel{~}{n}`$ state revives the possibility of sneutrino dark matter, as was explored in in the context of both lepton number conserving and lepton number violating models. In the lepton number conserving case, direct detection experiments require $`m_{\stackrel{~}{\nu }}<3`$ GeV. In contrast, the lepton number violating scalar mass term central to the present model allows for easy evasion of this bound.
Direct detection experiments detect ordinary sneutrino dark matter through $`Z`$ boson exchange. However, once a lepton number violating mass term is present, the CP-even state $`\stackrel{~}{\nu }_+`$ and CP-odd state $`\stackrel{~}{\nu }_{}`$ are no longer degenerate in mass. Moreover, scalar couplings to the $`Z`$ are off-diagonal, i.e., they couple $`\stackrel{~}{\nu }_+`$ to $`\stackrel{~}{\nu }_{}`$, but not $`\stackrel{~}{\nu }_{}`$ to $`\stackrel{~}{\nu }_{}`$. The scattering of $`\stackrel{~}{\nu }_{}`$ off of a nucleus through $`Z`$ exchange is kinematically forbidden for $`\mathrm{\Delta }m>\beta _h^2m_{}m_A/2(m_A+m_{})`$, where $`\mathrm{\Delta }m=m_+m_{}`$ is the mass splitting between the CP-even and CP-odd states, $`m_A`$ is the mass of the nucleus, and $`\beta _h=10^3`$ for virialized halo particles on average. Thus, even for sneutrino masses of $`O(100`$ GeV) , direct detection limits are essentially harmless, stipulating only that the lepton number violating mass is adequately large. For example, taking $`m_{\stackrel{~}{\nu }}=100`$ GeV and $`m_A=72`$ GeV for a Ge target, we simply need $`\mathrm{\Delta }m>20`$ keV to prevent direct detection. Because $`\mathrm{\Delta }m=\delta ^2/m_{\stackrel{~}{\nu }}`$, this corresponds to $`\delta >45\mathrm{M}\mathrm{e}\mathrm{V}`$, which is of the order of what we expect from $`m_I^{5/2}/M_{Pl}^{3/2}`$.
Sneutrino dark matter will still scatter from the nuclei via Higgs exchange. The cross section per nucleon for this is small, however, given by
$$\sigma 1.8\times 10^{43}\mathrm{sin}^22\theta $$
(12)
$$\times \left(\frac{A}{100\mathrm{G}\mathrm{e}\mathrm{V}}\right)^2\left(\frac{130\mathrm{G}\mathrm{e}\mathrm{V}}{m_h}\right)^4\left(\frac{100\mathrm{G}\mathrm{e}\mathrm{V}}{m_{\stackrel{~}{\nu }}}\right)^2\mathrm{cm}^2,$$
(13)
where
$$\stackrel{~}{\nu }_{DM}=\stackrel{~}{\nu }\mathrm{sin}\theta +\stackrel{~}{n}\mathrm{cos}\theta .$$
(14)
The current upper bound on $`\sigma `$ for a dark matter candidate with mass $`100\mathrm{G}\mathrm{e}\mathrm{V}`$ is $`10^{41}\mathrm{cm}^2`$, and this bound is expected to be lowered by orders of magnitude in the near future. Thus, this version of sneutrino dark matter may be detectable due to Higgs exchange in upcoming direct detection experiments.
Sneutrino dark matter with lepton number violation has been previously explored, precisely because it can evade direct detection limits. However a large mass splitting ($`\mathrm{\Delta }m_{\stackrel{~}{\nu }}\mathrm{GeV}`$) was necessary to suppress $`\stackrel{~}{\nu }_+\stackrel{~}{\nu }_{}`$ coannihilation in the early universe to yield an appreciable relic density. Here, because the lightest sneutrino is an admixture of left- and right-handed states, the overall annihilation rate via MSSM processes is suppressed by an additional factor of $`\mathrm{sin}^4\theta `$. As discussed in , acceptable relic abundances are obtained for a broad range of parameters. For example, for sneutrino masses less than $`M_W`$, the relic density is essentially determined by the annihilation rate via neutralino exchange, and one finds
$$\mathrm{\Omega }h^2\left(\frac{M_{\stackrel{~}{W}}}{100\mathrm{G}\mathrm{e}\mathrm{V}}\right)^2\left(\frac{\mathrm{sin}\theta }{0.16}\right)^4,$$
(15)
where $`h`$ is the reduced Hubble parameter and $`M_{\stackrel{~}{W}}`$ is the neutral wino mass. Such a simple approximate formula does not apply when the sneutrino mass is heavy enough so that production of $`W`$ and $`Z`$ pairs becomes relevant, but the abundance is still promising for reasonable parameter choices.
## VI Conclusions
While conventional see-saw models generate neutrino masses proportional to $`v^2`$, in theories with gravity mediated supersymmetry breaking, it is also possible to generate neutrino masses proportional to $`v^{3/2}`$. Models of this type arise when flavor symmetries protect the masses of standard model singlet states, but are broken in the supersymmetry breaking sector of the theory.
Additional suppression to these masses can arise from Yukawa-type suppressions, or from loop factors, resulting in values for the neutrino mass in accordance with observations from Superkamiokande of an up/down neutrino asymmetry.
The model developed in section III, which features a right-handed sneutrino at the weak scale, is phenomenologically rich, with dramatic changes to collider signatures and the possibility of sneutrino dark matter. Although the sneutrinos in this model evade current detection limits on dark matter, the possibility exists for their detection at a future experiment.
Acknowledgements
We thank R. Rattazzi for pointing out an error in a previous version of this work. This work was supported in part by the U.S. Department of Energy under Contracts DE-AC03-76SF00098, in part by the National Science Foundation under grant PHY-95-14797.
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# NMR properties of a one-dimensional Cu-O model
## I Introduction
One of the challenging characteristics of the cuprate materials is the importance of magnetic fluctuations and their effect on normal state transport and on superconductivity. This feature is naturally present in theoretical approaches emphasizing strong interactions in 2D. The role of magnetism can also be assessed in other scenarios promoting marginal or nearly antiferromagnetic Fermi liquid (NAFL) behavior. The NAFL framework has been used by Millis, Monien, and Pines to discuss NMR experiments: on the basis of the Mila-Rice and Shastry local terms describing spin fluctuations induced by the hyperfine interactions, these authors were able to compute various Knight shifts and nuclear relaxation times. However, there are critical remarks and further details to this theory, and also other theories introduced in order to interpret the behavior of the various Knight shifts and nuclear relaxation times in high-$`\mathrm{T}_\mathrm{c}`$ superconductors, which we do not discuss here.
In this paper, we will follow the basic assumption of Mila and Rice that it is necessary to include a sizeable isotropic hyperfine interaction term to fit the data of NMR experiments in High-$`\mathrm{T}_\mathrm{c}`$ cuprates. Thus in the following we will focus our interest on this contribution which comes mainly from the Fermi contact interaction between a nucleus and its surrounding partially filled s orbitals, namely the 4s orbital for copper and the 3s orbital for oxygen. Experiments have shown that the magnetic properties are described by a single-spin component model. In the weak interaction limit, this single-spin degree of freedom could be associated with the strongly hybridized Cu–3d—O–2p anti-bonding band, whereas in the strong interaction limit it is associated with the nearly localized Cu-3d spin. Within the local picture, the contrasting NMR behavior seen on the Cu and O sites arises from their different hyperfine form factors. These are nothing but the Fourier transforms of the Fermi contact interaction terms approximated by a sum over surrounding localized next-neighbor Cu–3d spins.
One of the (many) complications concerning the physics of cuprate materials is that there is still no consensus about what should be the correct theoretical approach to treat correlations in 2D: Is the ground state Fermi liquid or non-Fermi liquid like? Can one treat interactions perturbatively, or is it more appropriate to treat kinetic terms as corrections in the strongly interacting limit? In the latter category, working out a consistent treatment of the non-double occupancy constraint is still an open issue.
By contrast, one-dimensional systems offer a perfect testing ground for the study of magnetic fluctuations. Since in one dimension it is possible to treat correlation effects properly both in the limit of weak and strong interactions, such models allow to compute explicitly the dependence of the relaxation. This allows to get some feeling for the effects of doping. In addition to the insight that such study allows to gain for higher dimensional models, there are explicit realizations of one-dimensional systems, such as the Bechgaard salts or copper germanate compounds.
For these reasons, we choose to investigate hyperfine interactions in the one-dimensional version of the Cu-O model. This allows us to extract form factors both in the insulating and in the doped regime, without assuming a specific form of the Fermi contact interaction term. We can then compare the exact results with the predictions that the standard approximation schemes used in two dimensions would give in the one-dimensional situation.
The paper is organized as follows. In section II, we introduce the model in one dimension, as well as the three different approximations used to describe the magnetic relaxation processes related to the Fermi contact interaction term , namely that due to Mila-Rice, to Shastry, and to Bulut. In section III, we solve the full model for weak interactions as compared to the bandwidth. We obtain spin-spin correlation functions and discuss the Knight shift $`K`$ and the relaxation rate $`1/T_1`$ in detail. We compare our results with the prediction of the Bulut model, which is applicable for weak interactions. Section IV solves the problem in the opposite limit of very strong interactions using a Gutzwiller projection eliminating double occupation on the copper sites. We again compare our results with the one-dimensional extrapolation of the Mila-Rice and of the Shastry approximations. A general discussion of our results is presented in section V. Since NMR data on organic and inorganic quasi-1D compounds seem to give an essentially isotropic relaxation rate, the body of the paper mostly focuses on the isotropic contribution to the hyperfine interaction. Yet, for the sake of completeness and in view of the fact that $`K`$ can be anisotropic (see below in section V) we discuss the effect of anisotropic hyperfine terms in Appendix A: these terms only modify prefactors in the expressions of $`1/T_1`$ and of $`K`$. Appendix B and C offer details of our calculations.
## II Definition of the models
### A The four-band model
We consider a system with two different atoms per unit cell (Cu and O). In order to describe the ground state properties, we take into account the 3d and 2p orbitals on Cu and O respectively. The related hole states are represented in Fig. 1 and denoted by $`a`$ and $`b`$ in the following. Since the coupling to the nuclear spin via the Fermi contact interaction occurs only for partially filled s orbitals, we also have to retain the Cu–4s and O–3s shells (denoted by $`A`$ and $`B`$) to correctly obtain the desired NMR properties. The Hamiltonian describing the system can thus be written as
$$H=H_0+H_S+H_N,$$
(1)
where $`H_0`$ contains the electronically relevant orbitals $`a`$ and $`b`$. $`H_S`$ describes the coupling of orbitals $`a`$ and $`b`$ to orbitals $`A`$ and $`B`$. As indicated in Fig. 1, the orbitals $`A`$ and $`B`$ are basically filled and produce only small corrections to the electronic term represented by $`H_0`$, so that we will treat $`H_S`$ as a perturbation. Finally, $`H_N`$ describes the coupling of the orbitals $`A`$ and $`B`$ to the nuclear spins and will be treated as a small – time dependent – perturbation in linear response.
The main contribution, $`H_0`$, is given by
$$H_0=H_T+H_U,$$
(2)
where
$`H_T`$ $`=`$ $`{\displaystyle \underset{j}{}}ϵ_an_{aj}+ϵ_bn_{bj}`$ (3)
$``$ $`{\displaystyle \underset{j\sigma }{}}t_{ab}\left[a_{j\sigma }^{}(b_{j\sigma }+b_{j1,\sigma })+\text{h.c.}\right]`$ (4)
$`H_U`$ $`=`$ $`{\displaystyle \underset{j}{}}U_an_{aj}n_{aj}+{\displaystyle \underset{j}{}}U_bn_{bj}n_{bj}.`$ (5)
Here $`t_{ab}`$ describes the hopping between the Cu-3d and the O-2p orbitals with the phase conventions shown in Fig. 2. $`U_a`$ and $`U_b`$ are the local repulsions on the copper and oxygen sites as shown in Fig. 1. The Coulomb repulsions $`U_A`$ and $`U_B`$ can be ignored assuming that energy cost considerations discourage processes in which two holes are excited in one s orbital. A nearest neighbor interaction $`U_{ab}`$ could also be added to the model to generate a phase transition to a superconducting phase, but here our study deals with the vicinity of the half-filled case, i.e. near the antiferromagnetic phase, and we will ignore $`U_{ab}`$.
The coupling between the orbitals $`a,b`$ and $`A,B`$ reads
$`H_S`$ $`=`$ $`{\displaystyle \underset{j}{}}ϵ_An_{Aj}+ϵ_Bn_{Bj}`$ (6)
$`+`$ $`{\displaystyle \underset{j\sigma }{}}t_{Ba}\left[B_{j\sigma }^{}(a_{j+1,\sigma }a_{j\sigma })+\text{h.c.}\right]`$ (7)
$`+`$ $`{\displaystyle \underset{j\sigma }{}}t_{Ab}\left[A_{j\sigma }^{}(b_{j\sigma }+b_{j1,\sigma })+\text{h.c.}\right]`$ (8)
$`+`$ $`{\displaystyle \underset{j\sigma }{}}t_{Aa}\left[A_{j\sigma }^{}(a_{j+1,\sigma }+a_{j1,\sigma })+\text{h.c.}\right]`$ (9)
with the phase conventions of Fig. 2. The density operators $`n_{\eta j}`$ in (3), (5), and (6) are the standard ones
$$n_{\eta j}=\underset{\sigma }{}n_{\eta j\sigma }=\underset{\sigma }{}\eta _{j\sigma }^{}\eta _{j\sigma },$$
(10)
where $`\eta =a,b,A,B`$.
Finally, the isotropic coupling to the nuclear spins ($`𝐈`$ for the copper atom and $`𝐉`$ for the oxygen atom) is given by
$$H_N=\underset{j}{}C_A𝐈_j𝐒_{Aj}+C_B𝐉_j𝐒_{Bj},$$
(11)
where similarly to (10) the spin operators are given by
$$𝐒_{\eta j}=\frac{1}{2}\underset{\sigma _1\sigma _2}{}\eta _{j\sigma _1}^{}\sigma \sigma \sigma \sigma _{\sigma _1\sigma _2}\eta _{j\sigma _2}.$$
(12)
The coupling constant for the electron-nuclear interaction is given by $`C_\eta =(8\pi /3)|\psi _\eta (0)|^2\gamma _ng\mu _B\mathrm{}`$, which is proportional to the local hole density for the respective s orbitals at the origin. $`\sigma _{\sigma _1\sigma _2}`$ are the Pauli matrices.
### B Reduced models
In 1D, the Hamiltonian (1) can explicitly be expressed in terms of bose fields. This allows for a full treatment, in which $`H_0`$, $`H_S`$, and $`H_N`$ are all treated on an equal footing. This is the route we follow in section III. In higher dimension there is still no solution of the fully interacting problem. So various approximation schemes have been devised and applied to each of the pieces of $`H`$ separately. These lead to effective ”nuclear” Hamiltonians : the Mila-Rice and the Shastry model for strong interactions and the Bulut model for weak interactions. Typically, the analysis by Mila-Rice and Shastry starts from a partially projected Hamiltonian
$$H=\widehat{P}H_0\widehat{P}+H_S+H_N.$$
(13)
The first part gives the t-J model or at half-filling the Heisenberg model, which contains the dynamics related to the Cu-3d and O-2p orbitals. The second part contains the unprojected degrees of freedom related to the Cu-4s and O-3s orbitals, as well as the electron-nuclear interaction part. Further approximations for $`H_S`$ and $`H_N`$ lead to the Mila-Rice model or to the Shastry model (see below). In the weak interaction limit, Bulut et al. have proposed an RPA treatment of $`H_0`$ in combination with the effective electron-nuclear interaction term of Mila-Rice and Shastry.
Before we turn to the full solution of the model, let us review the main features of such approximations when applied to our one-dimensional system. This will allow us to contrast the predictions of the 1D version of these three models and the results obtained for the four-band model, which may provide some clue to the validity of these approaches for strongly correlated systems.
#### 1 The Mila-Rice model for strong interaction
The model defined in (1) is approximated by
$$HH_0^{Mi}+H_N^{Mi}.$$
(14)
$`H_0^{Mi}`$ is the approximation for (2) and denotes, at half-filling, a Heisenberg model for local Cu-3d spins generated by
$$H_0^{Mi}=\widehat{P}\left[H_0(U_b=0)\right]\widehat{P},$$
(15)
where $`\widehat{P}`$ is the Gutzwiller projection operator which prohibits doubly occupied Cu-3d states. The additional unprojected part (6) with $`t_{Aa}=t_{Ba}=0`$ and the electron-nuclear interaction part (11) are approximated by
$$H_N^{Mi}=\underset{j}{}C_A𝐈_j𝐒_{Aj}^{Mi},$$
(16)
where
$$𝐒_{Aj}^{Mi}=F_{Aa}^{Mi}\left(𝐒_{a,j1}+𝐒_{a,j+1}\right).$$
(17)
Thus $`𝐒_{Aj}^{Mi}`$ is the Mila-Rice approximation for the original spin $`𝐒_{Aj}`$ used to explain the NMR experiments measured on the copper sites. $`F_{Aa}^{Mi}=|\lambda _{Aa}^{Mi}|^2`$ denotes the effective overlap between one Cu-4s spin with a neighboring Cu-3d spin. In the Mila-Rice model only hopping processes via the O-2p orbitals are included, whereas the direct hopping between Cu-3d and Cu-4s orbitals is ignored. In Ref. Mila and Rice perform a quantum chemical analysis without including interaction effects, so we will do the same and propose for the amplitudes
$$\lambda _{Aa}^{Mi}=\frac{t_{Aa}t_{ab}}{(ϵ_aϵ_A)(ϵ_aϵ_b)}.$$
(18)
This result is obtained by a projection in real space of a Cu-4s orbital onto a neighboring Cu-3d orbital for $`U_a=0`$. The Fourier transform of (17) is given by
$$𝐒_{Ap}^{Mi}=F_{Aa}^{Mi}(p)𝐒_{ap}$$
(19)
with the Mila-Rice form factor
$$F_{Aa}^{Mi}(p)=2F_{Aa}^{Mi}\mathrm{cos}(pa).$$
(20)
#### 2 The Shastry model for strong interaction
The approximation proposed by Shastry in Ref. is given by
$$HH_0^{Sh}+H_N^{Sh},$$
(21)
where
$$H_0^{Sh}=\widehat{P}\left[H_0(U_b=0,U_a=\mathrm{})\right]\widehat{P}$$
(22)
leads to the Heisenberg model at half-filling and to the t-J model for a doped system with strong repulsion on the copper sites. The electron-nuclear interaction part reads
$$H_N^{Sh}=\underset{j}{}C_A𝐈_j𝐒_{Aj}^{Sh}+C_B𝐉_j𝐒_{Bj}^{Sh},$$
(23)
where the spins $`𝐒_{\eta j}^{Sh}`$ are approximated by a linear combination of unprojected Cu-3d orbitals
$`𝐒_{Aj}^{Sh}`$ $`=`$ $`F_{Aa}^{Sh}\left(𝐒_{a,j1}+𝐒_{a,j+1}\right)`$ (24)
$`+`$ $`F_{Ab}^{Sh}\left(𝐒_{bj}+𝐒_{b,j1}\right)`$ (25)
$``$ $`F_{Aa}^{Sh}\left(𝐒_{a,j1}+𝐒_{a,j+1}\right)`$ (26)
and
$$𝐒_{Bj}^{Sh}=F_{Ba}^{Sh}\left(𝐒_{aj}+𝐒_{a,j+1}\right)$$
(27)
with the coefficients
$`F_{Aa}^{Sh}=|\lambda _{Aa}^{Sh}|^2`$ $`=`$ $`\left({\displaystyle \frac{t_{Aa}}{ϵ_Aϵ_a}}\right)^2`$ (28)
$`F_{Ab}^{Sh}=|\lambda _{Ab}^{Sh}|^2`$ $`=`$ $`\left({\displaystyle \frac{t_{Ab}}{ϵ_Aϵ_b}}\right)^2`$ (29)
$`F_{Ba}^{Sh}=|\lambda _{Ba}^{Sh}|^2`$ $`=`$ $`\left({\displaystyle \frac{t_{Ba}}{ϵ_Bϵ_a}}\right)^2.`$ (30)
For finite doping, it is assumed in Ref. that the spin degrees of freedom related to $`𝐒_{bj}`$ are quenched in a Zhang-Rice singlet. This assumption justifies the second approximation done in (26). Further, Shastry includes only the direct couplings (28) up to second order proportional to $`t_{Aa}^2`$ for the relaxation of the nuclear copper spin and ignores the fourth order contributions proportional to $`t_{ab}^2t_{Ab}^2`$ via the O-2p orbital as proposed by Mila and Rice. The Fourier transform of the approximated spins $`𝐒_{\eta j}^{Sh}`$ reads
$$𝐒_{\eta p}^{Sh}=2F_{\eta a}^{Sh}(p)𝐒_{ap}$$
(31)
with
$`F_{Aa}^{Sh}(p)`$ $`=`$ $`2F_{Aa}^{Sh}\mathrm{cos}(pa)`$ (32)
$`F_{Ba}^{Sh}(p)`$ $`=`$ $`2F_{Ba}^{Sh}\mathrm{cos}(pa/2).`$ (33)
For the uniform contribution ($`p0`$) all form factors are finite, but for the antiferromagnetic wave vector ($`p\pi /a`$) the form factor vanishes for the oxygen sites, whereas it stays finite for the copper sites. Thus (20), (32), and (33) are the one-dimensional analogs of the NMR Mila-Rice and Shastry form factors for High-$`\mathrm{T}_\mathrm{c}`$ cuprates.
#### 3 The Bulut model for weak interaction
In Ref. , Bulut et al. used a weak interaction RPA calculation combined with an electron-nuclear interaction as proposed by Mila-Rice and Shastry to compute NMR related quantities for the High-$`\mathrm{T}_\mathrm{c}`$ cuprates. The one-dimensional analog reads
$$HH_0^{Bu}+H_N^{Bu}.$$
(34)
In two dimensions, $`H_0^{Bu}`$ is obtained by applying the RPA method to the 2D version of (2) with $`U_b=0`$. In one dimension, we can treat all the interaction terms ($`U_a,U_b`$) of the original Hamiltonian $`H_0`$ by means of bosonization and of renormalization-group theory. Finally, the electron-nuclear interaction term $`H_N^{Bu}`$ is given in analogy to (23) replacing the approximated spins by
$`𝐒_{Aj}^{Bu}`$ $`=`$ $`F_{Aa}^{Bu}\left(𝐒_{a,j1}+𝐒_{a,j+1}\right)`$ (35)
$`𝐒_{Bj}^{Bu}`$ $`=`$ $`F_{Ba}^{Bu}\left(𝐒_{aj}+𝐒_{a,j+1}\right).`$ (36)
In Ref. , the parameters $`F_{\eta a}^{Bu}`$ are undefined and could in general include all possible overlaps of the Cu-4s and O-3s orbitals with the Cu-3d orbitals in the sense of Mila-Rice and Shastry. The coefficient for the oxygen will be
$$F_{Ba}^{Bu}=|\lambda _{Ba}^{Bu}|^2=|\lambda _{Ba}^{Sh}|^2,$$
(37)
and is thus the same as that proposed by Shastry, whereas the coefficient for copper
$$F_{Aa}^{Bu}=|\lambda _{Aa}^{Bu}|^2=|\lambda _{Aa}^{Sh}+\lambda _{Aa}^{Mi}|^2$$
(38)
includes additional combined terms of third order proportional to $`t_{Aa}t_{Ab}t_{ab}`$, which are absent in the Mila-Rice and the Shastry model. The form factors correspond to
$`F_{Aa}^{Bu}(p)`$ $`=`$ $`2F_{Aa}^{Bu}\mathrm{cos}(pa)`$ (39)
$`F_{Ba}^{Bu}(p)`$ $`=`$ $`2F_{Ba}^{Bu}\mathrm{cos}(pa/2).`$ (40)
## III The weak interaction limit
Let us now solve the full model (1) when interactions are weak compared to the bandwidth. This allows us to use the bosonization technique for treating interactions in the undoped as well as in the doped case.
### A NMR properties of the four-band model
#### 1 Reduction to an effective single-band Hamiltonian
Instead of working with the basis $`a,b`$ it is more convenient to diagonalize (3) within a unit cell, and to introduce the bonding and anti-bonding bands. Using the transformation
$`a_{k\sigma }`$ $`=`$ $`\left[\mathrm{cos}(\gamma _k)\alpha _{k\sigma }\mathrm{sin}(\gamma _k)\beta _{k\sigma }\right]e^{i\frac{ka}{2}}`$ (41)
$`b_{k\sigma }`$ $`=`$ $`\mathrm{sin}(\gamma _k)\alpha _{k\sigma }+\mathrm{cos}(\gamma _k)\beta _{k\sigma }`$ (42)
with
$$\mathrm{tan}(2\gamma _k)=\frac{2t_{ab}}{ϵ}\mathrm{cos}(ka/2),\gamma _k[0,\frac{\pi }{4}[$$
(43)
the kinetic energy (3) becomes
$$H_T=\underset{k\sigma }{}\left[ϵ_\alpha (k)\alpha _{k\sigma }^{}\alpha _{k\sigma }+ϵ_\beta (k)\beta _{k\sigma }^{}\beta _{k\sigma }\right],$$
(44)
where the state $`|\alpha _{k\sigma }`$ refers to the lower Hubbard band with energy $`ϵ_\alpha (k)=ϵ/\mathrm{cos}(2\gamma _k)`$, and the state $`|\beta _{k\sigma }`$ to the upper one with energy $`ϵ_\beta (k)=ϵ/\mathrm{cos}(2\gamma _k)`$. In the absence of interactions the chemical potential $`\mu `$ lies in the $`\alpha `$-band both for the undoped and for the doped system, and one can ignore the $`\beta `$-band, which is at least $`ϵ_bϵ_a=2ϵ`$ higher in energy. The same property holds when the interaction terms (5) are added to (3), given that in the weak coupling limit $`U_b,U_a2t_{ab}^2/ϵ`$. Correlation effects in 1D will strongly affect the $`\alpha `$-band states, thus in the following we ignore the terms containing $`\beta `$-operators when substituting (41) into (5).
Substituting (41) into (6) and performing a first order perturbation theory with respect to $`H_S`$, all operators in (1) can be written as
$$\eta _{k\sigma }\lambda _{\eta \alpha }(k)\alpha _{k\sigma }$$
(45)
with
$`\lambda _{a\alpha }(k)`$ $`=`$ $`\mathrm{cos}(\gamma _k)e^{i\frac{ka}{2}}`$ (46)
$`\lambda _{b\alpha }(k)`$ $`=`$ $`\mathrm{sin}(\gamma _k)`$ (47)
$`\lambda _{A\alpha }(k)`$ $`=`$ $`{\displaystyle \frac{2t_{Aa}\mathrm{cos}(ka)}{ϵ_\alpha (k)ϵ_A}}\mathrm{cos}(\gamma _k)e^{i\frac{ka}{2}}`$ (48)
$`+`$ $`{\displaystyle \frac{2t_{Ab}\mathrm{cos}(ka/2)}{ϵ_\alpha (k)ϵ_A}}\mathrm{sin}(\gamma _k)e^{i\frac{ka}{2}}`$ (49)
$`\lambda _{B\alpha }(k)`$ $`=`$ $`{\displaystyle \frac{2it_{Ba}\mathrm{sin}(ka/2)}{ϵ_\alpha (k)ϵ_B}}\mathrm{cos}(\gamma _k).`$ (50)
Here we have assumed an unperturbed ground state $`|\alpha ^{(1)}|\alpha `$. Thus (45) implies that (1) reduces to an effective single-band Hamiltonian.
#### 2 The Continuum limit
We can now use the standard techniques in order to treat interacting one-dimensional systems. Restricting ourselves to the low energy physics regime we make the usual approximation valid for 1D systems, i.e we linearize the spectrum close to the Fermi points, as shown in Fig. 3. Then the Hamiltonian (2) is reduced to
$`H_T`$ $`=`$ $`{\displaystyle \underset{r=\pm ,q,\sigma }{}}rv_Fq\alpha _{rq\sigma }^{}\alpha _{rq\sigma }`$ (51)
$`H_U`$ $`=`$ $`{\displaystyle \underset{𝐫,𝐪}{}}{\displaystyle \frac{U(𝐫)}{N}}\alpha _{r_1,q_1+q_3}^{}\alpha _{r_2,q_2q_3}^{}\alpha _{r_3,q_2}\alpha _{r_4,q_1},`$ (52)
where $`𝐫=(r_1,r_2,r_3,r_4)`$ and $`𝐪=(q_1,q_2,q_3)`$. $`U(𝐫)`$ parameterizes the repulsive interaction in the continuum limit and is given in terms of the standard notations as
$`U(\pm ,\pm ,\pm ,\pm )`$ $``$ $`U_o`$ (53)
$`U(\pm ,,,\pm )`$ $``$ $`U_o`$ (54)
$`U(\pm ,,\pm ,)`$ $``$ $`U_s`$ (55)
$`U(\pm ,\pm ,,)`$ $``$ $`U_c.`$ (56)
$`U_o`$ refers to the two forward scattering processes, $`U_s`$ to the backward scattering, and $`U_c`$ to the umklapp scattering process that occurs at half-filling. The relation to the local repulsions defined in (5) is given by
$`U_o`$ $`=`$ $`U_{b\alpha }+U_{a\alpha }>0`$ (57)
$`U_s`$ $`=`$ $`U_{b\alpha }+U_{a\alpha }>0`$ (58)
$`U_c`$ $`=`$ $`U_{b\alpha }U_{a\alpha }<0,`$ (59)
where $`U_{\eta \alpha }`$ is the short notation for the projected Coulomb energies $`U_\eta |\lambda _{\eta \alpha }(k_F)|^4`$.
As usual for interacting one-dimensional systems, it is useful to introduce a boson representation of the fermion operators, related to the charge and spin density fluctuations. Since the technique is standard, we only recall the main steps and refer the reader to the literature. We rewrite the original density operators in terms of a linear combination of charge ($`\nu =c`$) and spin ($`\nu =s`$) density operators for each branch
$$\rho _{r\sigma }=(\rho _{rc}+\sigma \rho _{rs})/\sqrt{2}.$$
(60)
These density operators define the phase fields
$`\mathrm{\Phi }_\nu (x)`$ $`=`$ $`{\displaystyle \frac{i\pi }{L}}{\displaystyle \underset{r,q0}{}}{\displaystyle \frac{1}{q}}e^{a|q|/2iqx}\rho _{r\nu }`$ (61)
$`\mathrm{\Theta }_\nu (x)`$ $`=`$ $`{\displaystyle \frac{i\pi }{L}}{\displaystyle \underset{r,q0}{}}{\displaystyle \frac{r}{q}}e^{a|q|/2iqx}\rho _{r\nu }.`$ (62)
All operators can be expressed in terms of the boson fields (61), and the fermion operator reads:
$$\alpha _{r\sigma }(x)=\frac{1}{\sqrt{2\pi a}}e^{irk_Fx\frac{i}{\sqrt{2}}\left[r(\mathrm{\Phi }_c+\sigma \mathrm{\Phi }_s)(\mathrm{\Theta }_c+\sigma \mathrm{\Theta }_s)\right]}.$$
(63)
The complete Hamiltonian becomes
$$H=(H_0^c+H_U^c)+(H_0^s+H_U^s)+H_N,$$
(64)
where
$$H_0^\nu =\frac{dx}{2\pi }\left[(u_\nu K_\nu )(\pi \mathrm{\Pi }_\nu )^2+\left(\frac{u_\nu }{K_\nu }\right)(_x\mathrm{\Phi }_\nu )^2\right]$$
(65)
is a quadratic part containing only charge or spin degrees of freedom (with $`\nu =c,s`$). In (65), the variable $`\mathrm{\Pi }_\nu =_x\mathrm{\Theta }_\nu `$ is the momentum density conjugate to $`\mathrm{\Phi }_\nu `$, and thus they respect the commutation relation $`[\mathrm{\Phi }_\nu (x),\mathrm{\Pi }_\nu (x^{})]=i\delta (xx^{})`$. The interaction terms are given by
$`H_U^c`$ $`=`$ $`{\displaystyle 𝑑x\frac{2aU_{b\alpha }}{(2\pi a)^2}\mathrm{cos}[\sqrt{8}\mathrm{\Phi }_c\delta x]}`$ (66)
$``$ $`{\displaystyle 𝑑x\frac{2aU_{a\alpha }}{(2\pi a)^2}\mathrm{cos}[\sqrt{8}\mathrm{\Phi }_c\delta (xa/2)]}`$ (67)
$`H_U^s`$ $`=`$ $`{\displaystyle 𝑑x\frac{2aU_s}{(2\pi a)^2}\mathrm{cos}[\sqrt{8}\mathrm{\Phi }_s]}.`$ (68)
Here $`\delta =4k_F2\pi /a`$ is proportional to the doping of the system with respect to the half-filled case shown in Fig. 1 (for which $`k_F=\pi /2a`$). Using this representation we suppose to work with a fixed number of particles, since $`k_F`$ is directly related to the filling. Finally, the isotropic electron-nuclear interaction part could be written as
$$H_N=𝑑x\left[C_A𝐈(x)𝐒_A(x)+C_B𝐉(x)𝐒_B(x)\right],$$
(69)
where $`H_N`$ is the projection of (11) onto the $`\alpha `$-band using (12) and (45). The projected spin operators $`𝐒_\eta `$ are expressed in terms of (63). For example, the $`z`$-component of the spin operator $`𝐒_\eta `$ can be represented as a sum of $`p0`$ and $`p2k_F`$ components as
$$S_\eta ^z(x)=|\lambda _{\eta \alpha }(k_F)|^2[\overline{s}_\alpha (x)+\stackrel{~}{s}_{\eta \alpha }(x)],$$
(70)
where the non-oscillatory part is given by
$$\overline{s}_\alpha (x)=\frac{1}{\sqrt{2}\pi }(_x\mathrm{\Phi }_s)$$
(71)
and the oscillatory part by
$$\stackrel{~}{s}_{\eta \alpha }(x)=\frac{1}{\pi a}\mathrm{sin}[\sqrt{2}\mathrm{\Phi }_s]\mathrm{sin}[2k_F(xx_\eta )\sqrt{2}\mathrm{\Phi }_c].$$
(72)
The difference between the copper and the oxygen sites is reflected in the value of $`x_\eta `$ and affects the oscillatory part; indeed, for copper $`x_a=x_A=a/2`$ and for oxygen $`x_b=x_B=0`$ as a consequence of the different phase factors in (46).
In (65), the $`u_\nu `$ are the new velocities for the $`\nu `$-excitation and the $`K_\nu `$ are the Luttinger liquid parameters controlling the anomalous exponents in the correlation functions. For weak coupling, they are related to the interactions in (52) by
$`u_sK_s`$ $`=`$ $`u_cK_c=v_F`$ (73)
$`u_s/K_s`$ $`=`$ $`v_FaU_o/\pi `$ (74)
$`u_c/K_c`$ $`=`$ $`v_F+aU_o/\pi .`$ (75)
Since the Luttinger liquid representation is more general than the perturbative result for small interactions, it is also applicable when the interactions are strong. The quadratic Hamiltonian can be viewed in this case as an effective Hamiltonian describing the low-energy properties of the system, provided that the correct Luttinger liquid parameters are used. Such a smooth connection between weak and strong coupling has been proven for single-band models, and a similar Luttinger representation has been shown to work for the case of the two-band model. Equations (6469) define the four-band model, and the NMR properties can be computed through $`H_N`$.
#### 3 Correlation functions at zero temperature
We focus here on the spin-spin correlation functions relevant for NMR and for neutron scattering experiments. The general form of these functions is
$$R_{\eta \eta ^{}}(x,\tau )=T_\tau S_\eta ^{}^z(x,\tau )S_\eta ^z(0,0)S_\eta ^{}^zS_\eta ^z,$$
(76)
and it describes correlations between different orbitals $`\eta `$ and $`\eta ^{}`$ at different points in Euclidean space-time. Here we introduce the decomposition of this function into a non-oscillatory and an oscillatory part
$$R_{\eta \eta ^{}}(x,\tau )=\overline{R}_{\eta \eta ^{}}(x,\tau )+\mathrm{cos}(2k_Fx)\stackrel{~}{R}_{\eta \eta ^{}}(x,\tau ),$$
(77)
since the behavior of these functions will be very different for one-dimensional systems.
Because of the doping dependence in the cosine terms in (66), the behavior of the system will quite clearly be different for zero and for finite doping. At half-filling, one sees from (73), (B2), and (B3) that charge excitations are massive ($`c_m`$), whereas spin excitations are in the massless regime ($`s_o`$). One recovers the standard Mott or charge-transfer insulator with the massless excitations corresponding to a Heisenberg-like exchange. In the doped case, the term (66) is irrelevant because of the oscillatory factor $`\delta x`$. However, at short distances or for short times this term is still small, and the cosine term will influence the behavior of the system. We thus distinguish between two different regimes for the doped case: we assume that for intermediate distances ($`axl_\delta `$) the system remains in the ($`c_m,s_o`$)-phase as mentioned before for the half-filled case, and when distances are larger than $`l_\delta `$, the system will be in the ($`c_o,s_o`$)-phase because the umklapp process becomes ineffective. The characteristic length separating these two regimes denotes essentially the distance between two charge domain walls and is given by $`l_\delta =2\pi /\delta `$.
Due to the spin-charge separation in (64), each part of the correlation function (77) will factorize into independent averages over the spin ($`s_o`$) and the charge sector ($`c_o`$ or $`c_m`$) , and will only depend on the characteristic distance $`r_\nu =[(u_\nu \tau )^2+x^2]^{\frac{1}{2}}`$ between two points in Euclidean space-time (with $`\nu =s,c`$). Details about the correlation functions in the various regimes ($`c_i,s_i`$) are explained in Appendix B. Substituting (7072) in (76), the non-oscillatory contribution to the correlation function is given by
$$\overline{R}_{\eta \eta ^{}}=|\lambda _{\eta \alpha }|^2|\lambda _{\eta ^{}\alpha }|^2\overline{R}_\alpha (r_s),$$
(78)
where $`\overline{R}_\alpha (r_s)=(2\pi r_s)^2`$ depends only on the spin degrees of freedom and is thus completely independent of the coexisting charge phase. Notice that this function is also independent of the orbitals $`\eta `$ and $`\eta ^{}`$, and thus there is no fundamental difference between copper and oxygen contributions.
For the oscillatory part of the spin-spin correlation functions, the situation will be quite different. We restrict ourselves to the calculation of correlation functions between identical orbitals ($`\eta =\eta ^{}`$). Using averages over the charge and spin sectors of the Hamiltonian (64), these functions can be reexpressed as
$$\stackrel{~}{R}_{\eta \eta }^{c_m,s_o}=\frac{|\lambda _{\eta \alpha }|^4}{(2\pi a)^2}\stackrel{~}{R}_{\eta \alpha }^{c_m}(r_c)\stackrel{~}{R}_\alpha ^{s_o}(r_s)$$
(79)
in the massive charge regime and as
$$\stackrel{~}{R}_{\eta \eta }^{c_o,s_o}=\frac{|\lambda _{\eta \alpha }|^4}{(2\pi a)^2}\stackrel{~}{R}_\alpha ^{c_o}(r_c)\stackrel{~}{R}_\alpha ^{s_o}(r_s)$$
(80)
in the massless charge regime. The newly defined correlation functions in the massless phases ($`\nu _o`$) are given by
$$\stackrel{~}{R}_\alpha ^{\nu _o}(r_\nu )=(a/r_\nu )^{K_\nu ^{}}F(r_\nu ).$$
(81)
The function $`F(r_\nu )`$ describes the corrections to the Luttinger Liquid behavior which come from the flow to the fixed point . To lowest order, $`F(r_\nu )`$ can be approximated by $`1`$. The renormalized Luttinger liquid parameters $`K_\nu ^{}`$ for a spin symmetric model with repulsive interaction are restricted to
$$K_s^{}=1\mathrm{and}0K_c^{}1.$$
(82)
The value of the renormalized Luttinger liquid parameter $`K_c^{}`$ depends on the interactions. For weak interaction, $`K_c^{}`$ is close to $`1`$, and it decreases as interactions become more repulsive.
The correlation functions in (79) which are characterized by the massive charge phase are given by
$`\stackrel{~}{R}_{A\alpha }^{c_m}(r_c)`$ $`=`$ $`2\mathrm{cosh}[K_cK_0(m_cr_c)](m_ca)^{K_c}`$ (83)
$`\stackrel{~}{R}_{B\alpha }^{c_m}(r_c)`$ $`=`$ $`2\mathrm{sinh}[K_cK_0(m_cr_c)](m_ca)^{K_c}`$ (84)
and depend on the chosen orbital $`\eta `$. Thus, the behavior for copper and oxygen will be quite different. It depends on the distance $`r_c`$, the mass $`m_c`$, and the stiffness constant $`K_c`$. In general, for distances larger than $`l_{m_c}=1/m_c`$, the function $`\stackrel{~}{R}_{A\alpha }^{c_m}`$ for copper tends to a finite constant, whereas $`\stackrel{~}{R}_{B\alpha }^{c_m}`$ for oxygen tends exponentially to zero.
#### 4 The asymptotic expressions at finite temperature
In order to obtain the temperature dependent correlation function $`\stackrel{~}{R}_{\eta \eta }(x,\tau ,\beta )`$, we will only use the asymptotic expressions of (84). We recover a Luttinger liquid behavior, and the temperature dependence can easily be obtained with the help of the conformal symmetry; indeed, we only need to replace $`r_\nu (x,\tau )`$ by $`r_\nu (x,\tau ,\beta )`$ where
$$r_\nu (x,\tau ,\beta )=\frac{u_\nu \beta }{\pi }\sqrt{\mathrm{sinh}\left[\frac{xiu_\nu \tau }{u_\nu \beta /\pi }\right]\mathrm{sinh}\left[\frac{x+iu_\nu \tau }{u_\nu \beta /\pi }\right]}.$$
(85)
The relevant asymptotic expressions at half-filling and away from half-filling depend on the relative magnitudes of the various characteristic lengths of the system, namely the lengths related to the mass, $`l_{m_c}`$, and to the doping, $`l_\delta `$, as well as the thermal length $`l_\beta =\mathrm{min}\{(u_c\beta )^1,(u_s\beta )^1\}`$. For half-filling ($`\gamma =0`$) at low temperature, we are in the regime where $`l_\beta r_cl_{m_c}`$ and $`l_\delta =\mathrm{}`$, thus we can approximate the oscillatory charge contribution in (79), and the functions are simplified to
$$\stackrel{~}{R}_{\eta \eta }^0=C_\eta ^0|\lambda _{\eta \alpha }|^4\stackrel{~}{R}_\alpha ^0(r_s),$$
(86)
where the amplitude of the oscillatory part at half-filling is given by $`C_\eta ^0=\stackrel{~}{R}_{\eta \alpha }^{c_m}(\mathrm{}_c)`$. The remaining correlation function is independent of the orbital $`\eta `$ and given by $`\stackrel{~}{R}_\alpha ^0=(2\pi a)^2(a/r_s)`$.
The large distance limit of the corresponding expression for small doping ($`\gamma =\delta `$) and low temperature, where $`l_\beta r_cl_\delta l_{m_c}`$, looks like
$$\stackrel{~}{R}_{\eta \eta }^\delta =C_\eta ^\delta |\lambda _{\eta \alpha }|^4\stackrel{~}{R}_\alpha ^\delta (r_c,r_s).$$
(87)
In the doped regime, $`C_\eta ^\delta `$ is the amplitude $`\stackrel{~}{R}_{\eta \alpha }^{c_m}(l_\delta )`$ obtained in the massive phase at the crossover, as shown in Fig. 4. Like before, the remaining correlation function $`\stackrel{~}{R}_\alpha ^\delta =(2\pi a)^2(a/r_s)(a/r_c)^{K_c^{}}`$ is also independent of $`\eta `$ but shows dependence on spin and charge degrees of freedom. For larger doping rates ($`l_\delta <l_{m_c}`$), the difference between copper and oxygen sites vanishes.
#### 5 Knight shifts and relaxation rates
The standard expressions for the Knight shifts and for the relaxation rates resulting from a hyperfine coupling term like (11) are
$`K_\eta ^\gamma `$ $`=`$ $`{\displaystyle \frac{C_\eta }{\gamma _\eta \gamma _e\mathrm{}^2}}{\displaystyle \underset{\eta ^{}=a,b}{}}\chi _{\eta \eta ^{}}^\gamma (\omega =0,p0)`$ (88)
$`{\displaystyle \frac{1}{T_{1\eta }^\gamma }}`$ $`=`$ $`{\displaystyle \frac{C_\eta ^2}{\gamma _\eta \gamma _e\mathrm{}^2\beta }}{\displaystyle \underset{p}{}}{\displaystyle \frac{\mathrm{Im}\left[\chi _{\eta \eta }^\gamma (\omega _\eta ,p)\right]}{\omega _\eta }},`$ (89)
where $`\gamma =0`$ refers to the half-filled case and $`\gamma =\delta `$ to the doped case. $`\omega _\eta `$ denotes the electronic Zeeman frequency in orbital $`\eta `$, which is very small as compared to the energy scale of the purely electronic system fixed by the cutoff $`\lambda `$. For the Knight shifts the sum is restricted to the active orbitals $`a`$ and $`b`$. We can split up the susceptibility $`\chi _{\eta \eta ^{}}^\gamma `$ into the non-oscillatory $`\overline{\chi }_{\eta \eta ^{}}^\gamma `$ and the oscillatory contribution $`\stackrel{~}{\chi }_{\eta \eta ^{}}^\gamma `$ just like for the correlation functions in (77). Finally, the Knight shifts for the linearized four-band model in units of $`C_\eta /(\gamma _\eta \gamma _e\mathrm{}^2)`$ are given by
$$K_\eta ^\gamma =\overline{F}_{\eta \alpha }\overline{\chi }_\alpha (\omega =0,q0)$$
(90)
and the relaxation rates in units of $`C_\eta ^2/(\gamma _\eta \gamma _e\mathrm{}^2\omega _\eta )`$ by
$$\frac{1}{T_{1\eta }^\gamma }=\frac{1}{\beta }\underset{|q|<\lambda }{}\mathrm{Im}\left[(\overline{F}_{\eta \alpha })^2\overline{\chi }_\alpha (\omega _\eta ,q)+(\stackrel{~}{F}_{\eta \alpha }^\gamma )^2\stackrel{~}{\chi }_\alpha ^\gamma (\omega _\eta ,q)\right].$$
(91)
The susceptibilities $`\overline{\chi }_\alpha `$ and $`\stackrel{~}{\chi }_\alpha ^\gamma `$ in space-time can be obtained from
$`\overline{\chi }_\alpha (x,t)=2\theta (t)\mathrm{Im}\left[\overline{R}_\alpha (x,\tau ,\beta )\right]_{\tau =it+ϵ}`$ (92)
$`\stackrel{~}{\chi }_\alpha ^\gamma (x,t)=2\theta (t)\mathrm{Im}\left[\stackrel{~}{R}_\alpha ^\gamma (x,\tau ,\beta )\right]_{\tau =it+ϵ}`$ (93)
performing the continuation to real time. The $`\tau `$-ordered temperature dependent Green’s functions on the right-hand side are the same as in (78), (86) and (87), using (85). Thus, in general we can calculate (90) and (91) by performing the Fourier transform of (92) and (93). Here we restrict ourselves to the solutions obtained by the so called power counting method. The temperature dependences of the Knight shifts and of the relaxation rates are shown in Table I and the form factors $`\overline{F}_{\eta \alpha }`$ and $`\stackrel{~}{F}_{\eta \alpha }^\gamma `$ are given in Table II.
### B NMR properties of the Bulut model
In order to obtain the NMR properties of the 1D version of the Bulut model we perform the same procedure as before for the four-band model. The bosonized version of the Bulut model (34) is given by (64) replacing $`H_N`$ by
$$H_N^{Bu}=𝑑x\left[C_A𝐈(x)𝐒_A^{Bu}(x)+C_B𝐉(x)𝐒_B^{Bu}(x)\right],$$
(94)
where $`𝐒_\eta ^{Bu}`$ is the projection onto the $`\alpha `$-band using (35), (36), (12) and (45). For the z-component of the spin $`𝐒_\eta ^{Bu}`$ we get
$`S_A^{zBu}(x)`$ $`=`$ $`2|\lambda _{Aa}^{Bu}|^2|\lambda _{a\alpha }|^2[\overline{s}_\alpha (x)+\mathrm{cos}(2k_Fa)\stackrel{~}{s}_{a\alpha }(x)]`$ (95)
$`S_B^{zBu}(x)`$ $`=`$ $`2|\lambda _{Ba}^{Bu}|^2|\lambda _{a\alpha }|^2[\overline{s}_\alpha (x)+\mathrm{cos}(k_Fa)\stackrel{~}{s}_{a\alpha }(x)]`$ (96)
ignoring all gradient terms of the field $`\varphi _c`$. The spin operators $`\overline{s}_\alpha `$ and $`\stackrel{~}{s}_{\eta \alpha }`$ are defined as before in (71) and (72). The NMR properties for a hyperfine coupling like (94) are given by (90) and (91) by the replacements $`\overline{F}_{\eta \alpha }\overline{F}_{\eta \alpha }^{Bu}`$ and $`\stackrel{~}{F}_{\eta \alpha }^\gamma \stackrel{~}{F}_{\eta \alpha }^{\gamma Bu}`$. The values for the different form factors are shown in Table III. At this level of approximation, both the four-band and the Bulut model show exactly the same temperature dependence for the Knight shifts as well as for the relaxation rates; this dependence is different for the uniform contribution and for the oscillatory one (see Table I), as is well known for interacting one-dimensional systems. This effect has nothing to do with the various orbitals where the Knight shifts and the relaxation rates are measured.
### C Comparing the four-band model and the 1D Bulut model
First we focus on the coefficients of the Bulut and of the four-band model (compare Table II and III) related to the different projection procedures of the s-orbitals onto the ground state. For comparing both models, we investigate the limit $`t_{ab}(ϵ_bϵ_a)`$. Then, for the four-band model the projection of the s orbitals ($`A,B`$) onto the lowest band ($`\alpha `$) is strictly done in $`k`$-space and results in
$`|\lambda _{B\alpha }(k_F)|^2`$ $``$ $`2(\lambda _{Ba}^{Sh})^2`$ (98)
$``$ $`2\mathrm{cos}(k_Fa)(\lambda _{Ba}^{Sh})^2`$ (99)
$`|\lambda _{A\alpha }(k_F)|^2`$ $``$ $`2(\lambda _{Aa}^{Sh}+\lambda _{Aa}^{Mi})^2`$ (101)
$`+`$ $`4(\lambda _{Aa}^{Mi})^2`$ (102)
$`+`$ $`2\mathrm{cos}(2k_Fa)(\lambda _{Aa}^{Sh}+\lambda _{Aa}^{Mi})^2`$ (103)
$`+`$ $`8\mathrm{cos}(k_Fa)(\lambda _{Aa}^{Sh}+\lambda _{Aa}^{Mi})\lambda _{Aa}^{Mi},`$ (104)
whereas for the Bulut model it is a combination of real space and $`k`$-space projection yielding
$`2|\lambda _{Ba}^{Bu}|^2|\lambda _{a\alpha }(k_F)|^2`$ $``$ $`2(\lambda _{Ba}^{Sh})^2`$ (105)
$`2|\lambda _{Aa}^{Bu}|^2|\lambda _{a\alpha }(k_F)|^2`$ $``$ $`2(\lambda _{Aa}^{Sh}+\lambda _{Aa}^{Mi})^2.`$ (107)
The general solution for the projected O-3s orbital includes one more term (99) than the solution proposed by Bulut (105). This term corresponds to a dynamic contribution which includes a charge displacement. However, for a half-filled system the additional term vanishes and the two solutions become identical. By contrast, the projection procedure for the Cu-4s orbital produces a completely different behavior in the two models. For a half-filled system, the hopping processes via $`\lambda _{Aa}^{Bu}=\lambda _{Aa}^{Sh}+\lambda _{Aa}^{Mi}`$ contribute only in the Bulut model (107), whereas they are exactly canceled by the related dynamic terms (103) in the four-band model. Thus, for the four-band model at half-filling, only an additional local term (102) as well as a dynamic combined term (104) remain. The term (102) is the local analog to the transferred terms proposed by Mila-Rice, and the term (104) is a combination of Mila-Rice and Shastry terms which includes a charge displacement. It should be clear that our projection procedure is the right one for a system with small Coulomb interactions: First we diagonalize the tight-binding Hamiltonian dealing with extended wave functions, and then we treat the Coulomb energy approximately within this non-local basis. The approximation proposed by Bulut suffers from a mismatch between the local and the non-local point of view.
The second part of the oscillatory contribution to the form factors (compare Table II and III), which contains the dependence on the characteristic lengths related to the doping rate, $`l_\delta `$, as well as to the charge mass, $`l_{m_c}`$, is the crucial one. Away from half-filling, the four-band model shows a different behavior on the copper and on the oxygen, despite the fact that umklapp processes only contribute on short or intermediate scales. Indeed, the different hyperbolic dependencies of the two characteristic lengths $`l_{m_c}`$ and $`l_\delta `$ for copper and for oxygen (see Table II) affect measured quantities related to long distance or long time behavior. Instead, for the Bulut model the difference between copper and oxygen comes in only because of the special choice of a Mila-Rice-Shastry type electron-nuclear interaction term (94) and the related unconventional projection procedure which results in the different trigonometric form factors (see Table III). The influence of the charge mass $`m_c`$ is the same for copper and for oxygen, a fact which manifests itself by the same dependence on the characteristic length $`l_{m_c}`$.
Note that the four-band model leads to a very small contribution on the oxygen even at finite doping, because the contribution is exponentially suppressed in a way which depends on the ratio between $`l_{m_c}`$ and $`l_\delta `$, whereas the oscillatory contribution on the copper atom is nearly independent of the doping rate for long distances or times. For the ratio between copper and oxygen we distinguish between two regimes:
$$\mathrm{tanh}\left[K_cK_0\left(\frac{l_\delta }{l_{m_c}}\right)\right]\{\begin{array}{cc}1\mathrm{for}\hfill & l_\delta l_{m_c}\hfill \\ 0\mathrm{for}\hfill & l_\delta l_{m_c}\hfill \end{array}.$$
(108)
In the former regime, we recover the Luttinger Liquid behavior, since the infinite length $`l_{m_c}`$ stems from the vanishing of the umklapp process when $`U_{a\alpha }=U_{b\alpha }`$; in that case there is no fundamental difference between copper and oxygen anymore. Only the overlaps with the ground state remain different. The latter regime, where the fundamental difference occurs, will be reached exponentially as $`K_c\sqrt{\pi l_{m_c}/2l_\delta }\mathrm{exp}(l_\delta /l_{m_c})`$, and thus the oxygen does not see the antiferromagnetic fluctuations in this limit. Instead, for the Bulut model everything depends on the same correlation function, and the difference between copper and oxygen comes from the filtering factors. Thus, the oscillating contributions to the relaxation rates for the oxygen is always proportional to $`[0+(\pi a/2l_\delta )^2]`$, whereas the contributions on copper are reduced by a factor $`[1(\pi a/l_\delta )^2]`$. The ratio of the oscillating contribution to the relaxation rates is approximately given by $`(\pi a/2l_\delta )^2`$, and is completely independent of the details of the projected local Coulomb repulsions $`U_{a\alpha }`$ and $`U_{b\alpha }`$. It only depends on the doping rate and is proportional to $`(\delta a/4)^2`$. By contrast, the four-band model includes the effect of the Coulomb interactions through its dependence on $`l_{m_c}`$. In Table IV we show the ratios of the different Knight shifts and relaxation rates contributions.
## IV The strong interaction limit
For strong interactions the four-band system in (1) is much more difficult to solve. Yet, it is still possible to highlight the qualitative features of the transferred hyperfine coupling interaction, specifically for the half-filled case. To obtain the strong interaction limit of this model we can perform the Gutzwiller projection
$$H=\widehat{P}(H_0+H_S+H_N)\widehat{P}$$
(109)
which eliminates doubly occupied states in the Cu-3d orbitals from the Fock space. The projection is effectively performed on all three terms of (109), which are treated on equal footing. As far as the first part $`\widehat{P}H_0\widehat{P}`$ is concerned, two possible superexchange processes are generated, as shown in Fig. 5. In the strong interaction limit ($`U_a|ϵ_\eta ϵ_\eta ^{}|,U_bt_{\eta \eta ^{}}`$) the superexchange process in Fig. 5(a) is much more effective than the process in Fig. 5(b). For the basic system $`H_0`$, we only have to keep 3 states per unit cell, whereas for the four-band model (109), we end up with a system where we have to keep 27 spin-degenerate local states per unit cell $`j`$ with 4 tight-binding parameters $`t_{\eta \eta }`$ for a half-filled system (excluding doubly excited $`A,B`$-states, see Appendix C). For a doped system the number of states as well as the number of possible transitions increases very fast, as has been shown for a two-band model. A correct projection procedure such as (109) becomes very difficult to handle, and one must resort to some approximations. In any event, in the vicinity of the half-filled case where the projection can be explicitly used for the full Hamiltonian, we will analyze the differences between the predictions of the four-band model and those of the approximated Hamiltonians. So let us restrict our analysis to the half-filled case where only virtual double occupancies of the copper site are allowed and where electron-nuclear interaction processes require that the initial and the final charge distribution be the same. We deal with electron-nuclear interaction processes where effectively one local Cu-3d spin will be reversed and then relaxed by the thermodynamic fluctuations of the Heisenberg model. We decompose $`(H_0+H_S+H_N)`$ into $`(L+K)`$. $`L`$ includes all local and $`K`$ all kinetic contributions of the complete Hamiltonian $`H`$ introduced in (1). Then we can expand $`\widehat{P}H\widehat{P}`$ on the basis of the unperturbed eigenstates of $`L`$ and compute the projected local s-orbital spin operators like $`\widehat{P}𝐒\eta j\widehat{P}`$. For the details we refer to Appendix C and discuss only the final results.
First we analyze some relaxation processes for the oxygen atom. The process shown in Fig. 6(a) is a transferred (T) contribution proportional to
$$F_{B,T,(a)}=\left[\frac{t_{Ba}}{ϵ_B(ϵ_a+U_a)}\right]^2.$$
(110)
For the process shown in Fig. 6(b), we include a part of the superexchange process to avoid double occupation of the copper site, and the contribution is proportional to
$$F_{B,T,(b)}=\left[\frac{t_{ab}t_{Ba}}{(ϵ_aϵ_b)(ϵ_Bϵ_b)}\right]^2.$$
(111)
Then the lowest order contribution to the general form factor for the oxygen is given by
$$F_B(p)=2F_{B,L}+2F_{B,T}\mathrm{cos}(pa/2)$$
(112)
with
$`F_{B,L}`$ $`=`$ $`0`$ (113)
$`F_{B,T}`$ $`=`$ $`\underset{\mathrm{projected}\mathrm{Shastry}}{\underset{}{n_{B,T,(a)}F_{B,T,(a)}}}+n_{B,T,(b)}F_{B,T,(b)}+\mathrm{}.`$ (114)
$`n_{B,T,(i)}`$ denotes the combinatorial factor associated with all possible processes yielding a contribution $`F_{B,T,(i)}`$. The factor 2 for the left-right symmetry is not included in $`n_{B,T,(i)}`$. Like for the superexchange processes (Fig. 5) some processes are forbidden due to the Pauli principle. However, since all energy levels are assumed to be spin-independent the related amplitudes $`F_{B,T,(i)}`$ are the same.
For copper we also distinguish between the transferred (Fig. 7) and the local contributions (Fig. 8). The transferred contributions are proportional to
$`F_{A,T,(a)}`$ $`=`$ $`\left[{\displaystyle \frac{t_{Ab}t_{ab}}{(ϵ_Aϵ_b)(ϵ_A+ϵ_a2ϵ_bU_b)}}\right]^2`$ (115)
$`F_{A,T,(b)}`$ $`=`$ $`\left[{\displaystyle \frac{t_{Ab}t_{ab}}{(ϵ_Aϵ_b)(ϵ_Aϵ_aU_a)}}\right]^2`$ (116)
$`F_{A,T,(c)}`$ $`=`$ $`\left[{\displaystyle \frac{t_{Aa}}{ϵ_Aϵ_aU_a}}\right]^2`$ (117)
$`F_{A,T,(d)}`$ $`=`$ $`{\displaystyle \frac{t_{Ab}t_{ab}t_{Aa}}{(ϵ_Aϵ_b)^2(ϵ_aϵ_b)}},`$ (118)
whereas the local contributions are given by
$`F_{A,L,(a)}`$ $`=`$ $`\left[{\displaystyle \frac{t_{Ab}t_{ab}}{(ϵ_Aϵ_b)(ϵ_A+ϵ_a2ϵ_bU_b)}}\right]^2`$ (119)
$`F_{A,L,(b)}`$ $`=`$ $`\left[{\displaystyle \frac{t_{Ab}t_{ab}}{(ϵ_Aϵ_b)(ϵ_Aϵ_aU_a)}}\right]^2`$ (120)
$`F_{A,L,(c)}`$ $`=`$ $`\left[{\displaystyle \frac{t_{Ab}t_{ab}}{(ϵ_Aϵ_b)(ϵ_A+ϵ_a2ϵ_b)}}\right]^2.`$ (121)
Then, the general form factor for copper reads
$$F_A(p)=2F_{A,L}+2F_{A,T}\mathrm{cos}(pa)$$
(122)
with
$`F_{A,L}`$ $`=`$ $`n_{A,L,(a)}F_{A,L,(a)}+n_{A,L,(b)}F_{A,L,(b)}`$ (123)
$`+`$ $`n_{A,L,(c)}F_{A,L,(c)}+\mathrm{}`$ (124)
$`F_{A,T}`$ $`=`$ $`\underset{\mathrm{projected}\mathrm{Mila}\mathrm{Rice}}{\underset{}{n_{A,T,(a)}F_{A,T,(a)}+n_{A,T,(b)}F_{A,T,(b)}}}`$ (126)
$`+`$ $`\underset{\mathrm{projected}\mathrm{Shastry}}{\underset{}{n_{A,T,(c)}F_{A,T,(c)}}}+n_{A,T,(d)}F_{A,T,(d)}+\mathrm{}.`$ (127)
Let us now compare the predictions of the four-band model and those of the 1D Mila-Rice or Shastry models.
Using the projected expression for the oxygen instead of the unprojected one (30), only process (b) in (114) contributes in the strong interaction limit, whereas process (a) in (114) proposed by Shastry becomes negligible
$$F_{B,T,(a)}\stackrel{U_a\mathrm{}}{}0.$$
(128)
The form factor for the characteristic wave vectors ($`p=0`$ or $`p=\pi /a`$) is then reduced to
$`F_B(0)`$ $`=`$ $`2n_{B,T,(b)}F_{B,T,(b)}`$ (129)
$`F_B(\pi /a)`$ $`=`$ $`0.`$ (130)
Since only the relaxation process of the oxygen nuclear spin contributes, which corresponds to $`p0`$ , we recover the basic structure of the form factor of Shastry with modified amplitudes. Thus at half-filling, there is no fundamental difference for the oxygen between the general form factor (112) and the form factor proposed by Shastry (33).
In the strong coupling limit at half filling, the following contributions to the form factor for copper are suppressed:
$$F_{A,L,(b)},F_{A,T,(b)},F_{A,T,(c)}\stackrel{U_a\mathrm{}}{}0.$$
(131)
Thus the projected Shastry contribution (c) in (127) and one of the projected Mila-Rice contributions (b) in (127) as well as one of the projected local contributions (b) in (124) become negligible, and we end up with
$`F_{A,L}`$ $`=`$ $`n_{A,L,(a)}F_{A,L,(a)}+n_{A,L,(c)}F_{A,L,(c)}+\mathrm{}`$ (132)
$`F_{A,T}`$ $`=`$ $`n_{A,T,(a)}F_{A,T,(a)}+n_{A,T,(d)}F_{A,T,(d)}+\mathrm{}`$ (133)
for the local and for the transferred contributions to the general form factor (122), respectively. Thus the uniform part of the form factor is given by
$`F_A(0)`$ $`=`$ $`4n_{A,(a)}F_{A,(a)}+2n_{A,L,(c)}F_{A,L,(c)}`$ (134)
$`+`$ $`2n_{A,T(d)}F_{A,T(d)}+\mathrm{},`$ (135)
whereas the oscillating part reads
$$F_A(\pi /a)=2n_{A,L,(c)}F_{A,L,(c)}2n_{A,T,(d)}F_{A,T,(d)}+\mathrm{}.$$
(136)
We used the fact that $`n_{A,L,(a)}=n_{A,T,(a)}n_{A,(a)}`$ and $`F_{A,L,(a)}=F_{A,T,(a)}F_{A,(a)}`$. The uniform part includes contributions which are absent in the 1D version of the Mila-Rice and of the Shastry model. Furthermore, some terms proposed by Shastry turn out to be zero in the strongly interacting limit. For the oscillatory part the effects are much more drastic. The transferred terms proposed by Shastry vanish in the strong coupling regime, whereas other transferred terms, which come from a combination of Mila-Rice and Shastry processes, contribute. Besides, the transferred terms proposed by Mila and Rice are canceled by the equivalent local terms. Hence, in 1D, the general form factor differs both qualitatively and quantitatively from the form factors one would derive from the Mila-Rice or from the Shastry models.
## V Discussion and perspectives
In this paper, we have analyzed the 1D analogs of the hyperfine form factors proposed for NMR measurements of high-$`\mathrm{T}_\mathrm{c}`$ materials in the antiferromagnetic phase. We have focused on the situation where one deals with an antiferromagnet generated by a superexchange process via an oxygen atom located at the midpoint between two copper atoms and where the Fermi contact interaction is one of the main contributions to the possible electron-nuclear interaction terms. We have investigated a 1D Cu-O model including four orbitals per unit cell, namely the Cu-3d and the O-2p orbitals governing the ground state properties, as well as the Cu-4s and O-3s orbitals describing the isotropic Fermi contact interaction. In 1D, we were able to solve this model using only standard techniques without having to introduce any additional approximations for the hyperfine interaction term as proposed by Mila-Rice and by Shastry. Thus, we were able to compare our solutions of the four-band model with the predictions of the approximative models.
In the low interaction limit, we have calculated the resulting temperature dependence of the Knight shifts $`K`$ and of the relaxation rates $`1/T_1`$ for an undoped and for a doped system; in that limit the ground state is well described by the strongly hybridized Cu-3d–O-2p anti-bonding band the width of which is large as compared to all Coulomb interactions. For both models, the four-band and the approximative one (Bulut model), the temperature dependences are the same and show the typical power law behavior of one-dimensional interacting systems (see Table I). Within this scope we have shown for the four-band model that for an undoped and a slightly doped system copper and oxygen behave completely different for long distances or long times, when the temperature is low enough. The oxygen nuclei see only the Korringa-like contributions, since the antiferromagnetic contributions are exponentially suppressed depending on the ratio of the characteristic length related to the charge gap and the doping. In contrast, the copper nuclei always see both contributions, the Korringa-like contribution as well as the antiferromagnetic one. This fundamental difference between copper and oxygen vanishes gradually when the characteristic doping length or the characteristic thermal length becomes shorter than the length related to the charge gap (the difference goes away abruptly when the system develops a gap in the spin sector). This solution is at variance with the prediction of the related approximate model, where for oxygen the antiferromagnetic contributions to $`1/T_1`$ increase with doping like $`\delta ^2`$, whereas for copper they decrease proportionally to $`\delta ^2`$. Thus, the scenario where oxygen does not see the antiferromagnetic fluctuations is realized much more effectively in the four-band model than in the 1D version of the models proposed for the high-$`T_\mathrm{c}`$ materials. In 1D, such an unconventional scenario works, since even small interactions generate strong antiferromagnetic correlations due to the drastic reduction of the Fermi surface.
We have also considered the strong interaction limit. Performing a Gutzwiller projection onto the four-band model without further approximations for the electron-nuclear interaction term, we computed the various processes contributing to NMR. Our analysis was limited to the insulating phase (Heisenberg model), since even in 1D a full solution of the model for a doped system (t-J model with four orbitals per unit cell) is unavailable. In the strong interaction limit of the 1D Cu-O model, we were able to compare the form factors obtained for the four-band model with the predictions obtained for the approximate models (Mila-Rice model and Shastry model) investigating the different relaxation processes for the copper and oxygen nuclear spins. In this context, we have shown that neither the 1D analog of the Mila-Rice model nor the 1D analog of the Shastry model could describe the strong interaction limit at half-filling. In contrast to the usual assumption that only transferred contributions are relevant, we predict that both local and transferred contributions should be taken into account for describing the relaxation of the nuclear copper spin via an Cu-4s orbital. Furthermore, we have shown that for infinite local repulsions on the copper sites and small local repulsions on the oxygen sites, the contributions proposed by Mila-Rice and Shastry vanish. For the relaxation of the nuclear oxygen spin we recover the basic idea of transferred hyperfine couplings with slightly modified amplitudes, but once again the contribution proposed by Shastry vanishes for infinite repulsion on the copper site.
Both the strong and the weak coupling limits underscore the importance of keeping the full four-band model, at least in one dimension, in order to give an accurate description of the NMR properties. The method we used in the present paper to tackle such a model can thus be extended in various directions. First, it can be applied to study specific models which have a structure similar to the model Cu-O chain analyzed here. This is for example the case for ladder materials such as $`\mathrm{Sr}_{14\mathrm{x}}\mathrm{Ca}_\mathrm{x}\mathrm{Cu}_{24}\mathrm{O}_{41}`$. Analyses of the NMR material have so far been performed in terms of Mila-Rice-Shastry approximations. An analysis retaining the full four-band model, with the specific symmetries of these ladder systems, is currently in progress. Other systems for which our analysis can be relevant are TMTSF and TMTTF alloys. At stoichiometric composition they form an alternate stack. Let us now comment on anisotropic contributions to $`K`$ and to $`1/T_1`$; these can be produced by a dipolar hyperfine coupling (see Appendix A). They also stem from the specific structural details of a given compound which may lead to an anisotropic form for the susceptibility: in that situation the anisotropy of the $`p=0`$ component (92) will usually be different from that of the $`p=2k_F`$ part (93). In both the weak and the strong interaction limits, we find that – for low enough temperature – $`1/T_1`$ is mostly determined by (93), whereas $`K`$ is proportional to (92). The experimental observation that $`1/T_1`$ is essentially isotropic and that $`K`$ is anisotropic suggest that anisotropic effects are not too important for the $`p=2k_F`$ contributions but do affect the $`p=0`$ terms.
Another possible extension of our work concerns of course the two-dimensional systems. Although it is unclear how much of the weak coupling approach remains valid in higher dimension, our strong coupling analysis can straightforwardly be applied to higher dimensional structures. The main difference in that case between the 2D (or higher) and the 1D study presented here comes from the symmetry of the various orbitals. In the case of a Cu-O plane, in the presence of a Coulomb repulsion on the oxygen sites ($`U_{\text{O-2p}}0`$), the related amplitudes for the local processes (119) and (121) are not equal anymore, and a cancellation of these terms by symmetry arguments as assumed by Mila-Rice does not occur. Only the contribution like (120) will vanish by symmetry arguments. The transferred Mila-Rice contributions (115) via the O-2p orbital, which always cost the Coulomb energy $`U_{\text{O-2p}}`$, and the local processes (119) have exactly the same combinatorial factor and the same amplitude; thus the term (119) cancels out the term (115) for the antiferromagnetic wave vector. This suggests for $`U_{\text{Cu-3d}}\mathrm{}`$ that the antiferromagnetic contribution to the relaxation of the copper nuclei via an isotropic interaction comes from local terms (see (121)) and from new transferred combined terms of third order (see (118)), while the transferred contributions proposed up to now are absent (see (116) and (117)).
###### Acknowledgements.
This work was initiated by a suggestion from the late H.J. Schulz whom we wish to acknowledge here.
## A Anisotropic hyperfine couplings
Taking into account anisotropic hyperfine couplings related to the orbitals Cu-3d and O-2p we have to replace (11) by
$$H_N^{}=\underset{j\varsigma }{}C_AI_j^\varsigma S_{Aj}^\varsigma +C_a^\varsigma I_j^\varsigma S_{aj}^\varsigma +C_BJ_j^\varsigma S_{Bj}^\varsigma +C_b^\varsigma J_j^\varsigma S_{bj}^\varsigma $$
(A1)
The sum on $`\varsigma `$ is over components of the diagonal hyperfine tensors $`C_\eta ^\varsigma `$.
In the weak interaction limit of the four-band model we can perform the same calculations as done before, and we will end up with the bosonized expression (64), where now the electron-nuclear interaction is given by
$`H_N`$ $`=`$ $`{\displaystyle \underset{\varsigma }{}}{\displaystyle 𝑑x\left(C_A|\lambda _{A\alpha }|^2+C_a^\varsigma |\lambda _{a\alpha }|^2\right)I^\varsigma \left(\overline{S}_\alpha ^\varsigma +\stackrel{~}{S}_{a\alpha }^\varsigma \right)}`$ (A2)
$`+`$ $`{\displaystyle \underset{\varsigma }{}}{\displaystyle 𝑑x\left(C_B|\lambda _{B\alpha }|^2+C_b^\varsigma |\lambda _{b\alpha }|^2\right)J^\varsigma \left(\overline{S}_\alpha ^\varsigma +\stackrel{~}{S}_{b\alpha }^\varsigma \right)}.`$ (A3)
Here we used the fact that $`\stackrel{~}{S}_{A\alpha }^\varsigma =\stackrel{~}{S}_{a\alpha }^\varsigma `$ and $`\stackrel{~}{S}_{B\alpha }^\varsigma =\stackrel{~}{S}_{b\alpha }^\varsigma `$. In general the explicit bosonized expressions for the spin part of $`\stackrel{~}{S}_{\eta \alpha }^x`$ and $`\stackrel{~}{S}_{\eta \alpha }^y`$ in (72) differ from $`\stackrel{~}{S}_{\eta \alpha }^z`$, but finally for a spin-symmetric model there will be no influence on the correlation functions. Thus only the coefficients are slightly modified and vary for the different directions $`\varsigma =x,y,z`$. Formally the contributions to the $`\varsigma `$-directions of the Knight shifts $`K_{Cu}^\varsigma `$ and $`K_O^\varsigma `$, as well as the contributions to the relaxation times $`T_{1,Cu}^\varsigma `$ and $`T_{1,O}^\varsigma `$ are given by (90) and (91) performing the replacements
$`|\lambda _{A\alpha }|^2`$ $``$ $`|\lambda _{A\alpha }|^2+{\displaystyle \frac{C_a^\varsigma }{C_A}}|\lambda _{a\alpha }|^2`$ (A5)
$`|\lambda _{B\alpha }|^2`$ $``$ $`|\lambda _{B\alpha }|^2+{\displaystyle \frac{C_b^\varsigma }{C_B}}|\lambda _{b\alpha }|^2`$ (A6)
in the expressions of the form factors defined in Table II.
In the strong interaction limit of the four-band model the inclusion of anisotropic hyperfine interactions results in
$`H_N^{\prime \prime }`$ $`=`$ $`{\displaystyle \underset{j\varsigma }{}}\left(2C_AF_{A,L}+C_a^\varsigma F_{a,L}\right)I_j^\varsigma S_{aj}^\varsigma `$ (A7)
$`+`$ $`{\displaystyle \underset{j\varsigma }{}}F_{A,T}I_j^\varsigma (S_{a,j1}^\varsigma +S_{a,j+1}^\varsigma )`$ (A8)
$`+`$ $`{\displaystyle \underset{j\varsigma }{}}\left(C_BF_{B,T}+C_b^\varsigma F_{b,T}\right)J_j^\varsigma \left(S_{aj}^\varsigma +S_{a,j+1}^\varsigma \right).`$ (A9)
Here the new defined parameters which describe the additional couplings to the local Cu-3d spins are given by $`F_{a,L}=1`$ and $`F_{b,T}=t_{ab}^2/(ϵ_aϵ_b)^2`$, whereas all the others were defined in Section IV. For the copper atom the local contribution is modified, whereas for the oxygen atom it is the transferred one.
## B The sine-Gordon model
At half-filling the spin part as well as the charge part of the Hamiltonian (64) are described by a sine-Gordon model $`H_{SG}^\nu =H_0^\nu +H_U^\nu `$ where
$$H_U^\nu =\frac{2aU_\nu }{(2\pi a)^2}_0^L𝑑x\mathrm{cos}[\sqrt{8}\mathrm{\Phi }_\nu ]$$
(B1)
For this model two different regimes exist depending on the value of the parameter $`K_\nu `$. A massive regime ($`\nu _m`$) for
$$2\pi u_\nu (K_\nu 1)<|U_\nu |,$$
(B2)
where the perturbation of $`H_\nu ^U`$ is relevant, and a massless ($`\nu _o`$) for
$$2\pi u_\nu (K_\nu 1)>|U_\nu |,$$
(B3)
where the perturbation is irrelevant.
### 1 Massive regime ($`\nu _m`$)
When the cosine term is relevant, the conformal symmetry is lost and the elementary excitations become massive particles. To compute the correlation functions we can approximate the cosine term by
$$H_m^\nu =\frac{m_\nu ^2}{2}_0^L𝑑x(\mathrm{\Phi }_\nu \mathrm{\Phi }_\nu )^2,$$
(B4)
where the mass can be obtained from the exact solution of the sine-Gordon equation. For small $`U_\nu `$ one has
$$m_\nu =\left(\frac{4K_\nu |U_\nu |a}{\pi u_\nu }\right)^{\frac{1}{22K_\nu }}a^1.$$
(B5)
This Hamiltonian describes the fluctuations $`\delta \mathrm{\Phi }_\nu `$ of $`\mathrm{\Phi }_\nu `$ about its mean value $`\mathrm{\Phi }_\nu =0`$. For such a system the Green’s function $`T_\tau \mathrm{\Phi }_\nu (𝐫_\nu )\mathrm{\Phi }_\nu (\mathrm{𝟎})_{\nu _m}`$ of the Laplace operator defined on the domain $`A_\nu =[0<u_\nu \tau <u_\nu \beta ,0<x<L]`$ is given by
$$G^{\nu _m}(r_\nu )=\frac{K_\nu }{2}\mathrm{K}_0\left[m_\nu (r_\nu +a)\right].$$
(B6)
$`K_0`$ is the Bessel function of zero order.
### 2 Massless regime ($`\nu _o`$)
In this regime, the bare parameters are renormalized up to the fixed point values $`u_\nu u_\nu ^{}`$, $`K_\nu K_\nu ^{}`$, and $`U_\nu 0`$ without changing the basic Luttinger Liquid behavior of the unperturbed part $`H_0^\nu `$. For this model, the Green’s functions for the unperturbed part regularized for large distances by $`R_\nu `$ and for short distances by the lattice constant $`a`$ can be expressed as
$$G^{\nu _o}(r_\nu )=\frac{K_\nu }{2}\mathrm{ln}\left[R_\nu /(r_\nu +a)\right]$$
(B7)
or as the following limit
$$G^{\nu _o}(r_\nu )=\underset{m_\nu 0}{lim}G^{\nu _m}(r_\nu )$$
(B8)
### 3 Correlation functions
Typical spin-spin correlation functions of the original fermions defined in (63) are combinations of exponentials of $`\mathrm{\Phi }_\nu `$. For a Gaussian model these functions can be expressed in terms of the Green’s functions (B6) or (B7) depending on the phase $`\nu _i`$,
$`\mathrm{exp}[i\gamma _1\mathrm{\Phi }_\nu (1)]\mathrm{}\mathrm{exp}[i\gamma _N\mathrm{\Phi }_\nu (N)]_{\nu _i}=`$ (B9)
$`e^{_{n>m}^N\gamma _n\gamma _mG^{\nu _i}(r_\nu ^{nm})}e^{\frac{1}{2}_n^N\gamma _n^2G^{\nu _i}(r_\nu ^{nn})}.`$ (B10)
## C Local states in the strong interaction limit
The projected Hamiltonian (1) is expressed as
$$\widehat{P}H\widehat{P}=\widehat{P}(L+K)\widehat{P},$$
(C1)
where $`L`$ denotes the local system, whereas $`K`$ includes all possible hopping terms of $`H`$. The eigenstates of $`L`$ are given by
$$|n_1,n_2,\mathrm{},n_j,\mathrm{},n_N=\underset{j=1}{\overset{N}{}}|n_j,$$
(C2)
where $`n_j`$ labels the local states $`n`$ on site $`j`$. The local states and energies are shown in Table V. For simplicity we use the short notation
$$|0_1,0_2,\mathrm{},n_j,\mathrm{},m_i,\mathrm{}0_{N1},0_N|n_j,m_i$$
(C3)
(local ground state configurations are labeled by $`|0_j`$). The energy of such a state is given by
$$E_{n,m}=(N2)ϵ_0+ϵ_n+ϵ_m.$$
(C4)
Now, we can expand the projection operator $`\widehat{P}`$ onto the unperturbed eigenstates of $`L`$. Here for the half-filled case, we are only interested in the projection $`\widehat{P}`$ onto the state $`|0=_{j=1}^N|0_j`$ with the energy $`E_0=Nϵ_0`$, thus we get
$$\widehat{P}=\underset{i}{}\widehat{P}^{(i)},$$
(C5)
where the first orders are given by
$`\widehat{P}^{(0)}`$ $`=`$ $`\widehat{P}_P`$ (C6)
$`\widehat{P}^{(1)}`$ $`=`$ $`\widehat{P}_{PQ}+\widehat{P}_{QP}`$ (C7)
$`\widehat{P}^{(2)}`$ $`=`$ $`\widehat{P}_{PQQ}+\widehat{P}_{QPQ}+\widehat{P}_{QQP}`$ (C8)
$``$ $`(\widehat{P}_{PPQ^2}+\widehat{P}_{PQ^2P}+\widehat{P}_{Q^2PP})`$ (C9)
with
$`\widehat{P}_P`$ $`=\widehat{P}_0`$ (C10)
$`\widehat{P}_{PQ}`$ $`=\widehat{P}_0K\widehat{Q}_0{\displaystyle \frac{1}{E_0L}}\widehat{Q}_0`$ (C11)
$`\widehat{P}_{QP}`$ $`=\widehat{Q}_0{\displaystyle \frac{1}{E_0L}}\widehat{Q}_0K\widehat{P}_0`$ (C12)
$`\widehat{P}_{PQQ}`$ $`=\widehat{P}_0K\widehat{Q}_0{\displaystyle \frac{1}{E_0L}}\widehat{Q}_0K\widehat{Q}_0{\displaystyle \frac{1}{E_0L}}\widehat{Q}_0`$ (C13)
$`\widehat{P}_{QPQ}`$ $`=\widehat{Q}_0{\displaystyle \frac{1}{E_0L}}\widehat{Q}_0K\widehat{P}_0K\widehat{Q}_0{\displaystyle \frac{1}{E_0L}}\widehat{Q}_0`$ (C14)
$`\widehat{P}_{QQP}`$ $`=\widehat{Q}_0{\displaystyle \frac{1}{E_0L}}\widehat{Q}_0K\widehat{Q}_0{\displaystyle \frac{1}{E_0L}}\widehat{Q}_0K\widehat{P}_0`$ (C15)
$`\widehat{P}_{PPQ^2}`$ $`=\widehat{P}_0K\widehat{P}_0K\widehat{Q}_0{\displaystyle \frac{1}{(E_0L)^2}}\widehat{Q}_0`$ (C16)
$`\widehat{P}_{PQ^2P}`$ $`=\widehat{P}_0K\widehat{Q}_0{\displaystyle \frac{1}{(E_0L)^2}}\widehat{Q}_0K\widehat{P}_0`$ (C17)
$`\widehat{P}_{Q^2PP}`$ $`=\widehat{Q}_0{\displaystyle \frac{1}{(E_0L)^2}}\widehat{Q}_0K\widehat{P}_0K\widehat{P}_0`$ (C18)
The projection operator $`\widehat{Q}_0`$ denotes $`1\widehat{P}_0`$.
We can compute the projected electron-nuclear interaction term $`\widehat{P}H_N\widehat{P}`$. The projection affects only the electronic spins, and we have to evaluate projected local s orbital spin operators such as $`\widehat{P}𝐒_{\eta j}\widehat{P}`$. For example, the second order processes (see Fig. 6(a) and 7(c)) are given by
$$𝐒_{\eta j}^{(2)}=(\widehat{P}^{(0)}+\widehat{P}^{(1)})𝐒_{Aj}(\widehat{P}^{(0)}+\widehat{P}^{(1)})$$
(C19)
Introducing (C9) and (C18) in (C19) only
$$𝐒_{Aj}^{(2)}=\widehat{P}_{PQ}𝐒_{Aj}\widehat{P}_{QP}$$
(C20)
will contribute, because in a half-filled system the first hopping process brings the system out of the ground state $`|0`$ and the second one brings it back to a possible ground state configuration $`|0`$. Here only the projections
$`\widehat{P}_{PQ}`$ $`=`$ $`{\displaystyle \frac{t_{Aa}}{E_0E_{\overline{4},13}}}|0\overline{4}_j,13_{j\pm 1}|`$ (C21)
$`\widehat{P}_{QP}`$ $`=`$ $`{\displaystyle \frac{t_{Aa}}{E_0E_{\overline{4},13}}}|\overline{4}_j,13_{j\pm 1}0|`$ (C22)
could generate finite matrix elements for second order contributions to the Cu-4s spin $`𝐒_{Aj}`$. For the O-3s spin $`𝐒_{Bj}`$ it will be
$`\widehat{P}_{PQ}`$ $`=`$ $`{\displaystyle \frac{t_{Ba}}{E_0E_{12}}}|012_j|`$ (C23)
$`\widehat{P}_{QP}`$ $`=`$ $`{\displaystyle \frac{t_{Ba}}{E_0E_{12}}}|12_j0|`$ (C24)
or
$`\widehat{P}_{PQ}`$ $`=`$ $`{\displaystyle \frac{t_{Ba}}{E_0E_{\overline{2},13}}}|0\overline{2}_j,13_{j+1}|`$ (C25)
$`\widehat{P}_{QP}`$ $`=`$ $`{\displaystyle \frac{t_{Ba}}{E_0E_{\overline{2},13}}}|\overline{2}_j,13_{j+1}0|.`$ (C26)
Finally, the projected spins $`𝐒_{\eta j}^{(2)}`$ are given by
$`𝐒_{Aj}^{(2)}`$ $`=`$ $`|\lambda _{Aa}|^2\widehat{P}_0\left(𝐒_{a,j1}+𝐒_{a,j+1}\right)\widehat{P}_0`$ (C27)
$`𝐒_{Bj}^{(2)}`$ $`=`$ $`|\lambda _{Ba}|^2\widehat{P}_0\left(𝐒_{a,j}+𝐒_{a,j+1}\right)\widehat{P}_0.`$ (C28)
with
$`\lambda _{Aa}`$ $`=`$ $`{\displaystyle \frac{t_{Aa}}{ϵ_Aϵ_aU_a}}`$ (C29)
$`\lambda _{Ba}`$ $`=`$ $`\pm {\displaystyle \frac{t_{Ba}}{ϵ_Bϵ_aU_a}}.`$ (C30)
Higher order contributions could be computed using the same procedure as for the above examples.
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# The Contribution of Field OB Stars to the Ionization of the Diffuse Ionized Gas in M331footnote 11footnote 1Observations made with the Burrell Schmidt Telescope of the Warner and Swasey Observatory, Case Western Reserve University. ,2footnote 22footnote 2Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute. STScI is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract NAS 5-26555.
## 1 Introduction
An important component of the interstellar medium (ISM) in spiral galaxies is the diffuse ionized gas (DIG, also called WIM for warm ionized medium). DIG is a warm ($``$ 8000 K), diffuse (n<sub>e</sub>=0.1$``$0.2 cm<sup>-3</sup>) layer of ionized hydrogen which permeates the disks of spiral galaxies (see Walterbos & Braun 1996 for a review). In the Milky Way, where it is often referred to as the Reynolds layer, DIG accounts for almost all of the ionized gas mass, contributes about 30% of the local HI column, and fills at least 20% of the volume of the galaxy (Reynolds, 1991). DIG in several edge-on galaxies has been studied through H$`\alpha `$ imaging (Rand, Kulkarni, & Hester, 1990; Dettmar, 1990; Pildis, Bregman, & Schombert, 1994; Rand, 1996; Hoopes, Walterbos, & Rand, 1999), and the extent and brightness of halo DIG has been found to vary dramatically. In face-on galaxies, however, the most surprising discovery has been the similarity of the DIG content among galaxies spanning a large range of characteristics. The DIG components of many non-edge-on galaxies have been studied: M31 (Walterbos & Braun, 1994), M33 (Hester & Kulkarni, 1990; Greenawalt, 1998), M51 and M81 (Greenawalt et al., 1998; Wang, Heckman, & Lehnert, 1997, 1999), NGC 253 and NGC 300 (Hoopes, Walterbos, & Greenawalt, 1996), NGC 247 and NGC 7793 (Ferguson et al., 1996), and M101 (Wang, Heckman, & Lehnert, 1997, 1999; Greenawalt, 1998). Despite the range of parameters seen in these galaxies, the contribution of the DIG to the total H$`\alpha `$ luminosity (the diffuse fraction) is consistently between 30 and 50%. This is essentially a comparison of the DIG luminosity and the HII region luminosity, so the constant ratio establishes a link between the DIG and massive star formation.
Photoionization by massive stars and shock ionization by supernovae have both been proposed as ionization sources of the DIG. All other sources have been found to lack by far the required energy (Reynolds, 1984), although Webber (1998) finds that cosmic ray ionizations may be a somewhat more significant source than previously thought. The consistent diffuse fraction of 30 to 50% of the total H$`\alpha `$ luminosity implies that a similar fraction of the ionizing photons produced in a galaxy is required to power the DIG, if photoionization is responsible for the DIG. The diffuse fraction for M33 is 40% (Greenawalt, 1998), which leads to a minimum energy requirement (assuming all ionizing photons have $`\lambda `$=912 Å) for the DIG of 1.8$`\times `$10<sup>41</sup> erg s<sup>-1</sup>, or 1.7$`\times `$10<sup>-4</sup> erg s<sup>-1</sup> cm<sup>-2</sup> of disk, using a radius of 6 kpc for M33. This can be compared to 1.0$`\times `$10<sup>42</sup> erg s<sup>-1</sup> and 1$`\times `$10<sup>-4</sup> erg s<sup>-1</sup> cm<sup>-2</sup> for the Milky Way, using a radius of 15 kpc, which just equals the amount of energy provided by supernovae (Reynolds, 1984). In M33 the supernova rate is 1/360 years (Gordon, et al., 1998), which leads to 0.9$`\times `$10<sup>41</sup> erg s<sup>-1</sup>, using the canonical 10<sup>51</sup> erg for each supernova. This is a factor of 2 too low to account for the DIG, even assuming an unrealistic 100% efficiency in converting the energy into ionized gas. Also, the H$`\alpha `$ luminosity used here has not been corrected for internal extinction, which would make the power requirement even higher. Although some DIG may be shock ionized, supernovae cannot be the source of energy responsible for the bulk of the ionization of the DIG in M33 or the Milky Way, through shock ionization or any other source which taps supernovae energy, such as turbulent mixing layers (Slavin, Shull, & Begelman, 1993).
This leaves OB stars as the most likely source. Photoionization is consistent for the most part with spectroscopic observations of the DIG. Emission from \[SII\] 6716, 6731Å and \[NII\] 6548, 6584Å is enhanced relative to H$`\alpha `$ in the DIG when compared to HII regions, and \[OIII\] 4959, 5007Å is fainter relative to H$`\beta `$ (Rand, 1997; Greenawalt, Walterbos, & Braun, 1997; Haffner, Reynolds, & Tufte, 1999; Reynolds, 1985). This spectrum has been reproduced by photoionization models (Domgörgen & Mathis, 1994) using a dilute radiation field in a diffuse medium. Low limits on HeI 5876Å/H$`\alpha `$ in the Milky Way (Reynolds & Tufte, 1995) have yet to be explained, and in NGC 891 the HeI/H$`\alpha `$ ratio, while higher than in the Milky Way, is still too low to be consistent with other line ratios (Rand, 1997). In M31, however, the HeI emission is consistent with the Domgorgen & Mathis models (Greenawalt, Walterbos, & Braun, 1997), at least for the brightest DIG, and the HeI/H$`\alpha `$ ratio in several irregular galaxies is also consistent with photoionization (Martin & Kennicutt, 1997). A more challenging problem has been posed by high \[OIII\]/H$`\beta `$ ratios seen in the DIG of several galaxies (Wang, Heckman, & Lehnert, 1997; Rand, 1998), which still have not been reproduced by pure photoionization models. Nevertheless, the bulk of the evidence suggests that OB stars are the dominant ionization source for the DIG.
There is still uncertainty, however, regarding the location of the OB stars responsible for ionizing the DIG. Are the ionizing photons leaking out of density-bounded HII regions, or are there enough field OB stars to provide the ionizing photons? Either option is a deviation from commonly held views of the ISM. If field stars are the dominant source of ionizing photons, this requires a substantial population of massive stars outside of HII regions. Do these stars form in the field, do they manage to drift out of the cloud in which they formed, or is the gas swept away by the supernovae of even more massive stars? On the other hand, if the photons which are ionizing the DIG are leaking out of HII regions, an appreciable population of density-bounded HII regions is required. HII regions are usually treated as being radiation-bounded, which makes the determination of star formation rates from H$`\alpha `$ luminosities straightforward. Does such a population of leaky HII regions exist?
Evidence for field OB stars does exist (Torres-Peimbert, Lazcano-Araujo, & Peimbert, 1974; Humphreys & Sandage, 1980; Garmany, Conti, & Chiosi, 1982; Massey et al., 1995a). Massey et al. (1995a) defined field stars as those further from an HII region than a star could travel in the lifetime of an OB star ($`<`$ 10<sup>7</sup> years), a few tens of parsecs at most. Using ground-based UBV photometry, Patel & Wilson (1995) investigated the distribution of OB stars in M33. They concluded that 50% of the ionizing stars were found outside of HII regions. Such a study is difficult using ground-based optical data, where both crowding and the color degeneracy of hot stars pose serious problems. A systematic search for field OB stars is needed, and space based photometry in the ultraviolet can provide a strong test of the optical results.
The alternative to field OB stars is a significant population of density-bounded HII regions. Individual HII regions which are density-bounded do exist, such as the Orion nebula (Rubin et al., 1991) and the starburst in NGC 4214 (Leitherer et al., 1996b). Oey & Kennicutt (1997) compared the H$`\alpha `$ luminosities of several HII regions in the LMC with the predicted ionizing fluxes of the stars within. They found 2 out of 14 HII regions to be substantially leaking ionizing photons, and 3 others that were slightly leaky. The observed absence of Lyman continuum radiation leaking from spirals (Leitherer et al., 1995) suggests that such radiation is contained within HII regions, although it is conceivable that HII regions are just leaky enough to ionize the DIG, but not so much that an appreciable number of ionizing photons escapes the entire galaxy (Wang, Heckman, & Lehnert, 1997). There is evidence from the H$`\alpha `$ luminosity functions of HII regions in spirals that high luminosity HII regions may be density-bounded (Beckman et al., 2000).
Here we focus on identifying the stars responsible for ionizing the DIG. In order to isolate massive stars we turn to the far-ultraviolet (FUV), where OB stars dominate the stellar emission. The local group spiral M33 is an ideal galaxy to study the ionization of the DIG, due to its low inclination, vigorous star formation rate, and distance of only 0.84 Mpc (Freedman, Wilson, & Madore, 1991). In section 2 we discuss the data used to carry out this project, which include H$`\alpha `$ imaging of M33, FUV imaging using the Ultraviolet Imaging Telescope (UIT), and optical and FUV stellar photometry from the Wide Field Planetary Camera 2 (WFPC2) on Hubble Space Telescope (HST). In section 3 we briefly review the properties of the DIG in M33. In section 4 we discuss the FUV emission on large scales from the UIT image and its relation to the H$`\alpha `$ emission from both DIG and HII regions. In section 5 we present the results of the stellar photometry in HII regions and DIG. In section 6 we compare the FUV information from HST and UIT. Section 7 contains a discussion of the results and conclusions.
## 2 The Data
### 2.1 The H$`\alpha `$ Mosaic
Table 1 contains a summary of the observations. The ground-based data consist of H$`\alpha `$ and off-band continuum images of M33 obtained with the 0.6 meter Burrell-Schmidt telescope at Kitt Peak National Observatory in November 1995 (Greenawalt, 1998). The 2048<sup>2</sup> detector provided a field of view of 1.15$`\times `$1.15 degrees<sup>2</sup> at 2.028<sup>′′</sup>/pixel. The images were offset using the shift-and-stare technique to minimize flat-field irregularities. The final mosaic covers an area of about 1.75$`\times `$1.75 degrees<sup>2</sup>. We reduced the images using standard procedures in IRAF. Twilight flats from all nights of the observing run were co-added to create a “super flat” which was used to correct for gain variations. We calibrated the continuum image using the published R magnitude and the known shape and transmission of the continuum filter. The H$`\alpha `$ image was calibrated relative to the continuum, again using the shape and transmission of the line filters. The continuum image was scaled to the H$`\alpha `$ image using foreground stars and subtracted. The total H$`\alpha `$ flux we observed is 3.6 $`\times `$ 10<sup>-10</sup> ergs s<sup>-1</sup> cm<sup>-2</sup>. Spectroscopy of DIG and HII regions in M33 (Hoopes & Walterbos 2000, in preparation) shows that \[NII\] 6584Å emission is about 20% as strong as H$`\alpha `$ in both the DIG and HII regions. The lack of \[NII\] enhancement in the DIG and the lower ratios in both environments than seen in other galaxies are probably due to a lower Nitrogen abundance in M33 (Vílchez et al., 1988). The result is that at most only 5% of the light detected in the H$`\alpha `$ filter can arise from \[NII\] emission, and the amount of contamination is the same for both HII regions and DIG. The continuum-subtracted H$`\alpha `$ image is shown in figure 1.
### 2.2 The UIT Image
We obtained an archival UIT image of M33 taken on the Astro-1 space shuttle mission (Stecher et al., 1992, 1997). We used the FUV image, a 424 second exposure taken through the B1 filter ($`\lambda `$=1520Å, $`\mathrm{\Delta }\lambda `$=354Å). The image was originally produced on film and later digitized. The resolution of the image is about 3<sup>′′</sup>and the field of view is roughly 40. The H$`\alpha `$ image was convolved to similar resolution and put on the same grid. UIT was calibrated in the laboratory before launch and also in flight by comparing UIT images with previous UV space missions. The calibration uncertainty is about 15% (Stecher et al., 1992). The UIT image can be found in Landsman et al. (1992).
### 2.3 WFPC2 Images and Stellar Photometry
Lastly, we used archival HST WFPC2 images of regions in M33, from project GO6038. These images were originally obtained to study hot stars and star clusters in M33 (Chandar, Bianchi, & Ford, 1999). Figure 1 shows that the pointings are restricted to the inner $`10^{}`$ radius of M33, but all of the pointings contain HII regions and DIG. The pointings near the center cover regions that are more crowded than those covered in the pointings at larger radii. Spiral arms and inter-arm regions are both sampled. The only obvious bias introduced by the locations of the images is that they do not include the brightest HII regions in M33. The filters and exposure times are F170W (2$`\times `$900 s), F336W (2$`\times `$900 s), F439W(600 s) and F555W(160 s). The two separate exposures in the UV and FUV filters were combined to remove cosmic rays. The images were corrected for the $``$ 4% charge transfer efficiency problem of the WFPC2 by multiplying with a ramp image as described by Holtzman et al. (1995). The F170W and F336W images were also corrected for degradation of throughput after decontamination, again using the formula in Holtzman et al. (1995).
Photometry was performed on the HST images using the DAOPHOT package. For each frame, a PSF model was created using bright, isolated stars. The empirical PSF model was found using an iterative process. A constant PSF model was first computed and used to subtract the neighbors of the PSF stars. Then a linearly varying PSF model was computed on the image with the PSF neighbors subtracted, and then used to remove the PSF neighbors from the original image. Finally a quadratically varying model was computed from the PSF stars with the neighbors subtracted. If there were few PSF stars (less than about 15, which was usually the case in the F170W images) a constant psf model was used. It was also created iteratively, first making a model PSF, then subtracting any neighbors of the PSF stars and revising the PSF based on the subtracted image. We compiled a list of stars detected in all four bands. The magnitudes were put on the STMAG photometric system, described by Holtzman et al. (1995), using the most recently determined calibration for WFPC2. The calibration uncertainty is 1$``$2% in the optical filters, and around 10% in the F170 filter.
## 3 DIG in M33
The H$`\alpha `$ mosaic of M33 is discussed in detail in Greenawalt (1998). We briefly repeat the relevant results here. The diffuse fraction of M33 was found to be 40%, in good agreement with all other non-edge-on spirals studied so far. The Schmidt image reveals the DIG to be a complex network of filaments and arcs, superimposed upon a fainter, more diffuse background. As in other spirals, DIG is concentrated in regions of high star formation, and in the inner disk, but it is much more extended than the star forming regions, filling the disk with faint H$`\alpha `$ emission. Narrow-band imaging of the DIG in M33 (Greenawalt, 1998) reveals that it shows the same spectral signature of enhanced \[SII\] seen in other galaxies.
## 4 H$`\alpha `$UIT Analysis
The FUV/H$`\alpha `$ ratio in the DIG can constrain the source of ionization. For example, if the measured ratio in the DIG did not correspond to any reasonable average spectral type or population of ionizing stars, we would conclude that the ionizing photons in the DIG are not produced locally. To perform this test we measured the FUV flux (L<sub>1520</sub>) in the UIT image and the H$`\alpha `$ flux in the Schmidt image. For HII regions, we defined an aperture and a background annulus on the H$`\alpha `$ image, and then used the same aperture and annulus on the FUV image. For DIG, we measured the fluxes in 500 pc $`\times `$ 500 pc square apertures (in the plane of the sky), with the HII regions masked out. The HII region mask was defined on the H$`\alpha `$ image and applied to both the H$`\alpha `$ and FUV images. This is slightly different from the approach taken in Hoopes & Walterbos (1997). In that paper the fluxes were measured in smaller apertures, leading to the possibility that in HII regions an aperture may miss the concentrated cluster of OB stars while still containing H$`\alpha `$ emission, leading to very high H$`\alpha `$ /L<sub>1520</sub> ratios and thereby increasing the scatter in the distribution of ratios. Also in that paper no background was subtracted from the HII regions, contrary to the approach taken here. The background subtraction typically affects the H$`\alpha `$ flux of the HII regions by only 1$``$2%, but in the FUV image the contrast between HII regions and the field is not as high as in H$`\alpha `$, and the background subtraction can reduce the FUV flux by as much as 50%, and even more for faint HII regions.
We computed the number of ionizing photons (N<sub>Lyc</sub>) from the H$`\alpha `$ luminosity, assuming ionization equilibrium and case B recombination (Osterbrock, 1989) and neglecting absorption of ionizing photons by dust. The histograms in figure 2 show the resulting distribution of N<sub>Lyc</sub>/L<sub>1520</sub> ratios for the two environments. The histograms are shown sideways for comparison with models (see below). The observed ratios are shown; they have not been corrected for foreground or internal extinction. The N<sub>Lyc</sub>/L<sub>1520</sub> ratio in DIG is lower than that in HII regions. This fits a scenario in which the ionizing stars in the DIG are of later type than those in HII regions, since the later type stars would produce a lower ratio of ionizing to non-ionizing UV photons.
In figure 2 we compare the N<sub>Lyc</sub>/L<sub>1520</sub> ratios with models of evolving stellar populations. We used the Starburst99 evolutionary model to compute the ratio of ionizing photons to FUV flux (Leitherer et al., 1999). The model used a Salpeter (1995) initial mass function (IMF) with an upper mass cutoff of 120 $`M_{}`$. LMC-like (Z=0.008) metallicity models were calculated. The metallicity of M33 in the central region is about solar, but there is a steep gradient towards lower metallicity with distance from the center (Vílchez et al., 1988), so that overall M33 is subsolar in metallicity. Single burst and steady state models are presented. The reddening bars shown indicate the direction that the models would move if reddened by the given amount. Both the LMC (Howarth, 1983) and Galactic (Cardelli, Clayton, & Mathis, 1989) extinction laws are given, and the case of a foreground screen and a uniform mixture of dust and gas are shown for each extinction law. Massey et al. (1995b) found an average color excess of E(B-V)=0.16 in several fields of M33. The Galactic foreground produces a color excess of E(B-V)=0.03$`\pm `$0.02 at this location (McClure & Racine, 1969), and the rest is internal to M33. The extinction in the Massey et al. (1995b)) fields ranges from E(B-V)=0.09 to 0.33. We also show reddened versions of each model. The steady state model, which represents the population expected outside of HII regions, was reddened by E(B-V)=0.1, while the single burst model, which represents the population in an HII region, was reddened by E(B-V)=0.2, all using the LMC extinction law. The dashed line is reddened assuming a uniform mixture, and the dotted line is reddened assuming a foreground screen. Each reddened model also includes E(B-V)=0.03 of foreground Galactic reddening, using the Galactic extinction law.
The DIG distribution is consistent with an older burst population, or with a steady state, constant star formation model. The steady state model is probably a more accurate description of the stellar population in the DIG, but it is only an approximation. Figure 2 also shows the predictions of the steady state model using steeper IMF slopes than the usual Salpeter slope of $`\alpha `$=$``$2.35. Although the predicted ratios are below the observed ratio, with a reasonable amount of extinction the prediction for $`\alpha `$=$``$3.0 would match the observed DIG values well. This is interesting in light of the fact that Massey et al. (1995a) found a steeper IMF slope for OB stars in the field in the LMC and SMC. We will explore this issue further in section 6. The HII regions resemble a young burst population. The model values are dependent on the assumed IMF, and in small HII regions (which are the most numerous) that have low numbers of massive stars, the ionizing flux of a single O star may dominate the light from the UV emitting B stars. Figure 3 shows the same histograms along with the expected ratios from ionizing stars, taken from the CoStar stellar models (Schaerer & de Koter, 1997). We show models using low metallicity (Z=0.004) and solar metallicity (Z=0.020), to cover the extreme values. Again the DIG ratio can be explained by later type ionizing stars.
The N<sub>Lyc</sub>/L<sub>1520</sub> ratio in the DIG can be explained with a local ionization model, where the ionizing population in the DIG is older than that in HII regions. However, the UIT analysis is far from conclusive. The ionizing photons which produce the H$`\alpha `$ emission seen in the DIG may be produced elsewhere and then leak into the DIG, as in the leaky HII regions explanation for the ionization of the DIG. In this case the models in figure 2 and 3 would not apply. To test this more conclusively, we need to find out whether the ionizing photons are produced locally. We investigate this further in the next section.
If the difference in the FUV/H$`\alpha `$ ratio between DIG and HII regions were due to higher extinction in the HII regions, at least 0.9 magnitudes A<sub>V</sub> in excess of that in the DIG would be necessary to explain the result using a foreground screen model, and much higher for a uniform mixture. Greenawalt, Walterbos, & Braun (1997) found, on average, a 0.3 magnitude difference in extinction between HII regions and DIG in M31. Our stellar photometry (see the next section) indicates that there is on average little difference in the amount of extinction in the two environments. However, this is for only a very small part of the disk, so to be consistent with Greenawalt, Walterbos, & Braun (1997) we reddened the burst models more than the steady state models in figures 2 and 3.
## 5 HST Photometry
### 5.1 Spectral Classification of OB Stars
The most conclusive test of whether field stars can ionize the DIG would be to find the spectral types of all the OB stars outside of HII regions, add up the number of ionizing photons contributed by each star, and compare the total to that required by the H$`\alpha `$ luminosity of the DIG. Massey (1985) and Massey et al. (1995b) have shown that optical photometry is unable to distinguish between the different spectral types of ionizing stars, a distinction that must be made if accurate Lyman continuum luminosities are to be determined. This is because the peak of the Planck spectrum lies blueward of these filters for stars of 30,000 K or hotter, so optical filters sample the Rayleigh-Jeans portion of the spectrum which has a constant slope independent of temperature. Spectroscopy is the most reliable way to determine spectral types of massive stars, but obtaining spectroscopy of a large sample of field stars in M33 would require a prohibitive amount of telescope time. However, FUV photometry combined with optical photometry can be used to estimate spectral types of ionizing stars with reasonable accuracy. The F170W is closer to the peak of the Planck spectrum for temperatures above 30,000 K, providing an advantage over optical photometry alone.
The archival HST images cover DIG of varying brightness (see overlay of figure 1). However, they do not overlap any very luminous HII regions, only moderate and low-luminosity HII regions. This is a limitation of the data, and it prevents us from investigating whether the brightest HII regions are density-bounded. It should be noted that the HST images cover a much smaller area of the disk of M33 than the UIT image, so the focus of the analysis shifts here from a global to a local perspective. The HST pointings are all located in the inner disk of M33, within 10 from the nucleus, where the metallicity is close to solar (Vílchez et al., 1988). For this reason we assume solar metallicity when employing stellar models in the analysis of the WFPC2 images.
We put the H$`\alpha `$ image of M33 on the same grid as the WFPC2 frames, and then classified the stars in the WFPC2 images as being either in an HII region or in the field based on the H$`\alpha `$ surface brightness. An isophotal cut was used, but the level was adjusted to satisfactorily isolate the two environments based on their morphology. In the inner disk the cut level is as high as 200 pc cm<sup>-6</sup>, while in the outer disk it is as low as 60 pc cm<sup>-6</sup>. Individual HII regions were then picked from the masked image by eye. The rest of the image was counted as DIG. Only the Wide Field images were used, as there were no complete HII regions in the Planetary Camera images, and they do not cover a large enough area of DIG to analyze.
Holtzman et al. (1995) caution that the response curves of the main photometric filters used with WFPC2 (F336W, F439W, F555W), are sufficiently different from the groundbased analogs (UBV), that any reddening correction must be applied to the photometry in the HST filter system. Thus we derived reddening relations for the HST filter system used in this project, rather than transform the photometry to the groundbased system. We used the model stellar spectra of Lejeune, Cuisinier, & Buser (1997), and reddened them according to both the Galactic extinction law (Cardelli, Clayton, & Mathis, 1989) and the LMC law (Howarth, 1983). We used the Lejeune, Cuisinier, & Buser (1997) model spectra rather than the CoStar models used below because the Lejeune et al. models extend to spectral types later than B0, allowing more of the unevolved main sequence to be used in determining the extinction. The reddened spectra were multiplied by the HST filter response curves to obtain the magnitude. The derived relations agree with those found in Holtzman et al. (1995), although the F170W relation was not given. We found that the LMC extinction law worked more satisfactorily than the Galactic law for reddening all bands to match the observed colors. As before, the reddening includes a correction for E(B-V)=0.03 foreground reddening using the Galactic extinction law, and any excess reddening using the LMC extinction law.
Figure 4 shows representative color-magnitude diagrams (CMDs) for HII regions and DIG regions in M33. Isochrones of ages 0 and 7$`\times 10^6`$ years are also shown. The isochrones were generated using the Geneva stellar evolution models (Schaerer et al., 1993). In order to put the isochrones in the STMAG system, for each stellar mass at each age we found the closest match in T<sub>eff</sub> and $`logg`$ in the model stellar spectra of Lejeune, Cuisinier, & Buser (1997). The spectrum was then multiplied by the WFPC2 filter response curves to derive the observed magnitude. The isochrones are shown for reference only, they are not the best fit to the CMD. The purpose of these isochrones was to determine the extinction for each region. The amount of extinction needed to match the observed main sequence with the isochrone is taken as the average extinction for that region. The same was done for the remaining stars not in HII regions, which are field stars in the DIG. Using a single value for extinction in either environment is obviously an over-simplification, but it is the best solution for the available data. The average extinction found in HII regions was E(B-V)=0.14, and in DIG it was E(B-V)=0.13. We estimate that the extinction can be determined with this method with a 1-$`\sigma `$ accuracy of E(B-V)= $`\pm `$0.02.
We assigned spectral types based on the CoStar models of Schaerer & de Koter (1997). They provide spectral energy distributions for stars of mass 20, 25, 40, 60, 85, and 120 M at various ages in the Leitherer et al. (1996a) database. We used the solar metallicity (Z=0.020) models to match the metallicity where the HST pointings are located. We ran the model spectra through the HST filters, and linearly interpolated the magnitudes and other stellar properties (e.g. temperature, surface gravity, and ionizing photon luminosity) to the ages and masses between those given, to create a complete grid with 1 M mass resolution and 10<sup>5</sup> year time resolution. We then searched for the nearest match to the four magnitudes measured for each star. Thus for each star we go directly from the measured photometry to ionizing luminosity. The uncertainty in the ionizing luminosity is determined by the uncertainty in the photometry, ignoring uncertainty in the accuracy of the models, which is difficult to quantify.
We then inspected the color magnitude diagrams and the positions of the stars that were classified, noting any UV-bright stars that were not classified. In order to classify these stars, we adjusted the tolerance for the F555W magnitude, and if necessary for the F439W magnitude as well. In most cases the tolerance was increased to 2-$`\sigma `$, and in a few cases, 3-$`\sigma `$. A small number of stars were in crowded regions and obviously had mismatched F555W and F439W magnitudes (crowding usually does not affect the UV magnitudes since fewer stars show up in the UV images), so we used only F170W and F336W magnitudes for classification. Most stars were matched within 1-$`\sigma `$ of the four measured magnitudes.
We compared our photometric spectral types to previous spectroscopic spectral types to assess the accuracy of the photometric technique. We looked for overlapping stars in Massey et al. (1995b) and Massey et al. (1996), finding only two stars that overlap, both in Massey et al. (1996). The reason is that our technique only assigns spectral types to B0 stars and earlier, so while many B stars overlapped, we did not assign spectral types to them, so they could not be used for comparison. The properties that the photometric technique assigned to the two stars are shown in table 2, along with the properties of the spectral type assigned by Massey et al. (1996). These properties are taken from Vacca, Garmany, & Shull (1996). The properties that the photometric technique assigned to both stars match reasonably well with the properties expected for stars of the correct spectral type. It should be noted that UIT104 is a Wolf-Rayet star, and might not be expected to match the properties of an O9Ia star. Nevertheless, the agreement gives us reason to believe that the photometrically determined spectral types agree with those determined spectroscopically, to within about a spectral type.
### 5.2 Comparison of Predicted and Observed H$`\alpha `$ Luminosity
For each HII region we summed the predicted ionizing luminosity of all the stars included in the region. The predicted ionizing luminosity in the DIG was the cumulative luminosity of stars not included in an HII region. Each chip of each WFPC2 pointing was considered separately. We converted the predicted ionizing luminosity to a predicted H$`\alpha `$ luminosity L, assuming that all ionizing photons result in an ionization (i.e. no dust absorption) and case B recombination. The predicted L was then compared to the observed L. We measured the observed L from the H$`\alpha `$ image (figure 1). The boundaries of the HII regions were determined on the masked image, and a background correction for the surrounding DIG was subtracted. HII regions which were not completely contained within the HST pointing were not considered. The L for DIG was the sum of all the luminosity not in an HII region. The luminosities were corrected for extinction using the estimates of E(B-V) from the color-magnitude diagrams.
In figure 5 we compare the predicted L, based on the ionizing luminosity of the stars present, with the observed L from the H$`\alpha `$ image for both HII regions and DIG. The errors reflect a combination of uncertainty in photometry and in extinction. The effect of photometric errors on the ionizing luminosity was calculated by determining for each star the model with the highest N<sub>Lyc</sub> and the model with the lowest N<sub>Lyc</sub> that still matched within the photometric errors. When the best match luminosities were added for all stars in the region, the range found for each star was added in quadrature. The effects of uncertainty in extinction were quantified by varying the correction by E(B-V)=$`\pm `$0.02, and again calculating the total ionizing luminosity of the stars in a region. These two uncertainties were added together in quadrature to produce the error bars shown in figure 5.
Qualitatively, we find that the measured L for both HII regions and DIG agree well with the predicted luminosities. There is considerable scatter, but given the uncertainties in stellar models, extinction in the UV, and spectral classification without the use of spectra, the agreement is encouraging. The observed L is corrected for extinction, but the predicted L may be an overestimate since we ignore the possibility of ionizing photons being absorbed by dust. To correct for this we would have to know how much of the ionizing luminosity of a star is absorbed by dust. McKee & Williams (1997) estimate the fraction of dust absorbed in Galactic HII regions to be about 25%, perhaps slightly lower for smaller HII regions such as those we consider here.
Figure 5 shows that while the mean ratio is close to unity, some HII regions are under-predicted or over predicted. This is an indication that predictions for individual regions are not reliable, but that the average for all the regions is better. Thus a comparison of the average values for DIG and HII regions is more illuminating than individual regions. The average ratio of predicted to required L are given in table 3. For HII regions the predictions are consistent with the observed L. For DIG regions there are not enough ionizing photons emitted by field stars to account for all of the observed H$`\alpha `$ emission, but there are enough to provide a significant fraction of the ionization, about 40%.
To test the idea that FUV information is necessary to estimate the spectral types of ionizing stars, we re-classified the ionizing stars using only the F336W, F439W, and F555W magnitudes. The results are given in table 3. The predicted L is higher, and most HII regions are overpredicted, leading to a large excess of ionizing photons. If the FUV information had been ignored in the initial classification, many non-ionizing stars would have been classified as ionizing stars, leading to a further excess. The FUV information is crucial to accurately predict ionizing fluxes, and optical information alone will lead to a severe over prediction of ionizing fluxes of OB stars. This test underscores the difficulty of drawing conclusions about OB stars based on optical photometry (Patel & Wilson 1995, see also O’Dell, Hodge, & Kennicutt 1999). This may account for the discrepancy in the fraction of OB stars found outside of HII regions. We identified 116 ionizing stars in the five WFPC fields. Of these, 27% (31 stars) were in the DIG, lower than the $``$50% found by Patel & Wilson (1995). However, it is important to remember that we can only investigate a small fraction of the disk of M33 with the present HST data.
### 5.3 Characteristics of the HII Regions
Oey & Kennicutt (1997) found that 2 out of 14 HII regions in the LMC are leaking a sizeable number of ionizing photons. These two HII regions both have a ring-like morphology, leading to the question of whether HII morphology is correlated the amount of leakage. This is an important question in light of suggestions that ionizing photons escaping from superbubbles may be an important ionization source (Dove, Shull, & Ferarra, 2000). Softening of the spectrum by re-radiation of ionizing photons through chimney walls may explain the low HeI measurements (Norman, 1991). We classified the HII regions in the HST field based on their morphology in the H$`\alpha `$ image, in order to look for trends in the amount of leakage with morphology. Figure 6 shows this comparison. The HII regions are either compact, meaning center brightened and obviously discrete, or diffuse, which includes rings, filaments, and faint tenuous HII regions (the ring-like HII regions are shown as double open circles). It appears that compact and diffuse HII regions are equally likely to be density-bounded.
In figure 6 we also see that there is no significant trend for leakiness with L. It is very important to remember that the brightest HII regions in M33 were not imaged in the HST data, and only a few in our sample are more luminous than 10<sup>38</sup> erg s<sup>-1</sup> (corrected for internal extinction). Thus we cannot address the issue of whether the most luminous HII regions are density-bounded (Beckman et al., 2000). It is still interesting, however, that the two brightest HII regions in M33, NGC 604 (González Delgado & Pérez, 2000) and NGC 595 (Malamuth, Waller, & Parker, 1996), appear to be radiation-bounded. In a future paper we will investigate the brightest HII regions in M33.
## 6 UIT - HST Comparison
The WFPC2 images resolve stars above a limiting magnitude, which in the F170W filter corresponds to about a B0 star. The UIT image measures all of the FUV surface brightness of the entire population of UV emitting stars. The UIT image sensitivity extends to 50-100 Myr old populations (O’Connell, 1997), which corresponds to B and early A type stars, so a comparison of the UV flux measured by HST and UIT thus compares the stars of spectral type B0 and earlier to the entire OBA population. This is an age dependent ratio which can be modeled, and may be used to further characterize the stellar population in both HII regions and DIG. The attractive feature of this comparison is that it is largely independent of extinction, as both filters cover a similar range of wavelengths (however extinction must still be incorporated, see below).
We have performed this comparison for the HII regions and DIG on the HST fields. The FUV flux from the UIT image was measured using the same apertures used to measure the H$`\alpha `$ flux. We then summed up the FUV flux of the stars brighter than a magnitude limit in the same regions in the HST images. In order to keep a consistent limiting magnitude, we reddened the F170W magnitude of a B0 star (M<sub>F170W</sub>=16.92) using the extinction found from the CMDs and used throughout this analysis, and summed the flux from all stars brighter than this magnitude. This is the point where extinction becomes important in this comparison, the determination of the cutoff magnitude in the HST images. The relationship between flux and magnitude in the STMAG system is given by
$$M_{F170W}=2.5log(F_\lambda )21.1$$
(1)
(Holtzman et al., 1995). The flux was determined from the magnitudes without correcting for extinction, as the UIT flux is also uncorrected. The extinction is only used to determine the cutoff magnitude which corresponds to a B0 star.
A histogram of the ratios is shown in figure 7. The histogram shows a difference between the two populations, with DIG having a lower ratio of HST to UIT flux. This is expected if the DIG is ionized by an older population, with fewer stars earlier than B0 relative the number of stars later than B0. Also shown on the plot are models of the evolution of this ratio through time, constructed using Starburst99 (Leitherer et al., 1999). The UIT flux was modeled as an evolving population with a mass range from 1 to 120 M, a Salpeter (1995) IMF, and solar metallicity. In order to simulate the populations seen in the HST FUV images, we needed to include flux only from stars with M$``$20M, as this corresponds to the limiting spectral type of B0 in the HST images. To do this we ran the models again, except that mass range was 20 to 120 M, and the total mass was scaled down by 0.2 (so that the total mass, including stars with M$``$20M, was 10<sup>6</sup>M, as it was for the UIT simulation). This approach is simplistic, but provides a first-order model for comparison with the observations. As expected, the ratio of HST to UIT flux in the single burst model generally decreases as time passes and the most massive stars die. The observed ratios for HII regions match well with the predicted values for single burst populations. An older burst model reproduces the DIG ratios reasonably well.
We also ran steady state models as described above, but the resulting ratio did not agree well with the observed ratio for the DIG, contrary to the analysis in figure 2. The disagreement may indicate that if the steady state model is an accurate description of the stellar population in the DIG, the parameters used for the HII region stars may not apply to the field stars. Specifically, Massey et al. (1995a) derived the IMF for field stars in the LMC and SMC, and found a slope ranging from $`\alpha =4.7`$ to $`5.1`$ (in the notation where the Salpeter IMF slope $`\alpha `$ is $`2.35`$). We ran models using IMF slopes for massive stars (M$``$20M) of $`\alpha `$=$`3.5,3`$, and the standard Salpeter slope of $``$2.35, while keeping the slope for less massive stars at $`\alpha `$=$`2.35`$. The observed HST/UIT flux ratio agrees much better with the models using a steeper slope of $`\alpha `$=$`3`$ than the Salpeter slope. This supports the conclusions of Massey et al. (1995a) regarding the IMF slope for field OB stars. The observed ratios suggest that the IMF slope for field OB stars in M33 is not as steep as that in the LMC and SMC, but is still steeper than the IMF in HII regions.
It should be noted that the Starburst99 models are intended for starburst regions, i.e. regions with a large number of stars. Most of the HII regions we are investigating do not fit into this category, so the model cannot be expected to accurately predict the observed ratio. The main problem is that with small numbers of massive stars, the IMF is not well sampled. This may be the reason behind the spread in the observed distribution of HST/UIT ratios, and also might explain why some HII regions have ratios too high to be explained by the models. However, the predicted ratio agrees with the majority of HII regions, which suggests that these deviations average out for a large number of HII regions. Also note that the two filters are not identical, and this is probably the reason some of the HII regions have HST/UIT $`>`$1.
## 7 Discussion
Our most important result is that the OB stars outside of HII regions can account for 40% of the ionization of the DIG in M33. There are several points to keep in mind. One is the scale we are investigating, which is constrained by the size of the Wide Field chips of WFPC2, 325 pc across each chip. This means we are implicitly assuming that stars (field stars in particular) do not have an influence at distances greater than this, or that the number of photons escaping the image is balanced by the number of photons coming in from outside the image. In a uniform medium with n=0.2 cm<sup>-3</sup>, which is the density found for the DIG in the Galaxy, the Strömgren sphere of an O3V star has a diameter of about 770 pc (Osterbrock 1989, using the stellar ionizing flux in Vacca et al. 1996). However there are few of these stars outside of HII regions. One chip can completely contain the Strömgren sphere of an O8V star, and a B0V star ionizes a region about 170pc across. The density in the DIG in the inner regions of M33 may be higher than the density measured in the solar neighborhood, which would tend to make the ionized regions smaller. We cannot account for neighboring HII regions or ionizing stars outside of the WFPC2 field of view.
Another point to consider is the many uncertainties dealt with in this analysis. When working in the UV extinction is always a prime concern, and small changes in the adopted extinction, or in the extinction law, can cause large changes in the ionizing luminosity. When dealing with individual stars in external galaxies, crowding can often be a concern. The stars visible in the F170W filter are rarely crowded, but of course crowding may be present but not detectable, as in the case of binaries. Crowding may also be a more severe problem for OB associations in HII regions than for the more sparsely distributed stars in the DIG. Another source of uncertainty stems from the stellar atmosphere models which we use to assign Lyman continuum luminosities. The stellar models give excellent agreement for the non-ionizing spectrum of massive stars, but the ionizing spectrum is difficult to test without making assumptions. We have tried to be as conservative as possible in every aspect of this analysis. Given these uncertainties, it is remarkable that we find such good agreement between the predicted and observed L. This gives us confidence in the stellar models, and it also suggests that any mistakes we are making regarding the extinction and spectral classification are relatively small. It also suggests that we are detecting most, if not all, of the ionizing stars.
Without placing trust in the absolute accuracy of the predicted ionizing fluxes, it is possible to draw conclusions by comparing our predictions for DIG and HII regions. This comparison relies on the assumption that any errors in the prediction for the DIG will also be made for HII regions. The ratio of N<sub>Lyc</sub>/L is lower in the DIG than in HII regions, indicating that some leakage is necessary to explain the DIG. However, the difference is not so great that field stars can be neglected as a source of ionization. The N<sub>Lyc</sub>/L in the DIG is about 37% of the ratio in HII regions. If we normalize the N<sub>Lyc</sub>/L to 100% in HII regions, this relative approach would suggest that field stars can ionize at least 37% of the DIG, well within the uncertainty of our absolute determination.
In the regions covered in the WFPC2 pointings, the fraction of the total H$`\alpha `$ luminosity that comes from the DIG is 40%, which is also the diffuse fraction of the entire galaxy. In the five fields analyzed here, on average the field stars can account for 40%$`\pm `$12% of the ionization of the DIG (it is a potentially confusing coincidence that these two numbers are the same). This implies that only 30% of the ionizing photons emitted in HII regions need to escape to account for the remaining DIG, or put differently, the predicted L in HII regions should be 143% of the observed L in HII regions. The average ratio of predicted to observed L for HII regions is 107%$`\pm `$26%, so the amount of excess ionizing photons is not enough to explain all of the remaining DIG, with the maximum within the uncertainty being 133%. Simply adding up the observed and predicted L, we find that, within the uncertainty, 98% of the total observed L (DIG+HII) can be explained by field stars plus leakage (taking the maximum predicted N<sub>Lyc</sub> within the uncertainty). Also keep in mind that other processes may play some role in ionizing the DIG, such as turbulent mixing layers (Slavin, Shull, & Begelman, 1993) and shock ionization. Most likely a combination of these processes, plus photons leaking from HII regions, ionize the rest of the DIG.
There is also the possibility that a fraction of the ionizing photons are absorbed by dust. McKee & Williams (1997) found that about 25% of the ionizing photons emitted by stars within HII regions are absorbed by dust. If this is the case, we then predict only 80% of the ionizing photons necessary to explain the L of the HII regions. Since there is not even enough to explain the L from HII regions, we cannot then explain the remaining DIG ionization with leaky HII regions. However, since we know that the HII regions are ionized by the stars within, the discrepancy might be explained by a systematic error in determining the spectral types and ionizing luminosities from the photometry. In this case we could then scale the predicted HII region luminosity up to 100%, and scale up the predicted DIG luminosity accordingly to perform a relative comparison. We could not address the question of whether any excess ionizing luminosity from stars in HII regions exists. If an equal fraction of the ionizing photons emitted by field OB stars is absorbed by dust, the results of this paper would not be greatly affected.
Taking the maximum predicted N<sub>Lyc</sub> within the uncertainty, we find that there are no ionizing photons from stars left over to escape the galaxy altogether. Of course this depends on the contribution of other ionization sources to the DIG. If 20% of the DIG is ionized by another source, there can be as much as 4% of the stellar ionizing photons left over to escape the galaxy, within the error bars. The uncertainty in these numbers is such that we cannot place much weight on this limit. However, it is in agreement with Leitherer et al. (1995), who found that less than 3% of the ionizing photons escape from four starburst galaxies observed with HUT, although Hurwitz, Jelinsky, & Dixon (1997) suggest that the fraction could be as high as 57% for one of the galaxies by allowing for absorption by undetected components of the ISM. Deharveng et al. (1997) also find very little leakage, less than 1%. We must also remember that we are restricted to relatively small regions (about 650pc across), and are really covering only a small fraction of the disk of M33. We have know way of knowing whether excess photons emitted by stars in the image actually escape the galaxy, or simply ionize gas outside of the image. Similarly, there may be photons from outside of the image ionizing gas in the image.
Note that leaky HII regions can provide less than 60% of the ionization of the DIG in the regions of M33 studied here. This is in disagreement with the results of Oey & Kennicutt (1997), who found that the Lyman continuum escaping from HII regions could very likely account for all of the DIG emission in the LMC. This may point to a difference between the DIG in irregular galaxies and spirals. Martin & Kennicutt (1997) found that Helium is completely ionized in the DIG in several irregular galaxies, indicating ionization by very massive stars, type O7 or hotter. They concluded that density-bounded HII regions are the dominant source of ionization in those galaxies. A difference in the ionization source might explain the lower HeI 5876 Å/H$`\alpha `$ seen in some spirals (Reynolds & Tufte, 1995; Rand, 1997).
The field OB stars in M33 appear to have a steeper IMF than OB stars in HII regions, confirming the results of Massey et al. (1995a) for field OB stars in the LMC and SMC. The slope which best fits our observations is not as steep as that found in the LMC and SMC, but it is still different from that found in HII regions. The difference may imply a different formation mechanism for field stars, but does not necessarily mean that they form in the field. Field OB stars may have formed in an HII region and then drifted out of the dense cloud, or they may be the remnants of an OB association after the surrounding gas has dissipated through the actions of SNe and stellar winds. In either case, lower mass OB stars are more likely to become field stars, as they have longer lives and thus can drift farther, or live long enough to outlast the HII region. The result would be a steeper mass function for field OB stars, because the youngest and most massive stars are still associated with HII regions. In this case the varying IMF is used to simulate a steeper mass function for stars in the field, and it is not necessary that they have a different initial mass function. The steeper IMF merely represents the fact that the most massive HII region stars have a smaller chance of becoming field stars due to their short lifetimes. However, Massey et al. (1995a) carefully corrected for stars which may have drifted out of HII regions, and found stars as massive as 85 M in the field. It is difficult to explain how a star this massive became a field star if it did not form in the field.
We thank the anonymous referee for a very careful reading of the manuscript and comments which improved the presentation of our results. We are grateful to Bruce Greenawalt for obtaining the H$`\alpha `$ data, and to Richard Rand for the use of his narrow band filters. The archival UIT image was obtained through the NASA Data Archive and Distribution Service. We would like to acknowledge the UIT project for making their data available. The availability of the Starburst99 models of C. Leitherer and collaborators is greatly appreciated. This research benefitted from helpful discussions with Jon Holtzman, Salman Hameed, Bruce Greenawalt, David Thilker, Nichole King, and Vanessa Galarza. Support for this work was provided by NASA through grant number AR07645.01-96A from the Space Telescope Science Institute, by NASA grant NAG5-2426, a Cottrell Scholar Award from Research Corporation, and by the NSF through grant AST-9617014. CGH was supported by a grant from the New Mexico Space Grant Consortium.
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# Electronic transport through ballistic chaotic cavities: reflection symmetry, direct processes, and symmetry breaking
## I Introduction
The problem of chaotic wave scattering is of great interest in various branches of physics, like optics, nuclear, mesoscopic and microwave physics.
The study of quantum-mechanical scattering problems whose classical dynamics is chaotic has been further motivated by recent experiments on quantum electronic transport in microstructures consisting of a cavity connected to leads . We know that symmetries have very interesting effects on the properties of the electric conductance in mesoscopic systems: time-reversal and spin-rotational symmetries , as well as spatial-reflection symmetries have been studied in the literature.
The problem of electronic transport through asymmetric (AS) chaotic cavities is addressed in detail in Ref. in an independent-electron approximation. In that reference, the possibility of direct processes due to the presence of short paths is accounted for by specifying the average, or optical, $`S`$ matrix $`S`$, within an information-theoretic approach. The statistical distribution for the $`S`$ matrix is known as Poisson’s kernel, in which $`S`$ enters as a parameter. When $`S=0`$, i.e. in the absence of direct processes, the statistical distribution reduces to the invariant measure for the appropriate universality class.
Microstructures with reflection symmetry and a chaotic classical dynamics are studied in Refs. and . The analysis is performed in the absence of direct processes, so that the statistical distribution of the $`S`$ matrix is the invariant measure for the universality class in question and the relevant spatial symmetry: the latter is a symmetry of the full system under consideration, i.e. the cavity plus the two leads that connect the cavity to the outside.
One purpose of the present paper is to extend the study of Refs. and to include the presence of direct processes. We consider two-dimensional systems with spinless particles and concentrate on left-right (LR) symmetry only, i.e. symmetry under reflection through an axis perpendicular to the current. We also restrict the analysis to time-reversal-invariant (TRI) problems. One particular way of inducing direct reflections is by adding potential barriers between the symmetrically positioned waveguides and the cavity. If the two barriers are equal, the system is fully LR symmetric; if the barriers are different, we have a LR-symmetric cavity coupled asymmetrically to the outside: using the jargon of nuclear physicists , we shall refer to this type of symmetry breaking as “external mixing”, with an obvious meaning. An interesting question, amenable to experimental observation, is that of the interplay between the symmetry of the cavity and external mixing in the statistical distribution of the conductance of such a structure: the study of that interplay is the second main purpose of this paper. From an experimental point of view microwave cavities and acoustic systems might represent good candidates to study these questions.
That interplay may also be there and have interesting effects when $`S=0`$, as in the case of a LR-symmetric cavity coupled to the outside by two waveguides free of potential barriers but asymetrically located. This problem can be addressed from the point of view of the systems described in the previous paragraph in the following way. One may think of a LR-symmetric cavity coupled to the outside by four waveguides, also placed symmetrically. We can break the symmetry by providing the two waveguides on the right-hand side of the cavity, say, with identical barriers. The desired problem is then approached in the limit of impenetrable barriers.
The paper is organized as follows. In order to make the paper reasonably self-contained, we summarize in the next section a number of concepts that we shall be using throughout the paper, like the invariant measure and Poisson’s kernel for $`S`$ matrices and their application to chaotic scattering in AS cavities, and the invariant measure for LR-symmetric systems. Sec. III deals with the problem of fully LR-symmetric systems in the presence of direct processes. The distribution of the conductance is calculated for the particular case of one open channel in each lead and a diagonal optical matrix (implying direct reflections), and contrasted with the one obtained for an AS chaotic cavity and the same optical matrix $`S`$. Different barriers added to the two waveguides of an otherwise fully LR-symmetric system with no direct processes give rise to direct reflections and external mixing: the problem is studied in Sec. IV. Again, the conductance distribution is computed for the one-channel case and contrasted with the one obtained for an AS chaotic cavity with the same optical matrix $`S`$. The problem of external mixing in a LR-symmetric cavity with asymetrically positioned leads and $`S=0`$ is addressed in Sec. V. The conductance distribution is calculated and compared with the one arising from the invariant measure in the AS case. Finally, for the sake of completeness, we include a number of appendices where some of the results mentioned in the text are derived.
## II The S matrix and its statistical distribution
### A The scattering problem in the absence of spatial symmetries
A single-electron scattering problem can be described by the scattering matrix $`S`$, which in the stationary case relates the outgoing-wave to the incoming-wave amplitudes . For a ballistic cavity connected to two leads, each with $`N`$ transverse propagating modes (see Fig. 1), the $`S`$ matrix is $`n=2N`$-dimensional and has the structure
$$S=\left(\begin{array}{cc}r& t^{}\\ t& r^{}\end{array}\right),$$
(1)
where $`r`$, $`r^{}`$ are the $`N\times N`$ reflection matrices (for incidence from either lead) and $`t`$, $`t^{}`$ the corresponding transmission matrices.
From the $`S`$ matrix we can construct the total transmission coefficient, or spinless dimensionless conductance
$$T=\text{tr}\left(tt^{}\right),$$
(2)
which is proportional to the conductance of the cavity,
$$G=(2e^2/h)T,$$
(3)
the factor 2 arising from the two spin directions.
In Dyson’s scheme there are three basic symmetry classes. In the absence of any symmetry, the only restriction on $`S`$ is unitarity, i.e.
$$SS^{}=I,$$
(4)
resulting from the physical requirement of flux conservation. This is the “unitary” case, also designated as $`\beta =2`$. For orthogonal symmetry, or $`\beta =1`$, $`S`$ is symmetric, i.e.
$$S=S^T,$$
(5)
because one has either time-reversal invariance (TRI) and integral spin, or TRI, half-integral spin and rotational symmetry. In the “symplectic” case ($`\beta =4`$), $`S`$ is self-dual because of TRI with half-integral spin and no rotational symmetry. From now on we consider the scattering problem of “spinless” electrons, so that the case $`\beta =4`$ will not be touched upon.
A convenient parametrization of the $`S`$ matrix is the polar representation
$$S=\left(\begin{array}{cc}v_1& \hfill 0\\ 0& \hfill v_2\end{array}\right)\left(\begin{array}{cc}\sqrt{1\tau }& \hfill \sqrt{\tau }\\ \sqrt{\tau }& \hfill \sqrt{1\tau }\end{array}\right)\left(\begin{array}{cc}v_3& \hfill 0\\ 0& \hfill v_4\end{array}\right),$$
(6)
where $`\tau `$ stands for the $`N`$-dimensional diagonal matrix of eigenvalues $`\tau _a`$ ($`a=1,\mathrm{},N`$) of the Hermitian matrix $`tt^{}`$; $`v_i`$ ($`i=1,\mathrm{},4`$) are arbitrary $`N\times N`$ unitary matrices for $`\beta =2`$, with the restriction $`v_3=v_1^T`$, $`v_4=v_2^T`$ for $`\beta =1`$.
#### 1 The invariant measure
When the classical dynamics of the system is chaotic, a statistical analysis of the quantum-mechanical problem is called for. That analysis is performed in terms of “ensembles” of physical systems, described mathematically by an ensemble of $`S`$ matrices, endowed with a probability measure. The starting point to such an analysis is the concept of invariant measure, which is a precise formulation of the intuitive notion of equal a priori probabilities in the space of scattering matrices.
The invariant measure, to be designated as $`d\mu ^{(\beta )}(S)`$, is invariant under the symmetry operation which is relevant to the universality class under consideration , i.e.
$$d\mu ^{(\beta )}(S)=d\mu ^{(\beta )}(U_0SV_0).$$
(7)
Here, $`U_0`$, $`V_0`$ are arbitrary but fixed unitary matrices in the unitary case, while $`V_0=U_0^T`$ in the orthogonal one. Eq. (7) defines the Circular (Orthogonal, Unitary) Ensembles (COE, CUE), for $`\beta =1,2`$, respectively.
#### 2 Chaotic scattering by AS cavities
The information-theoretic approach of Refs. leads to the probability distribution known as Poisson’s kernel
$$dP_S^{(\beta )}(S)=\frac{[det(ISS^{})]^{(\beta n+2\beta )/2}}{|det(ISS^{})|^{\beta n+2\beta }}d\mu ^{(\beta )}(S),$$
(8)
where the invariant measure is assumed normalized, i.e.
$$𝑑\mu ^{(\beta )}(S)=1.$$
(9)
Here, $`n=2N`$ is the dimensionality of the $`S`$ matrix and $`S`$ is the averaged, or optical, $`S`$ matrix, which describes the promt response arising from direct processes.
In the absence of direct processes, $`S=0`$ and Poisson’s measure (8) reduces to the invariant measure for the universality class in question. In terms of the polar representation, the invariant measure can be written as
$$d\mu ^{(\beta )}\left(S\right)=p^{(\beta )}\left(\{\tau \}\right)\underset{a}{}d\tau _a\underset{i}{}d\mu \left(v_i\right).$$
(10)
Here, the joint probability density of {$`\tau `$} is
$$p^{(\beta )}\left(\{\tau \}\right)=C_\beta \underset{a<b}{}|\tau _a\tau _b|^\beta \underset{c}{}\tau _c^{(\beta 2)/2},$$
(11)
$`C_\beta `$ being a normalization constant and $`d\mu (v_i)`$ denoting the invariant measure on the unitary group $`U(N)`$ for matrices $`v_i`$.
For $`S0`$, a useful construction of Poisson’s ensemble is given in Refs. . Consider the system shown in Fig. 1: it consists of a cavity described by the $`n`$-dimensional scattering matrix $`S_0,`$ connected to two leads by the tunnel barriers described by the $`n\times n`$ scattering matrices
$`S_1`$ $`=`$ $`\left(\begin{array}{cc}r_1& \hfill t_1^{}\\ t_1& \hfill r_1^{}\end{array}\right),`$ (14)
$`S_2`$ $`=`$ $`\left(\begin{array}{cc}r_2& \hfill t_2^{}\\ t_2& \hfill r_2^{}\end{array}\right),`$ (17)
respectively. We bunch the two leads into a “superlead” and construct the $`2n\times 2n`$ scattering matrix $`S_b`$
$$S_b=\left(\begin{array}{cc}r_b& \hfill t_b^{}\\ t_b& \hfill r_b^{}\end{array}\right)=\left(\begin{array}{cccc}r_1\hfill & 0\hfill & t_1^{}\hfill & 0\hfill \\ 0\hfill & r_2^{}\hfill & 0\hfill & t_2\hfill \\ t_1\hfill & 0\hfill & r_1^{}\hfill & 0\hfill \\ 0\hfill & t_2^{}\hfill & 0\hfill & r_2\hfill \end{array}\right).$$
(18)
Here, the various blocks ($`r_b`$, etc.) are $`n`$-dimensional. The scattering matrix $`S_0`$ for the cavity can be written in terms of the scattering matrix $`S`$ for the full system $`\left\{\mathrm{cavity}+\mathrm{barriers}\right\}`$ as
$$S_0=\frac{1}{t_b^{}}\left(Sr_b\right)\frac{1}{Ir_b^{}S}t_b^{}.$$
(19)
One can prove that between the invariant measures for $`S_0`$ and for $`S`$ we have the Jacobian
$$d\mu ^{(\beta )}(S_0)=\frac{[det(ISS^{})]^{(\beta n+2\beta )/2}}{|det(ISS^{})|^{\beta n+2\beta }}d\mu ^{(\beta )}(S).$$
(20)
Now, if the matrix $`S_0`$ for the cavity is distributed according to the invariant measure, i.e. $`d\mu ^{(\beta )}(S_0)`$, the distribution of the transformed $`S`$ satisfies
$$dP(S)=d\mu ^{(\beta )}(S_0)$$
(21)
and we obtain Eq. (8), the optical $`S`$ being given by the $`n`$-dimensional matrix
$$S=r_b=\left(\begin{array}{cc}r_1\hfill & 0\hfill \\ 0\hfill & r_2^{}\hfill \end{array}\right).$$
(22)
The $`N=1`$, $`\beta =1`$ case. The $`T`$ distribution. We now consider the distribution of the $`S`$ matrix for the system shown in Fig. 1 for the case $`N=1`$ and $`\beta =1`$. The matrices $`S_0`$ of the ballistic cavity, $`S_1`$ and $`S_2`$ of the two tunnel barriers and $`S`$ \[related through Eq. (19)\] are $`2\times 2`$ and have the structure (1) with $`t^{}=t`$. In the polar representation (6) we have three independent parameters $`\tau `$, $`\varphi `$, $`\psi `$, where we have written $`v_1=e^{i\varphi }`$, $`v_2=e^{i\psi }`$. The range of variation of these parameters is taken to be
$$\begin{array}{cc}\hfill \tau & [0,1],\hfill \\ \hfill \varphi ,\psi & [0,2\pi ].\hfill \end{array}$$
(23)
In terms of them, $`S`$ can be written as
$$S=\left(\begin{array}{cc}r& \hfill t\\ t& \hfill r^{}\end{array}\right)=\left[\begin{array}{cc}\sqrt{1\tau }e^{2i\varphi }& \hfill \sqrt{\tau }e^{i(\varphi +\psi )}\\ \sqrt{\tau }e^{i(\varphi +\psi )}& \hfill \sqrt{1\tau }e^{2i\psi }\end{array}\right].$$
(24)
and the invariant measure of Eqs. (10) and (11) as
$$d\mu (S)=\frac{d\tau }{2\sqrt{\tau }}\frac{d\varphi }{2\pi }\frac{d\psi }{2\pi }.$$
(25)
The distribution of $`S`$ is given by Poisson’s kernel, with the optical $`S`$ matrix
$$S=r_b=\left(\begin{array}{cc}r_1& \hfill 0\\ 0& \hfill r_2^{}\end{array}\right).$$
(26)
Substituting $`S`$ in Eq. (8), Poisson’s measure can be written as
$$dP_{r_1,r_2^{}}(S)=\frac{\left[\left(1|r_1|^2\right)\left(1|r_2^{}|^2\right)\right]^{3/2}}{|(1rr_1^{})(1r^{}r{}_{2}{}^{}^{})t^2r_{1}^{}{}_{}{}^{}r_{2}^{}{}_{}{}^{}|^3}d\mu (S).$$
(27)
By definition, the resulting distribution of the transmission coefficient $`T`$ can be expressed as the integral
$$w_{r_1,r_2^{}}(T)=\delta \left(T\tau \right)𝑑P_{r_1,r_2^{}}(S).$$
(28)
For this distribution, Ref. gives the expression
$$\begin{array}{cc}w_{r_1,r_2^{}}(T)=\frac{1}{2\sqrt{T}}\left[\left(1\left|r_1\right|^2\right)\left(1\left|r_2^{}\right|^2\right)\right]^{3/2}\hfill & \\ \times \frac{1}{\left|\left(e^{i\phi }+\left|r_1\right|\sqrt{1T}\right)\left(e^{i\psi }+\left|r_2^{}\right|\sqrt{1T}\right)\left|r_1\right|\left|r_2^{}\right|T\right|^3}_{\phi ,\psi },\hfill & \end{array}$$
(29)
where $`\mathrm{}_{\phi ,\psi }`$ denotes an average over the variables $`\phi `$ and $`\psi `$ over the interval $`[0,2\pi ]`$. When $`r_1=r_2^{}=0`$, the above expression (29) reduces to
$$w_{0,0}(T)=\frac{1}{2\sqrt{T}},$$
(30)
as it should. Fig. 2 ahead (Sect. III) shows with dotted lines the evolution of $`w_{r_1,r_2^{}}(T)`$ for $`r_1=r_2^{}=r`$ with the parameter $`r`$, obtained from Eq. (29) by numerical integration. That distribution tends to $`\delta (T)`$ as $`r0`$.
To further illustrate the physics resulting from the $`S`$-matrix distribution (27) we analyze the special case $`r_1=0`$, so that the right barrier is the only one present. For this case, Eqs. (27), (25) give, for the joint probability distribution of the parameters $`\tau `$, $`\varphi `$, $`\psi `$, the expression
$$dP_{0,r_2^{}}(S)=\frac{\left(1|r_2^{}|^2\right)^{3/2}}{|1\sqrt{1\tau }e^{2i\psi }r{}_{2}{}^{}^{}|^3}\frac{d\tau }{2\sqrt{\tau }}\frac{d\varphi }{2\pi }\frac{d\psi }{2\pi }.$$
(31)
We first notice that the angular variable $`\varphi `$ is uniformly distributed for all $`r_2^{}`$. In this particular case the $`T`$ probability density of Eq. (29) can be integrated analytically, to give
$$w_{0,r_2^{}}(T)=\frac{\left(1\left|r_2^{}\right|^2\right)^{3/2}}{2\sqrt{T}}{}_{2}{}^{}F_{1}^{}(3/2;3/2;1;\left|r_2^{}\right|^2\left(1T\right)),$$
(32)
$`{}_{2}{}^{}F_{1}^{}`$ being a hypergeometric function .
As a check, we consider two limiting situations. Firstly, for $`r_2^{}=0`$ we have a ballistic cavity without prompt response. The probability distribution for $`S`$, $`dP_{0,0}(S)`$ \[see Eq. (31)\], goes back to the invariant measure (25), as it should. Secondly, we obstruct the right lead by making the barrier there a perfect reflector. As a result, $`r_2^{}=1`$ and it can be shown (see Appendix A) that $`dP_{0,r_2^{}}(S)`$ reduces to
$$dP_{0,1}(S)=\delta \left(\tau \right)d\tau \frac{d\varphi }{2\pi }\frac{1}{2}\left[\delta \left(\psi \frac{\pi }{2}\right)+\delta \left(\psi 3\frac{\pi }{2}\right)\right]d\psi ,$$
(33)
where the angles in the arguments of the delta functions are defined modulo $`2\pi `$. We see from the above expression that the distribution of $`\tau `$ is a one-sided delta function at zero, i.e.
$$w(T)=\delta (T),$$
(34)
so that the transmission tends to zero, as expected. Also, the distribution of $`\psi `$ consists of delta functions centered at $`\pi /2`$ and 3$`\pi /2`$, so as to ensure the vanishing of the wave function at the impenetrable barrier. In contrast, as already noted, the variable $`\varphi `$ is uniformly distributed from 0 to $`2\pi `$. In this limiting case we end up with a ballistic cavity connected to just one lead: thus the resulting 1-dimensional $`S`$ matrix $`r=e^{2i\varphi }`$ is distributed according to the invariant measure.
Now we go back to the intermediate case in which $`r_2^{}`$ in Eq. (32) is real and $`1<r_2^{}<0`$. We show in Fig. 5 ahead (Sect. IV) with dotted lines the evolution of the $`T`$ distribution for several values of $`r_2^{}`$, obtained from the analytical result (32).
### B The scattering problem for TRI, LR-symmetric systems
In the presence of additional symmetries, for fixed values for all quantum numbers of the full symmetry group the invariant ensemble is one of the three circular ensembles in Dyson’s scheme. Thus for reflection symmetric systems $`S`$ is block diagonal in a basis of definite parity with respect to reflections, with a circular ensemble in each block .
For a system with TRI and LR symmetry the general form of the $`S`$ matrix is
$$S=\left(\begin{array}{cc}r& \hfill t\\ t& \hfill r\end{array}\right),$$
(35)
with
$`r`$ $`=`$ $`r^T`$ (37)
$`t`$ $`=`$ $`t^T.`$ (38)
All the matrices with the structure (35) can be simultaneoulsy brought to block-diagonal form using the rotation matrix
$$R_0=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}\hfill I_N& \hfill I_N\\ \hfill I_N& \hfill I_N\end{array}\right),$$
(39)
where $`I_N`$ is the $`N`$-dimensional unit matrix. In fact
$$S^{}=R_0SR_0^T=\left[\begin{array}{cc}s^{(+)}& \hfill 0\\ 0& \hfill s^{()}\end{array}\right],$$
(40)
with
$$s^{(\pm )}=r\pm t.$$
(41)
Since $`S`$ is unitary and symmetric, so are $`S^{}`$ and the two $`N\times N`$ matrices $`s^{(\pm )}`$. While $`S`$ has the restricted form (35), $`s^{(\pm )}`$ are the most general $`N\times N`$ unitary and symmetric, i.e. $`\beta =1`$, matrices.
#### 1 The invariant measure
The invariant measure for $`S`$ matrices with the structure (35) was found in Refs. , based on the consideration that two arbitrary unitary symmetric matrices $`s^{(\pm )}`$ can generate the most general unitary $`S`$ matrix with the structure (35). The invariant measure for matrices of the form (35) can be written as
$$d\widehat{\mu }^{(1)}(S)=d\mu ^{(1)}(s^{(+)})d\mu ^{(1)}(s^{()}),$$
(42)
where $`d\mu ^{(1)}(s^{(\pm )})`$ is the invariant measure discussed above for unitary and symmetric matrices ($`\beta =1`$) in the absence of spatial symmetries.
#### 2 Chaotic scattering by systems with full LR symmetry in the absence of direct processes
It has been found that single-electron scattering by classically chaotic cavities with LR symmetry and in the absence of direct processes is well described by the invariant measure discussed above.
The $`N=1`$ case. The $`T`$ distribution. Ref. finds the distribution of the total transmission coefficient $`T`$ for the one-channel case ($`N=1`$) arising from the invariant measure (42) as
$$w(T)=\frac{1}{\pi \sqrt{T(1T)}}.$$
(43)
## III Systems with TRI and full LR symmetry in the presence of direct processes
In this section we study a TRI system with full LR symmetry, just as in Sec. II B, but now admitting the possibility of direct processes. For the systems analyzed in Sec. II B, the average (or optical) $`S`$ matrix $`S`$ vanishes, indicating the absence of a prompt response, whereas now $`S0`$.
The $`S`$ matrix has the structure of Eq. (35), and so does $`S`$, i.e.
$$S=\left(\begin{array}{cc}r& \hfill t\\ t& \hfill r\end{array}\right),$$
(44)
with
$`r`$ $`=`$ $`r^T`$ (46)
$`t`$ $`=`$ $`t^T`$ (47)
being $`N\times N`$ blocks. Both $`S`$ and $`S`$ can be brought to a block-diagonal form by the rotation matrix (39): $`S`$ becomes $`S^{}`$ of Eq. (40) and $`S`$ becomes
$$S^{}=R_0SR_0^T=\left[\begin{array}{cc}s^{(+)}& \hfill 0\\ 0& \hfill s^{()}\end{array}\right].$$
(48)
As we noticed right below Eq. (41), $`s^{(\pm )}`$ are the most general $`N\times N`$ unitary and symmetric matrices; they thus belong to the $`\beta =1`$ universality class. Their distribution is given by two statistically independent Poisson’s kernels of the form (8), with $`s^{(\pm )}`$ as their optical matrices. Denoting by $`d\widehat{P}_S(S)`$ the $`S`$ matrix distribution, we have
$$d\widehat{P}_S(S)=dP_{s^{(+)}}(s^{(+)})dP_{s^{()}}(s^{()}),$$
(49)
where
$`dP_{s^{(\pm )}}(s^{(\pm )})`$ (50)
$`={\displaystyle \frac{\left[det\left(I_Ns^{(\pm )}s^{(\pm )}^{}\right)\right]^{(N+1)/2}}{\left|det\left(I_Ns^{(\pm )}s^{(\pm )}^{}\right)\right|^{N+1}}}d\mu ^{(1)}(s^{(\pm )}).`$ (51)
We can thus write $`d\widehat{P}_S(S)`$ as
$`d\widehat{P}_S(S)={\displaystyle \frac{\left[det\left(I_Ns^{(+)}s^{(+)}^{}\right)\right]^{(N+1)/2}}{\left|det\left(I_Ns^{(+)}s^{(+)}^{}\right)\right|^{N+1}}}`$ (52)
$`\times {\displaystyle \frac{\left[det\left(I_Ns^{()}s^{()}^{}\right)\right]^{(N+1)/2}}{\left|det\left(I_Ns^{()}s^{()}^{}\right)\right|^{N+1}}}d\widehat{\mu }^{(1)}(S)`$ (53)
where $`d\widehat{\mu }^{(1)}(S)`$ is defined in Eq. (42).
The special case of no direct transmission, $`t=0`$, i.e.
$$S=\left(\begin{array}{cc}r& \hfill 0\\ 0& \hfill r\end{array}\right),$$
(54)
can be written as
$`d\widehat{P}_r(S)`$ (55)
$`={\displaystyle \frac{\left[det\left(I_Nrr^{}\right)\right]^{N+1}}{\left|det\left(I_Ns^{(+)}r^{}\right)\right|^{N+1}\left|det\left(I_Ns^{()}r^{}\right)\right|^{N+1}}}`$ (56)
$`\times d\widehat{\mu }^{(1)}(S).`$ (57)
Physically, this case could be realized by fully LR-symmetric structures with no direct processes, to which identical barriers (with the $`\beta =1`$ symmetry) are added in the two leads, each with a reflection matrix \[see Eq. (22)\]
$$r_1=r_2^{}=r.$$
(58)
The situation is illustrated in Fig. 3 ahead but with equals barriers.
### A The $`N=1`$ case. The $`T`$ distribution.
In this case, Eq. (57) reduces to
$$d\widehat{P}_r(S)=\frac{\left[1rr^{}\right]^2}{\left|1s^{(+)}r^{}\right|^2\left|1s^{()}r^{}\right|^2}d\widehat{\mu }^{(1)}(S),$$
(59)
where $`r`$ and $`s^{(\pm )}`$ are now $`1\times 1`$ matrices, i.e. just complex numbers. The distribution of $`T`$ can be obtained from the general expression (28). For $`r`$ real, some of the relevant steps are found in App. B, the final result being
$$w_r(T)=\frac{1}{\pi \sqrt{T\left(1T\right)}}\frac{\left(1+r^2\right)\left(1r^2\right)}{\left(1+r^2\right)^24r^2\left(1T\right)}.$$
(60)
The distribution (60) is plotted in Fig. 2 for several values of $`r`$ and compared, in the same figure, with the distribution corresponding to an AS cavity with the same $`S`$, as given by Eq. (29).
For $`r=0`$, the distribution of Eq. (60) reduces to that of Eq. (43), which is symmetric with respect to $`T=\frac{1}{2}`$, so that $`T`$ and $`R=1T`$ are identically distributed; this feature is lost when $`r0`$, as small $`T`$’s become more probable. As $`r1`$, both distributions shown in the figure (i.e. for LR-symmetric and AS systems) tend to $`\delta (T)`$.
## IV Breaking the reflection symmetry of cavities by direct processes
In this section we study a TRI configuration consisting of a ballistic cavity with LR symmetry and scattering matrix $`S_0`$, connected to two symmetrically positioned waveguides by means of barriers described by $`S_1`$ and $`S_2`$, respectively; in general, the barriers are allowed to be different. This arrangement introduces direct reflections and a breakdown of the reflection symmetry (see Fig. 3). As a result, while the scattering matrix $`S_0`$ of the cavity plus the symmetrically positioned waveguides, but not including the barriers, has the restricted structure (35), (II B), the scattering matrix $`S`$ of the total system including the barriers has the more general form (1), (5). Now, $`S`$ is generated from $`S_0`$ through the inverse of the relation (19); thus, varying $`S_0`$ across its manifold of independent parameters, but keeping the barriers fixed, generates a matrix $`S`$ that varies over a manifold with the same dimensionality. In what follows we restrict ourselves to the one-channel case ($`N=1`$) in each lead. The matrices $`S_0`$ can be expressed in terms of two independent continuous parameters (plus a discrete parameter $`\sigma `$), as in Eq. (69) below, while $`S`$ has the more general form (24); thus there should be an algebraic relation connecting the three continuous parameters $`\tau `$, $`\varphi `$, $`\psi `$ appearing in the latter equation.
We want $`S_0`$ to be distributed according to the invariant measure $`d\widehat{\mu }(S_0)`$. In principle, the transformation between $`S_0`$ and $`S`$ (for fixed $`S_1`$ and $`S_2`$) defines uniquely the resulting statistical distribution of $`S`$, to be called $`d\widehat{P}(S)`$ \[see Eq. (90)\]; for that purpose one could find the Jacobian of the transformation relating $`S`$ to $`S_0`$, both matrices being subject to the restrictions explained in the previous paragraph. In what follows, though, we find it convenient to compute $`d\widehat{P}(S)`$ proceeding along a simpler route, taking advantage of the Jacobian between unrestricted $`S`$ matrices that we already know from Eq. (20). In fact, the measure $`d\widehat{\mu }(S_0)`$ can be first expressed as the measure $`d\mu (S_0)`$ of unrestricted $`S_0`$ matrices of the form of Eq. (65) below, times the appropriate delta functions that provide the required restriction \[see Eq. (66)\] among the three parameters $`\tau _0`$, $`\varphi _0`$, $`\psi _0`$. Next, Eq. (20) expresses $`d\mu (S_0)`$ in terms of $`d\mu (S)`$, the factor in front of $`d\mu (S)`$ in Eq. (20) being the Jacobian of the transformation from unrestricted $`S_0`$ to unrestricted $`S`$ matrices. Finally, the identity (90) gives the required distribution $`d\widehat{P}(S)`$ for the $`S`$ matrices. We proceed to implement this scheme in detail.
The relation between the scattering matrix $`S_0`$ for the cavity and the matrix $`S`$ for the full system is given by Eq. (19), with
$$t_b=t_b^{}=\left(\begin{array}{cc}t_1& \hfill 0\\ 0& \hfill t_2\end{array}\right),$$
(61)
$$r_b=\left(\begin{array}{cc}r_1& \hfill 0\\ 0& \hfill r_2^{}\end{array}\right),r_b^{}=\left(\begin{array}{cc}r_1^{}& \hfill 0\\ 0& \hfill r_2\end{array}\right).$$
(62)
Here all the matrices are 2-dimensional, so that the various entries are just complex numbers. We thus have
$$S_0=t_{b}^{}{}_{}{}^{1}(Sr_b)\frac{1}{I_2r_{b}^{}{}_{}{}^{}S}t_{b}^{}{}_{}{}^{}.$$
(63)
The matrix $`S`$ has the structure (24), while $`S_0`$ has the structure (35), i.e.
$$S_0=\left(\begin{array}{cc}r_0& \hfill t_0\\ t_0& \hfill r_0\end{array}\right).$$
(64)
It will be useful to write $`S_0`$ of Eq. (64) in the polar representation (24) as
$$S_0=\left[\begin{array}{cc}\sqrt{1\tau _0}e^{2i\varphi _0}& \hfill \sqrt{\tau _0}e^{i(\varphi _0+\psi _0)}\\ \sqrt{\tau _0}e^{i(\varphi _0+\psi _0)}& \hfill \sqrt{1\tau _0}e^{2i\psi _0}\end{array}\right].$$
(65)
However, the three parameters $`\tau _0`$, $`\varphi _0`$ and $`\psi _0`$ are not independent. In fact, the structure of $`S_0`$ given in Eq. (64) implies a relation between the two angles $`\varphi _0`$ and $`\psi _0`$, i.e.
$$e^{2i\psi _0}=e^{2i\varphi _0},$$
(66)
or, taking the square root on both sides
$$e^{i\psi _0}=i\sigma e^{i\varphi _0},$$
(67)
where $`\sigma =\pm 1`$. Equivalently
$$\psi _0=\varphi _0+\sigma \frac{\pi }{2},\mathrm{mod}(2\pi ).$$
(68)
The most general form of $`S_0`$ is thus
$$S_0=\left[\begin{array}{cc}\sqrt{1\tau _0}& \hfill i\sigma \sqrt{\tau _0}\\ i\sigma \sqrt{\tau _0}& \hfill \sqrt{1\tau _0}\end{array}\right]e^{2i\varphi _0},$$
(69)
written in terms of the independent parameters $`\tau _0`$, $`\varphi _0`$ and the discrete variable $`\sigma `$, which have the range of variation
$`\tau _0`$ $``$ $`[0,1],`$ (71)
$`\varphi _0`$ $``$ $`[0,2\pi ],`$ (72)
$`\sigma `$ $`=`$ $`\pm 1.`$ (73)
From (39)-(41), the matrix $`S_0`$ can be diagonalized by a $`\pi /4`$ rotation to give
$$S_0^{}=\left[\begin{array}{cc}e^{i\theta _0^{(+)}}& \hfill 0\\ 0& \hfill e^{i\theta _0^{()}}\end{array}\right],$$
(74)
where
$$e^{i\theta _0^{(\pm )}}=r_0\pm t_0=e^{2i\varphi _0\pm i\sigma \beta _0}$$
(75)
and
$$\beta _0=\mathrm{tan}^1\sqrt{\frac{\tau _0}{1\tau _0}},\frac{\pi }{2}\beta _0\frac{\pi }{2}.$$
(76)
With the range of variation (IV) for $`\tau _0`$, $`\varphi _0`$ and $`\sigma `$, $`e^{i\theta _0^{(+)}}`$ and $`e^{i\theta _0^{()}}`$ cover twice the torus defined by the two angles $`\theta _0^{(+)}`$, $`\theta _0^{()}`$.
Eq. (75) is a transformation from the parameters $`\tau _0`$, $`\varphi _0`$ and $`\sigma `$ to the parameters $`\theta _0^{(+)}`$, $`\theta _0^{()}`$, whose Jacobian can be written as
$$\frac{1}{2}\frac{d\theta _0^{(+)}}{2\pi }\frac{d\theta _0^{()}}{2\pi }=\frac{1}{2}\frac{d\tau _0}{\pi \sqrt{\tau _0(1\tau _0)}}\frac{d\varphi _0}{2\pi }.$$
(77)
Both sides of this last equation integrate to $`1`$ if the left-hand side is integrated in the region $`\theta _0^{(+)},\theta _0^{()}[0,2\pi ]`$ and multiplied by $`2`$ to account for the fact that the region is visited twice, and the right-hand side is integrated in the region specified by (IV).
According to Eq. (42), the left-hand side of Eq. (77) represents the invariant measure for $`S_0`$ matrices with LR symmetry. A function $`f(\tau _0,\varphi _0,\sigma )`$ can be translated into a function $`\stackrel{~}{f}(\theta _0^{(+)},\theta _0^{()})`$ using the transformation (75); its average over the $`S_0`$ invariant measure can thus be written as
$`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\theta _0^{(+)}}{2\pi }}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\theta _0^{()}}{2\pi }}\stackrel{~}{f}(\theta _0^{(+)},\theta _0^{()})`$ (78)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma =\pm 1}{}}{\displaystyle _0^1}{\displaystyle \frac{d\tau _0}{\pi \sqrt{\tau _0(1\tau _0)}}}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _0}{2\pi }}f(\tau _0,\varphi _0,\sigma ).`$ (79)
Here, on the left-hand side we integrate over the torus $`\theta _0^{(+)}`$, $`\theta _0^{()}`$ only once. Suppose now that we are given a function $`F(\tau _0,\varphi _0,\psi _0)=`$ $`F^{}(\tau _0,\varphi _0,e^{i\psi _0})`$ of the three parameters appearing in Eq. (65) and we want to compute its average over the above measure. First, we make use of Eq. (67) to eliminate $`\psi _0`$ and write
$`F(\tau _0,\varphi _0,\psi _0)`$ $`=`$ $`F^{}(\tau _0,\varphi _0,e^{i\psi _0})=F^{}(\tau _0,\varphi _0,i\sigma e^{i\varphi _0})`$ (80)
$`=`$ $`f(\tau _0,\varphi _0,\sigma )=\stackrel{~}{f}(\theta _0^{(+)},\theta _0^{()}),`$ (81)
where $`f(\tau _0,\varphi _0,\sigma )`$ and $`\stackrel{~}{f}(\theta _0^{(+)},\theta _0^{()})`$, have the same meaning as in (79) above. The average of this function can thus be written as in (79) and subsequently as
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma =\pm 1}{}}{\displaystyle _0^1}{\displaystyle \frac{d\tau _0}{\pi \sqrt{\tau _0(1\tau _0)}}}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _0}{2\pi }}F^{}(\tau _0,\varphi _0,i\sigma e^{i\varphi _0})`$ (82)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^1}{\displaystyle \frac{d\tau _0}{\pi \sqrt{\tau _0(1\tau _0)}}}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _0}{2\pi }}{\displaystyle _0^{2\pi }}𝑑\psi _0`$ (83)
$`\times `$ $`\left[\delta \left(\psi _0\varphi _0{\displaystyle \frac{\pi }{2}}\right)+\delta \left(\psi _0\varphi _03{\displaystyle \frac{\pi }{2}}\right)\right]F(\tau _0,\varphi _0,\psi _0),`$ (84)
where we have used Eq. (68). Comparing the left hand-side of (79) to the right-hand side of (LABEL:av(F)\_1) we thus write
$`{\displaystyle \frac{d\theta _0^{(+)}}{2\pi }}{\displaystyle \frac{d\theta _0^{()}}{2\pi }}`$ $``$ $`{\displaystyle \frac{2\left[\delta \left(\psi _0\varphi _0\frac{\pi }{2}\right)+\delta \left(\psi _0\varphi _03\frac{\pi }{2}\right)\right]}{\sqrt{1\tau _0}}}`$ (86)
$`\times `$ $`{\displaystyle \frac{d\tau _0}{2\sqrt{\tau _0}}}{\displaystyle \frac{d\varphi _0}{2\pi }}{\displaystyle \frac{d\psi _0}{2\pi }},`$ (87)
where the symbol $``$ indicates that the two measures are equivalent when the left and right-hand sides are used to integrate the functions $`\stackrel{~}{f}(\theta _0^{(+)},\theta _0^{()})`$ and $`F(\tau _0,\varphi _0,\psi _0)`$, respectively, defined above. Obviously, the angles in the argument of the delta functions above are defined modulo $`2\pi `$. As we have already noticed, the left-hand side of Eq. (87) is the invariant measure $`d\widehat{\mu }(S_0)`$ for scattering matrices $`S_0`$ of the form (74), i.e. for a LR-symmetric cavity. On the other hand, Eq. (25) shows that the last line of Eq. (87) is the invariant measure $`d\mu (S_0)`$ for scattering matrices $`S_0`$ of the more general form (65). The relation between the two measures is thus
$$d\widehat{\mu }(S_0)\frac{2\left[\delta \left(\psi _0\varphi _0\frac{\pi }{2}\right)+\delta \left(\psi _0\varphi _03\frac{\pi }{2}\right)\right]}{\sqrt{1\tau _0}}d\mu (S_0).$$
(88)
Here, the delta functions restrict the space of unitary and symmetric matrices to the subspace of matrices of the form (69).
As was explained at the beginning of this section, we now express $`d\mu (S_0)`$ in terms of $`d\mu (S)`$ using Eq. (20). That equation reads, for the present case,
$$d\mu (S_0)=\frac{\left[det\left(I_2r_br_b^{}\right)\right]^{3/2}}{\left|det\left(I_2Sr_b^{}\right)\right|^3}d\mu (S).$$
(89)
We substitute this last equation into Eq. (88) and use Eq. (25) to express $`d\mu (S)`$ in the polar representation. We also note that the measure$`d\widehat{\mu }(S_0)`$ appearing on the left-hand side of Eq. (88), i.e. the differential probability associated with the matrices $`S_0`$ \[having the form (64)\] for the LR-symmetric cavity, must coincide with the differential probability $`d\widehat{P}_{r_b}(S)`$ we are looking for, associated with the transformed matrices $`S`$ \[having the form (24), but with the appropriate restrictions\], i.e.
$$d\widehat{P}_{r_b}(S)=d\widehat{\mu }(S_0).$$
(90)
We thus have
$`d\widehat{P}_{r_b}(S)`$ $``$ $`2{\displaystyle \frac{\delta \left(\psi _0\varphi _0\frac{\pi }{2}\right)+\delta \left(\psi _0\varphi _03\frac{\pi }{2}\right)}{\sqrt{1\tau _0}}}`$ (91)
$`\times `$ $`{\displaystyle \frac{\left[det\left(I_2r_br_b^{}\right)\right]^{3/2}}{\left|det\left(I_2Sr_b^{}\right)\right|^3}}{\displaystyle \frac{d\tau }{2\sqrt{\tau }}}{\displaystyle \frac{d\varphi }{2\pi }}{\displaystyle \frac{d\psi }{2\pi }}.`$ (92)
There remains to express the variables $`\psi _0`$, $`\varphi _0`$, $`\tau _0`$ appearing in the delta-function arguments in terms of $`\psi `$, $`\varphi `$, $`\tau `$. This is done in App. C for the particular case in which barrier 1 is transparent, so that its scattering matrix $`S_1`$ of Eq. (14) is the Pauli matrix $`\sigma _x`$, and barrier 2 is described by Eq. (17) with real matrix elements. The result is
$$d\widehat{P}_{0,r_2^{}}(S)p_{r_2^{}}(\tau ,\varphi ,\psi )d\tau d\varphi d\psi ,$$
(93)
with
$`p_{r_2^{}}(\tau ,\varphi ,\psi )={\displaystyle \frac{\left(1r_{2}^{}{}_{}{}^{2}\right)^{3/2}\left|\sqrt{1\tau }r_2^{}e^{2i\varphi }\right|}{(2\pi )^2\sqrt{\tau }\left|\sqrt{1\tau }\left(1r_{2}^{}{}_{}{}^{2}\right)r_2^{}\tau e^{2i\varphi }\right|^2}}`$ (94)
$`\times \left[\delta \left(\psi \varphi \alpha (\varphi ){\displaystyle \frac{\pi }{2}}\right)+\delta \left(\psi \varphi \alpha (\varphi )3{\displaystyle \frac{\pi }{2}}\right)\right],`$ (95)
(96)
$`\alpha (\varphi )`$ being given by Eq. (C13). We recall that the angles in the arguments of the delta functions are defined modulo $`2\pi `$.
As a first check, set $`r_2^{}=0`$, corresponding to the case of no barriers. We obtain
$`d\widehat{P}_{0,0}(S){\displaystyle \frac{d\tau }{\pi \sqrt{\tau \left(1\tau \right)}}}{\displaystyle \frac{d\varphi }{2\pi }}`$ (97)
$`\times `$ $`{\displaystyle \frac{1}{2}}\left[\delta \left(\psi \varphi {\displaystyle \frac{\pi }{2}}\right)+\delta \left(\psi \varphi 3{\displaystyle \frac{\pi }{2}}\right)\right]d\psi .`$ (98)
Thanks to the delta functions, we recover the situation of LR symmetry. As expected, the right-hand side of Eq. (98) is the invariant measure defined for that symmetry, Eq. (87). As a second check, we analyze the case $`r_2^{}1`$, that corresponds to obstructing the waveguide on the right. We show in Appendix D that (93) gives in this case
$$d\widehat{P}_{0,1}\left(S\right)\delta \left(\tau \right)d\tau \frac{d\varphi }{2\pi }\frac{1}{2}\left[\delta \left(\psi \frac{\pi }{2}\right)+\delta \left(\psi 3\frac{\pi }{2}\right)\right]d\psi .$$
(99)
The conductance distribution reduces to a one-sided delta function at zero, as it should. Notice that the variable $`\varphi `$ is uniformly distributed in the two extreme cases $`r_2^{}=0`$ and $`r_2^{}=1`$; this is not so for an arbitrary value of $`r_2^{}`$. In the limiting case $`r_2^{}=1`$ we end up with a LR-symmetric ballistic cavity connected to just one lead (see Fig. 4): the resulting $`1`$-dimensional $`S`$ matrix, i.e. $`r=e^{2i\varphi }`$, is distributed according to the invariant measure: there is thus no memory left of the LR symmetry of the cavity. In fact, the right-hand side of Eq. (99) is identical to that of Sect. II A 2, Eq. (33), for an AS cavity with the right waveguide obstructed. As we shall see later on, in Sec. V, this is a peculiarity of the 1-channel case.
To get the joint distribution of $`\tau `$ and $`\varphi `$ for arbitrary $`r_2^{}`$ we integrate (96) over $`\psi `$. We find
$`q_{r_2^{}}(\tau ,\varphi )`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle \frac{\left(1r_{2}^{}{}_{}{}^{2}\right)^{3/2}}{\sqrt{\tau }}}`$ (100)
$`\times `$ $`{\displaystyle \frac{\left|\sqrt{1\tau }r_2^{}e^{2i\varphi }\right|}{\left|\sqrt{1\tau }\left(1r_{2}^{}{}_{}{}^{2}\right)r_2^{}\tau e^{2i\varphi }\right|^2}}.`$ (101)
The $`T=\tau `$ distribution $`w(T)`$ is obtained by integrating $`q_{r_2^{}}(\tau ,\varphi )`$ over $`\varphi `$. Fig. 5 shows (heavy lines) the evolution of $`w(T)`$ with the parameter $`r_2^{}`$. In the same figure we compare that sequence of distributions with those corresponding to an AS cavity (dotted lines) with the same $`r_2^{}`$. In the former case the system “remembers” in a rather conspicuous way that, although the resulting configuration is asymmetric, the cavity has LR symmetry.
## V Breaking the reflection symmetry of cavities with an asymmetric position of the waveguides
In the present section we study the effect of external mixing of LR symmetry in the absence of direct processes, i.e. for $`S=0`$: the problem will be that of a LR-symmetric cavity connected to two asymmetrically positioned waveguides in the absence of barriers. We proceed as follows. We first consider the LR-symmetric cavity connected to four symmetrically positioned waveguides (each supporting one open channel) by means of four, in general different, barriers, as shown in Fig. 6. The two barriers on the left side are then removed, while those on the right are made perfect reflectors.
We call $`S_0`$ the matrix associated with the LR-symmetric cavity connected to the four symmetrically located waveguides in the absence of barriers. The matrix $`S`$ in the presence of the four barriers is then given by Eq. (19), i.e.
$$S=r_b+t_b^{}\frac{1}{I_4S_0r_b^{}}S_0t_b,$$
(102)
where
$$t_b=t_b^{}=\left(\begin{array}{cccc}t_1& 0& 0& 0\\ 0& t_2& 0& 0\\ 0& 0& t_3& 0\\ 0& 0& 0& t_4\end{array}\right)$$
(103)
$$r_b=\left(\begin{array}{cccc}r_1& 0& 0& 0\\ 0& r_2& 0& 0\\ 0& 0& r_3^{}& 0\\ 0& 0& 0& r_4^{}\end{array}\right),r_b^{}=\left(\begin{array}{cccc}r_1^{}& 0& 0& 0\\ 0& r_2^{}& 0& 0\\ 0& 0& r_3& 0\\ 0& 0& 0& r_4\end{array}\right).$$
(104)
As explained above, we now open the left waveguides and block the right ones by means of perfect reflectors, so that
$$t_b=t_b^{}=\left(\begin{array}{cc}I_2& \hfill 0_2\\ 0_2& \hfill 0_2\end{array}\right),r_b=r_b^{}=\left(\begin{array}{cc}0_2& \hfill 0_2\\ 0_2& \hfill I_2\end{array}\right),$$
(105)
where $`I_2`$ and $`0_2`$ denote the two-dimensional unit and zero matrices, respectively.
The $`4\times 4`$ matrix $`S_0`$ has the structure (35), i.e.
$$S_0=\left(\begin{array}{cc}r_0& \hfill t_0\\ t_0& \hfill r_0\end{array}\right),$$
(106)
where $`r_0`$ and $`t_0`$ are 2-dimensional matrices. The matrix $`S`$ of Eq. (102) then reads
$$S=\left(\begin{array}{cc}r_0t_0\frac{1}{I_2+r_0}t_0& \hfill 0\\ 0& \hfill I_2\end{array}\right).$$
(107)
The 1-1 block of the above expression is the $`2\times 2`$ scattering matrix of the final system consisting of a LR-symmetric ballistic cavity connected to two waveguides on the left (see Fig. 7), i.e
$$s=r_0t_0\frac{1}{I_2+r_0}t_0.$$
(108)
Using the result (41) we can express $`r_0`$ and $`t_0`$ as
$$r_0=\frac{1}{2}\left[s^{(+)}+s^{()}\right],$$
(109)
$$t_0=\frac{1}{2}\left[s^{(+)}s^{()}\right],$$
(110)
where $`s^{(\pm )}`$ are $`2\times 2`$ unitary and symmetric matrices.
A numerical calculation was performed, in which 4-dimensional $`S_0`$ matrices were generated with a distribution corresponding to their invariant measure: this was done by constructing an ensemble of $`s^{(\pm )}`$ matrices, Eqs. (109), (110), distributed as two independent $`COE`$’s. From Eq. (108), the resulting $`S`$ matrices were then evaluated.
The distribution of the resulting transmission coefficient $`T`$ is shown in Fig. 8. For comparison, the distribution $`1/2\sqrt{T}`$ corresponding to an AS cavity connected to two one-channel waveguide and with $`S=0`$ is shown with dotted line. Although the LR-symmetric cavity with external symmetry breaking has a $`T`$ distribution very close to $`1/2\sqrt{T}`$, there is a statistically significant deviation which indicates that the resulting system has a memory of the point symmetry of the cavity. This is to be contrasted with the result mentioned right before Eq. (101) for the single 1-channel cavity illustrated in Fig. 4.
## VI Results and conclusions
One of the main purposes of the present paper has been the extension of previous studies on transport through ballistic chaotic cavities with reflection symmetry to include the presence of direct processes. In Sect. III we treated the problem of fully left-right (LR) symmetric systems in the presence of direct processes. The statistical distribution of the $`S`$ matrix, found analytically in Eq. (53), consists of the product of two Poisson’s kernels with the optical matrices $`s^{(+)}`$ and $`s^{()}`$, respectively. For no direct transmission processes, $`t=0`$, and real direct reflections $`r`$, we calculated analytically the distribution of the transmission coefficient $`w(T)`$ for the one-channel case. The difference with the $`T`$ distribution for an asymmetric cavity (AS) with the same optical matrix $`S`$, which is striking for $`r=0`$, becomes less dramatic as $`|r|`$ increases: that evolution is shown in Fig. 2.
The other main purpose of this work has been the study of LR-symmetry breaking by an asymmetric coupling of a LR-symmetric cavity to the outside. Two ways of producing “external mixing” of the spatial symmetry were analyzed:
##### a
In Sect. IV we studied the effect of breaking the reflection symmetry of a cavity by direct processes. The system consists of a ballistic cavity with reflection symmetry connected to two symmetrically positioned waveguides by means of barriers which, in general, are allowed to be different (Fig. 3). We found analytically, in Eqs. (93) and (96), the statistical distribution of the $`S`$ matrix for the one-channel case in each waveguide and, for simplicity, when only the barrier in the right waveguide is present ($`r_2^{}0`$). The $`T`$ distribution is strikingly different from that for the fully AS case (i.e the one in which the cavity itself is AS) having the same optical $`S`$ matrix, as shown in Fig. 5 for various values of $`r_2^{}0`$. We conclude that this two-waveguide system, although asymmetric with respect to the LR operation, has a memory of the reflection symmetry of the cavity from which it is constructed. In the limit $`r_2^{}1`$ the right waveguide is blocked and we end up with a LR-symmetric ballistic cavity connected, without any barrier, to just one lead, supporting one open channel (see Fig. 4). We found that the resulting 1-dimensional matrix $`S=e^{i\theta }`$ is distributed according to its invariant measure (i.e., $`\theta `$ is uniformly distributed) and, as a result, has no memory left of the LR symmetry of the cavity: this was found, though, to be a peculiarity of the one-waveguide–one-channel case (in fact, see the end of next paragraph).
##### b
In Sect. V we studied, in the absence of direct processes, the effect of external mixing of LR symmetry induced by an asymmetric position of the waveguides. The result is a LR-symmetric cavity connected, without any barriers, to two waveguides on its left-hand side (see Fig. 7). Let $`T`$ denote the total transmission coefficient between those two waveguides; its distribution $`w(T)`$ was calculated numerically for the one-channel case in each waveguide and compared, in Fig. 8, with $`1/2\sqrt{T}`$, the $`T`$ distribution arising from the invariant measure $`d\mu ^{(\beta =1)}(S)`$ for AS systems. Although the difference between the two distributions is quite small, it is statistically significant. This problem is clearly equivalent to having, on one side of the cavity, just one waveguide (coupled to the cavity without any barrier) supporting two open channels. In this one-waveguide–two-channel problem the resulting $`S`$ matrix is thus distributed very closely to its invariant measure, the difference exhibiting some memory left of the reflection symmetry of the cavity.
Two additional points are worth mentioning. First, from an experimental point of view, we notice that microwave cavities and acoustic systems might represent good possibilities to study the interplay between the symmetry of the cavity and external mixing in the statistical distribution of the conductance of such a structure. Finally, the problem described in b above is relevant to the study of transport between two one-channel leads connected by a “double” Cayley tree . In fact, under suitable circumstances the two problems can be mapped unto each other. This problem will be reported elsewhere.
###### Acknowledgements.
One of the authors (M. M.) wishes to acknowledge supports by DGAPA-UNAM, and CONACyT, México.
## A Derivation of Eq. (2.29)
We saw in Sec. II A 2 that the distribution $`dP_S(S)`$ of the scattering matrix of a cavity connected to two waveguides, where the one on the right of the cavity has a barrier, is given by
$$dP_{0,r_2^{}}(S)=\frac{\left(1|r_2^{}|^2\right)^{3/2}}{\left|1\sqrt{1\tau }e^{2i\psi }r_{2}^{}{}_{}{}^{}\right|^3}\frac{d\tau }{2\sqrt{\tau }}\frac{d\varphi }{2\pi }\frac{d\psi }{2\pi }.$$
(A1)
To see the behaviour of $`dP_{0,r_2^{}}`$ for $`r_2^{}=1`$, let $`r_2^{}`$ be a real number: assume for simplicity $`r_2^{}=\mathrm{cos}ϵ`$; we are interested in the limit $`ϵ0`$. Also, let us introduce the positive parameter $`\eta 1`$ in order to avoid the integrable singularity in $`\tau `$. Of course, we will take the limit $`\eta 0`$ later on. Because the variable $`\varphi `$ is uniformly distributed, the joint probability distribution of $`\tau `$ and $`\psi `$ can be written as
$$p_{\eta ,ϵ}(\tau ,\psi )=\frac{C_\eta }{4\pi \sqrt{\tau +\eta ^2}}\frac{\left|\mathrm{sin}ϵ\right|^3}{\left|1+\mathrm{cos}ϵ\sqrt{1\tau }e^{2i\psi }\right|^3},$$
(A2)
where $`C_\eta `$ is a normalization constant which depends on the parameter $`\eta `$.
We have the following properties of $`p_{\eta ,ϵ}(\tau ,\psi )`$:
1. From (A2) we see that
$$p_{0,0}(\tau ,\psi )=\underset{\eta 0}{lim}\underset{ϵ0}{lim}p_{\eta ,ϵ}(\tau ,\psi )=0$$
(A3)
for all $`\tau `$ and $`\psi `$, except for $`\tau =0`$ and $`\psi =\frac{\pi }{2},3\frac{\pi }{2}`$, where the denominator is zero
$$\left|1+\sqrt{1\tau }e^{2i\psi }\right|^3=0.$$
(A4)
2. For $`\tau =0`$ and $`\psi =\frac{\pi }{2},3\frac{\pi }{2}`$ we have
$`p_{0,0}\left(\tau =0,\psi ={\displaystyle \frac{\pi }{2}},3{\displaystyle \frac{\pi }{2}}\right)`$ (A5)
$`=\underset{\eta 0}{lim}\underset{ϵ0}{lim}p_{\eta ,ϵ}\left(\tau =0,\psi ={\displaystyle \frac{\pi }{2}},3{\displaystyle \frac{\pi }{2}}\right)`$ (A6)
$`=\underset{\eta 0}{lim}\underset{ϵ0}{lim}{\displaystyle \frac{C_\eta }{4\pi \eta }}\left|\mathrm{cot}{\displaystyle \frac{ϵ}{2}}\right|^3\mathrm{}`$ (A7)
3. The function $`p_{0,0}(\tau ,\psi )`$ is normalized to the unity:
$`{\displaystyle _0^1}𝑑\tau {\displaystyle _0^{2\pi }}𝑑\psi p_{0,0}(\tau ,\psi )`$ (A8)
$`=\underset{\eta ,ϵ0}{lim}{\displaystyle _0^1}𝑑\tau {\displaystyle _0^{2\pi }}𝑑\psi p_{\eta ,ϵ}(\tau ,\psi )=1`$ (A9)
Then the only function which satisfies those conditions is
$$p_{0,0}(\tau ,\psi )=\delta \left(\tau \right)\frac{1}{2}\left[\delta \left(\psi \frac{\pi }{2}\right)+\delta \left(\psi \frac{3\pi }{2}\right)\right].$$
(A10)
Finally, the distribution of the $`S`$ matrix in the above limit is given by Eq. (33).
## B Derivation of Eq. (3.11)
For $`r`$ real and $`s^{(\pm )}=e^{i\theta _\pm }`$, Eq. (59) can be written as
$$d\widehat{P}_r(S)=\frac{1r^2}{\left|1re^{i\theta _+}\right|^2}\frac{1r^2}{\left|1re^{i\theta _{}}\right|^2}\frac{d\theta _+}{2\pi }\frac{d\theta _{}}{2\pi }.$$
(B1)
The transmission amplitude is given by \[see Eq. (41)\]:
$$t=\frac{1}{2}\left(e^{i\theta _+}e^{i\theta _{}}\right),$$
(B2)
and the transmission coefficient is written as
$$T=\left|t\right|^2=\frac{1}{2}\left[1\mathrm{cos}\left(\theta _+\theta _{}\right)\right].$$
(B3)
The $`T`$ distribution $`w_r(T)`$ is obtained from
$$w_r(T)=\delta \left(T\frac{1}{2}\left[1\mathrm{cos}\left(\theta _+\theta _{}\right)\right]\right)𝑑\widehat{P}_r(S).$$
(B4)
We make the following change of variables in order to solve the integral:
$$\begin{array}{cc}\theta & =\frac{1}{2}\left(\theta _+\theta _{}\right),\\ \theta ^{}& =\frac{1}{2}\left(\theta _+\theta _{}\right);\end{array}$$
(B5)
the range of variation are: for $`\theta ^{}(0,2\pi )`$, $`\theta (\theta ^{},\theta ^{})`$ and for $`\theta ^{}(\pi ,2\pi )`$, $`\theta (2\pi +\theta ^{},2\pi \theta ^{})`$.
Substituting (B1) in (B4), considering the fact that the integrand is an even function of $`\theta `$ and writing the delta function in terms of its roots in the variable $`\theta `$ we have
$$\delta \left(T\mathrm{sin}^2\theta \right)=\frac{1}{2\sqrt{T\left(1T\right)}}\left[\delta \left(\theta \theta _1\right)+\delta \left(\theta \theta _2\right)\right],$$
(B6)
where $`\theta _2=\pi \theta _1`$ and $`\theta _1=\mathrm{arcsin}\sqrt{T}`$; finally, after some algebra, $`w_r(T)`$ can be written as a sum of two terms:
$$w_r(T)=\frac{\left(1r^2\right)^2}{\pi ^2\sqrt{T\left(1T\right)}}\left[I_1(T,r)+I_2(T,r)\right],$$
(B7)
where, for $`k=1,2`$,
$`I_k(T,r)`$ (B8)
$`=`$ $`{\displaystyle _0^\pi }𝑑\theta ^{}{\displaystyle _0^\theta ^{}}𝑑\theta {\displaystyle \frac{1}{\left[\left(1+r^2\right)2r\mathrm{cos}\left(\theta ^{}+\theta \right)\right]}}`$ (B9)
$`\times `$ $`{\displaystyle \frac{\delta \left(\theta \theta _k\right)}{\left[\left(1+r^2\right)2r\mathrm{cos}\left(\theta ^{}\theta \right)\right]}}.`$ (B10)
Again, after some algebra the sum of the two integrals give a single one:
$$I_1(T,r)+I_2(T,r)=\frac{1}{c}_0^\pi \frac{d\theta ^{}}{ab\mathrm{cos}\theta ^{}+\mathrm{cos}^2\theta ^{}},$$
(B11)
where
$`a`$ $`=`$ $`{\displaystyle \frac{1}{c}}\left[\left(1+r^2\right)^24r^2T\right],`$ (B12)
$`b`$ $`=`$ $`{\displaystyle \frac{4}{c}}r\left(1+r^2\right)\sqrt{1T},`$ (B13)
$`c`$ $`=`$ $`4r^2.`$ (B14)
Now, making the change of variable $`x=\mathrm{cos}\theta ^{}`$, (B7) can be written as
$$w_r(T)=\frac{\left(1r^2\right)^2}{4r^2\pi ^2\sqrt{T\left(1T\right)}}\left[I_+(T,r)+I_{}(T,r)\right],$$
(B15)
where now
$$I_\pm (T,r)=_0^1\frac{dx}{\sqrt{1x^2}\left(a\pm bx+x^2\right)}.$$
(B16)
By means of a change of variables
$$u=\frac{x+\left(A+B\right)}{x+\left(AB\right)},$$
(B18)
$$v=\frac{x\left(A+B\right)}{x\left(AB\right)},$$
(B19)
where
$`A`$ $`=`$ $`{\displaystyle \frac{1}{b}}\left(1+a\right),`$ (B20)
$`B`$ $`=`$ $`{\displaystyle \frac{1}{b}}\sqrt{\left(1+a\right)^2b^2},`$ (B21)
the indefinite integrals, $`Indef_\pm `$, corresponding to each one of the above, can be transformed to
$$Indef_+=\frac{2B}{\sqrt{C}D}\frac{\left|u+1\right|}{\sqrt{u^2+p}\left(u^2+q\right)}𝑑u,$$
(B23)
$$Indef_{}=\frac{2B}{\sqrt{C}D}\frac{\left|v+1\right|}{\sqrt{v^2+p}\left(v^2+q\right)}𝑑v.$$
(B24)
where
$`p`$ $`=`$ $`{\displaystyle \frac{ab\left(B+A\right)+\left(B+A\right)^2}{a+b\left(BA\right)+\left(BA\right)^2}}`$ (B25)
$`q`$ $`=`$ $`{\displaystyle \frac{1\left(B+A\right)^2}{1\left(BA\right)^2}}`$ (B26)
and
$$\begin{array}{cc}C& =1(BA)^2,\hfill \\ D& =a+b(BA)+(BA)^2.\hfill \end{array}$$
(B27)
Although the integrals (B16) seem to give the same result under the change $`bb`$, they do not, because the cutoff $`x_u=BA`$ in (B23), and $`x_v=AB`$ in (B24), are different. One must be careful evaluating the integrals in the limits. The results are
$`I_+(T,r)`$ $`=`$ $`{\displaystyle \frac{2B}{\sqrt{C}D\sqrt{pq}}}[\mathrm{arctan}\left({\displaystyle \frac{1r^2}{2r\sqrt{1T}}}\right)`$ (B29)
$``$ $`{\displaystyle \frac{1}{2\sqrt{p}}}\mathrm{ln}\left({\displaystyle \frac{1+r^2+2r\sqrt{T}}{1+r^22r\sqrt{T}}}\right)],`$ (B30)
$`I_{}(T,r)`$ $`=`$ $`{\displaystyle \frac{2B}{\sqrt{C}D\sqrt{pq}}}[\pi \mathrm{arctan}\left({\displaystyle \frac{1r^2}{2r\sqrt{1T}}}\right)`$ (B31)
$`+`$ $`{\displaystyle \frac{1}{2\sqrt{p}}}\mathrm{ln}\left({\displaystyle \frac{1+r^2+2r\sqrt{T}}{1+r^22r\sqrt{T}}}\right)].`$ (B32)
Now, we substitute the sum of equations (B) in (B15) to obtain the result
$$w_r(T)=\frac{\left(1r^2\right)^2}{4r^2\pi ^2\sqrt{T\left(1T\right)}}\frac{2\pi B}{\sqrt{C}D\sqrt{pq}};$$
(B33)
using Eqs. (B13), (B21), (B25), (B26) and (B27) the final result (60) is obtained.
## C Derivation of Eqs. (4.23), (4.24)
For the particular case in which barrier 1 is transparent \[see Fig. 3\], so that its scattering matrix $`S_1`$ of Eq. (14) is the Pauli matrix $`\sigma _x`$, and barrier 2 is described by Eq. (17) with real matrix elements, Eq. (92) can be written as
$`d\widehat{P}_{0,r_2^{}}(S)`$ $``$ $`2{\displaystyle \frac{\delta \left(\psi _0\varphi _0\frac{\pi }{2}\right)+\delta \left(\psi _0\varphi _03\frac{\pi }{2}\right)}{\sqrt{1\tau _0}}}`$ (C1)
$`\times `$ $`{\displaystyle \frac{\left(1r_{2}^{}{}_{}{}^{2}\right)^{3/2}}{\left|1\sqrt{1\tau }r_2^{}e^{2i\psi }\right|^3}}{\displaystyle \frac{d\tau }{2\sqrt{\tau }}}{\displaystyle \frac{d\varphi }{2\pi }}{\displaystyle \frac{d\psi }{2\pi }}.`$ (C2)
Also, the transformation $`S_0(S)`$ given by Eq. (63) can be written in terms of its elements as follows:
$`r_0`$ $`=`$ $`{\displaystyle \frac{1}{1r_2^{}r^{}}}\left[r(1r_2^{}r^{})+r_2^{}t^2\right],`$ (C3)
$`r_0^{}`$ $`=`$ $`{\displaystyle \frac{1}{1r_2^{}r^{}}}\left(r^{}r_2^{}\right),`$ (C4)
$`t_0`$ $`=`$ $`{\displaystyle \frac{1}{1r_2^{}r^{}}}t_2t,`$ (C5)
or in terms of the independent parameters \[see Eqs.(24) and (65)\] as:
$`\sqrt{1\tau _0}e^{2i\varphi _0}`$ $`=`$ $`e^{2i\varphi }{\displaystyle \frac{\sqrt{1\tau }r_2^{}e^{2i\psi }}{1r_2^{}\sqrt{1\tau }e^{2i\psi }}}`$ (C7)
$`\sqrt{1\tau _0}e^{2i\psi _0}`$ $`=`$ $`e^{2i\psi }{\displaystyle \frac{\sqrt{1\tau }r_2^{}e^{2i\psi }}{1r_2^{}\sqrt{1\tau }e^{2i\psi }}}`$ (C8)
$`\sqrt{\tau _0}e^{i(\varphi _0+\psi _0)}`$ $`=`$ $`{\displaystyle \frac{t_2\sqrt{\tau }e^{i(\varphi +\psi )}}{1r_2^{}\sqrt{1\tau }e^{2i\psi }}}`$ (C9)
From (C7) or (C8) we find
$$\sqrt{1\tau ^{(0)}}=\frac{\left|\sqrt{1\tau }r_2^{}e^{2i\psi }\right|}{\left|1\sqrt{1\tau }r_2^{}e^{2i\psi }\right|};$$
(C10)
also, dividing the (C7) by (C8) we obtain
$$e^{2i(\psi _0\varphi _0)}=e^{2i(\psi \varphi )}\frac{\sqrt{1\tau }r_2^{}e^{2i\psi }}{\sqrt{1\tau }r_2^{}e^{2i\psi }}.$$
(C11)
Because the roots of the delta functions appearing in Eq. (92) satisfy $`e^{2i(\psi _0\varphi _0)}=1`$, from (C11) we find
$$e^{2i\psi }=e^{2i\varphi }e^{2i\alpha (\varphi )}$$
(C12)
where
$$e^{i\alpha (\varphi )}=\frac{\sqrt{1\tau }r_2^{}e^{2i\varphi }}{\left|\sqrt{1\tau }r_2^{}e^{2i\varphi }\right|}.$$
(C13)
Then, we have the conditions for $`\psi `$:
$$\begin{array}{cc}\psi \varphi \alpha (\varphi )=\frac{\pi }{2}\hfill & \mathrm{for}\psi _0\varphi _0=\frac{\pi }{2}\hfill \\ \psi \varphi \alpha (\varphi )=3\frac{\pi }{2}\hfill & \mathrm{for}\psi _0\varphi _0=3\frac{\pi }{2}.\hfill \end{array}$$
(C14)
The Jacobian for the transformation $`\psi _0\psi `$ is
$$\left|\frac{}{\psi }\left(\psi _0\varphi _0\right)\right|=\frac{\left|(1\tau )r_{2}^{}{}_{}{}^{2}\right|}{\left|\sqrt{1\tau }r_2^{}e^{2i\psi }\right|^2}.$$
(C15)
Then we write
$`\delta \left(\psi _0\varphi _0{\displaystyle \frac{2n+1}{2}}\pi \right)`$ (C16)
$`=`$ $`{\displaystyle \frac{\left|\sqrt{1\tau }r_2^{}e^{2i\psi }\right|^2}{\left|(1\tau )r_{2}^{}{}_{}{}^{2}\right|}}\delta \left[\psi \varphi \alpha (\varphi ){\displaystyle \frac{2n+1}{2}}\pi \right],`$ (C17)
for $`n=0,1`$.
From (C12) and (C13) we find
$`\sqrt{1\tau }r_2^{}e^{2i\psi }`$ $`=`$ $`{\displaystyle \frac{(1\tau )r_{2}^{}{}_{}{}^{2}}{\sqrt{1\tau }r_2^{}e^{2i\varphi }}},`$ (C19)
$`1r_2^{}\sqrt{1\tau }e^{2i\psi }`$ $`=`$ $`{\displaystyle \frac{\sqrt{1\tau }\left(1r_{2}^{}{}_{}{}^{2}\right)r_2^{}\tau e^{2i\varphi }}{\sqrt{1\tau }r_2^{}e^{2i\varphi }}}.`$ (C20)
Finally, substituting Eqs. (C10), (LABEL:deltapsi0-phi0), (C19) and (C20) in Eq. (C2), we arrive to
$$d\widehat{P}_{0,r_2^{}}(\tau ,\varphi ,\psi )p_{r_2^{}}(\tau ,\varphi ,\psi )d\tau d\varphi d\psi ,$$
(C21)
where $`p_{r_2^{}}(\tau ,\varphi ,\psi )`$ is given by Eq. (96).
## D Derivation of Eq. (4.26)
In Sect. IV we find the joint distribution of $`\tau `$, $`\varphi `$ and $`\psi `$ \[Eq. (96)\]. From that it is easy to integrate over $`\psi `$ to find the joint distribution of $`\tau `$ and $`\varphi `$ to be
$$q_{r_2^{}}(\tau ,\varphi )=\frac{\left(1r_{2}^{}{}_{}{}^{2}\right)^{3/2}\left|\sqrt{1\tau }r_2^{}e^{2i\varphi }\right|}{(2\pi )^2\sqrt{\tau }\left|\sqrt{1\tau }\left(1r_{2}^{}{}_{}{}^{2}\right)r_2^{}\tau e^{2i\varphi }\right|^2}$$
(D1)
As in App. A we assume for simplicity $`r_2^{}=\mathrm{cos}ϵ`$; again we introduce the parameter $`\eta 1`$. Of course, we will take the limits $`\eta `$, $`ϵ0`$; then
$$q_{\eta ,ϵ}(\tau ,\varphi )=\frac{C_\eta }{2\pi ^2}\frac{|\mathrm{sin}ϵ|^3}{\sqrt{\tau +\eta ^2}}\frac{\left|\sqrt{1\tau }+\mathrm{cos}ϵe^{2i\varphi }\right|}{\left|\sqrt{1\tau }\mathrm{sin}^2ϵ+\tau \mathrm{cos}ϵe^{2i\varphi }\right|^2}.$$
(D2)
where $`C_\eta `$ is a normalization constant which depends on $`\eta `$.
Again, as before we have the following properties for $`q_{\eta ,ϵ}(\tau ,\varphi )`$:
1. From (D2) we see that
$$q_{0,0}(\tau ,\varphi )=\underset{\eta 0}{lim}\underset{ϵ0}{lim}q_{\eta ,ϵ}(\tau ,\varphi )=0$$
(D3)
for all $`\tau `$ and $`\varphi `$, except for $`\tau =0`$ and $`\varphi =\frac{\pi }{2},3\frac{\pi }{2}`$, where the denominatror is zero:
$$\left|\sqrt{1\tau }\mathrm{sin}^2ϵ+\tau \mathrm{cos}ϵe^{2i\varphi }\right|^2=0$$
(D4)
2. $`\tau 0`$ and $`\varphi `$.
It is easy to see from (D2) that in this case
$$q_{0,0}(\tau 0,\varphi )=\underset{\eta }{lim}\underset{ϵ0}{lim}q_{\eta ,ϵ}(\tau ,\varphi )=0.$$
(D5)
3. For $`\tau =0`$ and $`\varphi =\frac{\pi }{2},3\frac{\pi }{2}`$ we have
$`q_{0,0}\left(\tau =0,\varphi ={\displaystyle \frac{\pi }{2}},3{\displaystyle \frac{\pi }{2}}\right)`$ (D6)
$`=\underset{\eta 0}{lim}\underset{ϵ0}{lim}q_{\eta ,ϵ}\left(\tau =0,\varphi ={\displaystyle \frac{\pi }{2}},3{\displaystyle \frac{\pi }{2}}\right)`$ (D7)
$`=\underset{\eta 0}{lim}\underset{ϵ0}{lim}{\displaystyle \frac{C_\eta }{2\pi ^2\eta }}\left|\mathrm{tan}{\displaystyle \frac{ϵ}{2}}\right|`$ (D8)
4. For $`\tau =0`$, $`\varphi \frac{\pi }{2},3\frac{\pi }{2}`$ we obtain
$`q_{0,0}\left(\tau =0,\varphi {\displaystyle \frac{\pi }{2}},3{\displaystyle \frac{\pi }{2}}\right)`$ (D9)
$`=\underset{\eta 0}{lim}\underset{ϵ0}{lim}q_{\eta ,ϵ}\left(\tau =0,\varphi {\displaystyle \frac{\pi }{2}},3{\displaystyle \frac{\pi }{2}}\right)`$ (D10)
$`=\underset{\eta 0}{lim}\underset{ϵ0}{lim}{\displaystyle \frac{C_\eta }{2\pi ^2\eta }}{\displaystyle \frac{\left|1+\mathrm{cos}ϵe^{2i\varphi }\right|}{\left|\mathrm{sin}ϵ\right|}}\mathrm{}.`$ (D11)
5. Also, the function $`q_{0,0}(\tau ,\varphi )`$ is normalized to the unity:
$`{\displaystyle _0^1}𝑑\tau {\displaystyle _0^{2\pi }}q_{0,0}(\tau ,\varphi )`$ (D12)
$`=\underset{\eta ,ϵ0}{lim}{\displaystyle _0^1}𝑑\tau {\displaystyle _0^{2\pi }}q_{\eta ,ϵ}(\tau ,\varphi )𝑑\tau 𝑑\varphi =1.`$ (D13)
These conditions defines the function
$$q_{0,0}(\tau ,\varphi )=\delta \left(\tau \right)\frac{1}{2\pi }.$$
(D14)
We thus arrive at Eq. (99).
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# 1. Introduction
## 1. Introduction
Landau levels have been introduced in 1930 (see \[LL58\]). They found an important physical application only quite recently, after the discovery of the Quantum Hall Effect (see \[Aok87, Fub92, Fer94\] and references therein). More recently it has been recognized that the theory of Landau levels provides a general bridge between Feynman path integrals and “geometric quantization” in all cases where the classical phase space is equipped with a complex structure which makes it a Kaehler manifold (see \[KO89\]). From this general viewpoint, or to get a more realistic description of conducting thin films, it is important to understand the case of a finite region with suitable boundary conditions. If these correspond to a compact (smooth) manifold without boundary the quantization condition of Kostant-Souriau (see \[SW76\] and references therein) or, equivalently, Dirac’s quantization condition for monopole charges require that the total magnetic flux be quantized, i.e. it must be an integral multiple $`N`$ of the universal constant $`hc/e`$. At the same time, the degeneracy of the ground state is finite and coincides with $`N`$, except for a topological correction (half the Euler characteristic of the manifold). A similar, approximate, result can be obtained by a semi-classical argument (see \[LL58\]).
Consider now the simple case of a rectangular area with sides $`L_1,L_2`$ and periodic boundary conditions; the problem is formulated on a toroidal surface with a transverse magnetic field, whose flux is $`BL_1L_2`$. Of course this fact implies the presence of magnetic charges, hence Dirac’s quantization. The problem is: what is the symmetry of the Hamiltonian? We expect that the classical symmetry of the torus ($`S^1\times S^1`$) be realized as a projective representation, the two infinitesimal generators satisfying Heisenberg algebra with a central charge $`\mathrm{}B`$ (at least this is what happens in the non compact $`R^2`$ case). But this is clearly incompatible (Wintner’s theorem) with finite degeneracy of energy levels! While we cannot expect a spontaneous symmetry breaking in a system with a finite number of degrees of freedom, we know from Geometric Quantization that not all classical symmetries survive at the quantum level, only those which are lifted at the pre-quantum level and, secondly, respect the polarization (in Landau level language, those symmetries are preserved which leave the first Landau level invariant). The problem is: what exactly is happening on the torus?
To get an answer, we shall reconstruct the explicit form of the Landau levels in terms of sections of the hermitian line bundle associated to the principal bundle with connection given by the magnetic potential $`𝑨`$. The language of fibre bundles is the natural one to describe gauge fields and it is becoming more familiar to physicists especially after the advent of modern string theory. We shall explicitly construct the transition functions of the line bundle and find a natural orthonormal basis of holomorphic sections, which turn out to be Jacobi-$`\theta `$-functions. By inspection it turns out that translation invariance is broken to a discrete subgroup $`Z_N\times Z_N`$, $`N`$ being the monopole charge. This fact has the counterpart that the Hermitian operators which correspond to infinitesimal translations (in the non–compact case) do not leave the Landau levels invariant, i.e. they do not commute with the Hamiltonian: while formally commuting with the Hamiltonian as differential operators, they fail to respect the boundary conditions, given by the bundle transition functions.
## 2. Magnetic field on the torus.
Let $`𝕋^2=^2/^2`$ denote the two-torus; we describe it in physical terms by identifying an atlas of four local charts specified as follows:
1. $`𝒰_\alpha =\{0<x<L_1,\mathrm{\hspace{0.33em}0}<y<L_2\}`$
2. $`𝒰_\beta =\{\overline{x}<x<L_1+\overline{x},\mathrm{\hspace{0.33em}0}<y<L_2\}`$
3. $`𝒰_\gamma =\{0<x<L_1,\overline{y}<y<L_2+\overline{y}\}`$
4. $`𝒰_\delta =\{\overline{x}<x<L_1+\overline{x},\overline{y}<y<L_2+\overline{y}\}`$
with some choice of constants $`\overline{x}`$ and $`\overline{y}`$. A uniform magnetic field $`B`$ transverse to the surface $`𝕋^2`$ is represented by the translation invariant two-form $`𝑩=Bdxdy`$. The connection form $`𝑨`$, representing the magnetic potential, is defined in each local chart in such a way that $`d𝑨=𝑩`$. It is well–known that a global one-form on $`𝕋^2`$ satisfying this condition does not exist, since otherwise $`_{𝕋^2}𝑩=_{𝕋^2}𝑑𝑨=0`$, by Stokes theorem, while it holds $`_{𝕋^2}𝑩=BL_1L_2`$. The problem is essentially the same as the presence of a Dirac string in the case of a three-dimensional magnetic monopole. For the sake of simplicity, we may define the local connection forms by the same formula $`𝑨=\frac{1}{2}B(xdyydx)`$, since the coordinates $`x,y`$ are indeed differentiable within each local chart <sup>3</sup><sup>3</sup>3The correct mathematical language to describe such a setup is that of algebraic geometry; a nice introduction for physicists can be found for instance in \[Alv85\]. In this paper we try to keep the mathematical jargon to a minimum.. To characterize the connection form completely we have to identify the transition functions which relate $`𝑨_i`$ to $`𝑨_j`$ for any pair $`(i,j)`$ in the set $`\{\alpha ,\beta ,\gamma ,\delta \}`$ and for each connected component of the overlap $`𝒰_i𝒰_j`$; we have
$`𝑨_\beta (x,y)`$ $`=`$ $`𝑨_\alpha (x,y)(\overline{x}<x<L_1)`$
$`𝑨_\gamma (x,y)`$ $`=`$ $`𝑨_\alpha (x,y)(\overline{y}<y<L_2)`$
$`𝑨_\delta (x,y)`$ $`=`$ $`𝑨_\alpha (x,y)(\overline{x}<x<L_1,\overline{y}<y<L_2)`$
$`𝑨_\beta (x+L_1,y)`$ $`=`$ $`𝑨_\alpha (x,y)+d\left(\frac{1}{2}BL_1y+\phi _{\alpha \beta }\right)(0<x<\overline{x})`$
$`𝑨_\gamma (x,y+L_2)`$ $`=`$ $`𝑨_\alpha (x,y)+d\left(\frac{1}{2}BL_2x+\phi _{\alpha \gamma }\right)(0<y<\overline{y})`$
$`𝑨_\delta (x+L_1,y+L_2)`$ $`=`$ $`𝑨_\alpha (x,y)+d(\frac{1}{2}B(L_1yL_2x)+\phi _{\alpha \delta })(0<x<\overline{x},0<y<\overline{y})`$
and similar transition functions for the other cases. The constants $`\phi _{ij}`$ are arbitrary at this level; they will however play a crucial role in the lifting to the associated line bundle which describes the quantum wave functions<sup>4</sup><sup>4</sup>4The constants $`\phi _{\alpha \beta }`$ are connected to the fundamental cocycle $`c_{\alpha \beta \gamma }`$ of Ref.\[SW76, Alv85\]..
## 3. The holomorphic gauge
We now make a gauge transformation to a special gauge which is particularly convenient in the quantization process. Let us introduce complex coordinates $`z=x+iy,\overline{z}=xiy`$. The magnetic potential is given by
(1)
$$𝑨(z,\overline{z})=\frac{1}{2i}B\overline{z}dz\frac{1}{4i}Bd|z|^2$$
which shows that by a gauge transformation we can adopt a holomorphic form
$$𝑨^h=\frac{1}{2i}B\overline{z}dz,$$
for which we have the transition functions
$`𝑨_\beta ^h(z)`$ $`=`$ $`𝑨_\alpha ^h(z)(\overline{x}<x<L_1)`$
$`𝑨_\gamma ^h(z)`$ $`=`$ $`𝑨_\alpha ^h(z)(\overline{y}<y<L_2)`$
$`𝑨_\delta ^h(z)`$ $`=`$ $`𝑨_\alpha ^h(z)(\overline{x}<x<L_1,\overline{y}<y<L_2)`$
$`𝑨_\beta ^h(z+L_1)`$ $`=`$ $`𝑨_\alpha ^h(z)+d\left(i\frac{1}{2}BL_1z+\phi _{\alpha \beta }\right)(0<x<\overline{x})`$
$`𝑨_\gamma ^h(z+iL_2)`$ $`=`$ $`𝑨_\alpha ^h(z)+d\left(\frac{1}{2}BL_2z+\phi _{\alpha \gamma }\right)(0<y<\overline{y})`$
$`𝑨_\delta ^h(z+L_1+iL_2)`$ $`=`$ $`𝑨_\alpha ^h(z)+d(i\frac{1}{2}B(L_1iL_2)z+\phi _{\alpha \delta })(0<x<\overline{x},0<y<\overline{y})`$
with some new choice of constants $`\phi _{ij}`$.
## 4. Quantization
The Hamiltonian for a charged particle is given by the minimal-coupling prescription. The local expression as a differential operator must be complemented by suitable boundary conditions which ensure selfadjointness. This is easily done in terms of a line bundle associated to $`𝑨`$ as defined in the previous section. The physical principle to adopt is the gauge principle, according to which
$$\left(i\mathrm{}_\mu \frac{e}{c}A_\mu \right)\psi $$
is covariant under gauge transformations, in particular under the transition from one chart to another (here $`A_\mu `$ are the components of the gauge potential one-form $`𝑨=A_\mu dx^\mu `$). This can be done directly in the holomorphic gauge, which is our choice for the sequel. As usual, the complex line bundle has transition functions obtained by exponentiating those which characterize $`𝑨^h`$:
$`\psi _\beta (z)`$ $`=`$ $`\psi _\alpha (z)(\overline{x}<x<L_1)`$
$`\psi _\gamma (z)`$ $`=`$ $`\psi _\alpha (z)(\overline{y}<y<L_2)`$
$`\psi _\delta (z)`$ $`=`$ $`\psi _\alpha (z)(\overline{x}<x<L_1,\overline{y}<y<L_2)`$
$`\psi _\beta (z+L_1)`$ $`=`$ $`\psi _\alpha (z)\mathrm{exp}\left\{{\displaystyle \frac{eBL_1}{2\mathrm{}c}}z+\varphi _{\alpha \beta }\right\}(0<x<\overline{x})`$
$`\psi _\gamma (z+iL_2)`$ $`=`$ $`\psi _\alpha (z)\mathrm{exp}\left\{{\displaystyle \frac{ieBL_2}{2\mathrm{}c}}z+\varphi _{\alpha \gamma }\right\}(0<y<\overline{y})`$
$`\psi _\delta (z+L_1+iL_2)`$ $`=`$ $`\psi _\alpha (z)\mathrm{exp}\{{\displaystyle \frac{eB(L_1iL_2)}{2\mathrm{}c}}z+\varphi _{\alpha \delta }\}(0<x<\overline{x},0<y<\overline{y})`$
where we have redefined the constants $`\phi \varphi `$ to absorb a common factor $`ieB/\mathrm{}c`$. It is clear that both $`\overline{}\psi `$ and $`(z)\psi `$ transform in the same way as $`\psi `$. We have to stress here that while $`\varphi _{ij}`$ are totally arbitrary, they must be chosen once for all to define the Hamiltonian; as we shall show, different choices correspond in general to unitarily equivalent, yet distinct, operators. The situation is rather different from the well-known Aharonov-Bohm case, where the various admissible boundary conditions yield inequivalent Hamiltonians.
The local expression of the Hamiltonian in terms of complex coordinates is easily found to be
(2)
$$H^h=4\frac{\mathrm{}^2}{2m}\left(\frac{eB}{2mc\mathrm{}}\overline{z}\right)\overline{}$$
($`/z,\overline{}/\overline{z}`$) where we dropped a zero point energy term $`\mathrm{}\omega `$.
We must now introduce the Hermitian structure which allows to define the quantum inner product between wave functions. It is readily seen (e.g. starting from the Euclidean inner product in the real gauge and performing the gauge transformation to the holomorphic case) that the Hermitian structure is given by
$$h(\psi _1,\psi _2)=\mathrm{exp}\{\frac{eB}{2\mathrm{}c}|z|^2\}\overline{\psi _1}\psi _2.$$
in terms of which we can define the quantum inner product
$$\psi _1|\psi _2=_{𝕋^2}h(\psi _1,\psi _2)[dz]$$
where $`[dz]\frac{1}{2i}\overline{dz}dz`$. It is easy to check that there is a smooth match $`h(\psi _i,\psi _i)=h(\psi _j,\psi _j)`$ on each $`𝒰_i𝒰_j`$ provided that $`\mathrm{}(\varphi _{ij})`$ be suitably chosen. To simplify the notation, let us introduce natural units adapted to the problem: let us use $`(\mathrm{}/m\omega )^{{\scriptscriptstyle \frac{1}{2}}}`$ as length unit, where $`\omega =eB/2mc`$ is Larmor’s frequency. Then we get the new transition functions which make $`h(\psi ,\varphi )`$ smooth. At this point we can drop the chart index from the wave-function: from our convention, there is an open set common to all local chart where the wave function is the same in all local charts and the transition functions merely represent the boundary conditions to be imposed on $`\psi `$.
(3)
$$\begin{array}{cc}\hfill \psi (z+L_1)& =\psi (z)\mathrm{exp}\left\{L_1z+\frac{1}{2}L_1^2+i\delta _1\right\}\hfill \\ \hfill \psi (z+iL_2)& =\psi (z)\mathrm{exp}\left\{iL_2z+\frac{1}{2}L_2^2+i\delta _2\right\}\hfill \end{array}$$
It is also easily checked that these b.c. make the Hamiltonian hermitian. (Hint: make use of the complex integration by parts in the form $`_{𝕋^2}\overline{dz}dz\overline{\varphi (z)}\overline{}\psi (z)=_{𝕋^2}\overline{\varphi (z)}d\left(\psi dz\right)=\overline{\varphi (z)}\psi (z)𝑑z_{𝕋^2}\overline{dz}dz\overline{\varphi (z)}\psi `$ ).
However, there is a consistency condition to be satisfied, which stems from a general theorem about hermitian line bundles due to A.Weil (see \[Wei58\]; a simple proof taken from \[SW76\] is reproduced in Appendix B). In our case it can be found as follows: by successively applying the previous relations we get
(4)
$$\begin{array}{cc}\hfill \psi (z+L_1+iL_2)& =\psi (z+L_1)\mathrm{exp}\{iL_2(z+L_1)+\frac{1}{2}L_2^2+\delta _2\}\hfill \\ & =\psi (z)\mathrm{exp}\{(L_1iL_2)z+\frac{1}{2}|L_1+iL_2|^2+i\delta _1+i\delta _2iL_1L_2\}\hfill \\ & =\psi (z+iL_2)\mathrm{exp}\{L_1(z+iL_2)+\frac{1}{2}L_1^2+\delta _1\}\hfill \\ & =\psi (z)\mathrm{exp}\{(L_1iL_2)z+\frac{1}{2}|L_1+iL_2|^2+i\delta _1+i\delta _2+iL_1L_2\}\hfill \end{array}$$
hence
$$2L_1L_2=2N\pi ,$$
which is Dirac-Weil-Kostant-Souriau quantization condition. Let us conclude this section by giving the explicit expression for $`\psi |H|\psi `$, which exhibits $`H`$ as a positive operator:
$$\psi |H|\psi =_{𝕋^2}[dz]e^{|z|^2}|\overline{}\psi |^2,$$
a general result for quantum mechanics on Kaehler manifolds \[KO89\], from which we get the general result that the ground state coincides with the subspace of holomorphic sections ($`\overline{}\psi =0`$).
## 5. Finite–dimensional Landau levels
We can now compute the solutions of Schroedinger equation belonging to the ground state. These are given by holomorphic functions satisfying the boundary conditions (3). Let us choose $`\delta _1=\delta _2=0`$. Setting $`\psi (z)=\mathrm{exp}\{\frac{1}{2}z^2\}\theta (z)`$, we find that $`\theta `$ must be periodic with real period $`L_1`$ hence it can be expanded in a Fourier series $`\theta (z)=_{n=\mathrm{}}^{\mathrm{}}c_ne^{2\pi inz/L_1}`$). It follows
$$\begin{array}{c}\hfill \psi (z+iL_2)=e^{{\scriptscriptstyle \frac{1}{2}}(z+iL_2)^2}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}c_ne^{2\pi inz/L_1}e^{2\pi nL_2/L_1}\\ \hfill =e^{{\scriptscriptstyle \frac{1}{2}}z^2}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}c_ne^{2\pi inz/L_1}e^{iL_2z+{\scriptscriptstyle \frac{1}{2}}L_2^2}.\end{array}$$
which gives
$$e^{2iL_2zL_2^2}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}c_ne^{2\pi inz/L_1}e^{2n\pi L_2/L_1}=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}c_ne^{2\pi inz/L_1}.$$
Making use of Dirac’s quantization ($`L_2=N\pi /L_1`$) we get the condition
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}c_ne^{2\pi i(n+N)z/L_1}e^{2n\pi L_2/L_1L_2^2}=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}c_ne^{2\pi inz/L_1}$$
which is readily transformed into the recurrence relation
$$c_n=c_{nN}e^{2n\pi L_2/L_1+2N\pi L_2/L_1L_2^2}$$
whose solution is
$$c_n=e^{\pi n^2L_2/(L_1N)}b_n$$
where $`b_n`$ is such that $`b_n=b_{n+N}`$. Hence there are $`N`$ orthogonal solutions given by
$$\{\psi _\nu (z)=𝒩_\nu e^{{\scriptscriptstyle \frac{1}{2}}z^2}\underset{n\nu mod(N)}{}\mathrm{exp}\{\frac{\pi n^2L_2}{NL_1}+2n\pi iz/L_1\}|\nu =0,1,\mathrm{},N1\}.$$
We can obtain a new representation in terms Gaussian functions, very convenient for a practical calculation of $`\psi `$, by applying Poisson’ summation formula (see for e.g. \[Lig64\]). We find
$$\psi _\nu (z)=𝒩_\nu e^{{\scriptscriptstyle \frac{1}{2}}z^2+2\pi i\nu z/L_1}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{exp}\{(z+nL_1/N+i\nu L_2/N)^2\}.$$
Higher levels can be simply obtained by applying the covariant creation operator $`\overline{z}`$ to each $`\psi _\nu `$.
## 6. Translation symmetry breaking
The main question which started this investigation was the following: what happens to the translation symmetry of the torus? The question is motivated from the fact that unitary translations are realized as projective representations with a “central charge” given by the magnetic field strength. Hence they cannot live in a finite dimensional space (see Appendix A). Before going to analyze the problem in great detail, just observe that under the assumption that such a translation symmetry would nevertheless survive in some way, we should see it as a property of the ground state, i.e. there must exist a finite unitary matrix $`t_{\mu \nu }`$ such that $`(T_a\psi _\nu )(z)=_\mu t_{\nu \mu }(a)\psi _\mu (z)`$. It would follow that the density matrix $`\rho _N(z)=_\nu |\psi _\nu (z)|^2`$ should then be translation invariant, i.e. constant on the torus. If we calculate $`\rho _N`$ for the first few values of $`N`$ we immediately find that this is not so. The density $`\rho `$ exhibits a series of regularly spaced bumps, precisely at the location $`(n_1L_1+n_2L_2)/N`$ (see Fig.s 1-2, where the deviation from uniformity is plotted for the first two Landau levels at various values of the magnetic charge).
As is clear from the pictures, translation symmetry is broken, presumably to $`Z_N\times Z_N`$, but the breaking tends to be weaker at high $`N`$ (a variation of $`O(10^N)`$). Is there a simple explanation of this symmetry breaking? The point is that we can easily implement compact translations in the same way as we can do in the non-compact case. The unitary operators are given by
$$(T_a\psi )(z)=e^{\overline{a}z{\scriptscriptstyle \frac{1}{2}}|a|^2}\psi (za)$$
where the value of $`\psi `$ should be found through the twisted periodicity conditions given in Eq. (3). It is readily checked that
1. $`T_a`$ formally commute with the Hamiltonian, i.e. with the differential operator of Eq. (2);
2. $`T_aT_bT_aT_b=\mathrm{exp}\{\overline{a}ba\overline{b}\}`$;
3. $`T_a`$ does not in general leave the ground state invariant, i.e. invariance is maintained only if $`Na`$ is trivial, that is $`a=(n_1L_1+in_2L_2)/N`$;
4. the formal infinitesimal generators of $`T_a`$, namely $`𝔭_1=izi(+\overline{})`$ and $`𝔭_2=izi(\overline{})`$ do not leave the space of sections (Eq. 3) invariant.
To begin with the last statement, it is clear that we may consider the linear combinations $``$ and $`z\overline{}`$, neither of which is such as to transform sections into sections. From the group point of view, let $`\mathrm{}`$ be a translation in $`_2`$, i.e. $`\mathrm{}=k_1L_1+ik_2L_2,k`$. Let us consider $`T_a\psi `$. We find
$$\begin{array}{cc}\hfill (T_a\psi )(z+\mathrm{})& =\mathrm{exp}\{\overline{a}(z+\mathrm{})\frac{1}{2}|a|^2\}\mathrm{exp}\{\overline{\mathrm{}}(za)+\frac{1}{2}|\mathrm{}|^2\}\psi (za)\hfill \\ & =(T_a\psi )(z)\mathrm{exp}\{\overline{\mathrm{}}z+\frac{1}{2}|\mathrm{}|^2+\overline{a}\mathrm{}a\overline{\mathrm{}}\}.\hfill \end{array}$$
We conclude that a translated section satisfies boundary conditions with a different choice of the constants $`\delta _1,\delta _2`$, hence the bundle structure is not invariant under translation, except for
$$\overline{a}\mathrm{}a\overline{\mathrm{}}=2i\mathrm{}\{\overline{a}\mathrm{}\}2\pi i$$
which occurs precisely when $`a=(n_1L_1+in_2L_2)/N`$ ($`\mathrm{}\{\overline{a}\mathrm{}\}=(n_1k_2n_2k_1)L_1L_2/N=(n_1k_2n_2k_1)\pi `$, by Dirac’s quantization).
## 7. Conclusions
The problem of a constant magnetic field transversal to a torus raises the problem of translational symmetry. By quantizing the system according to the standard mathematical formulation of gauge theory we have shown that the symmetry is broken to $`Z_N\times Z_N`$. The conclusion to which one is led by this result is that the ambiguity in quantization, namely the two arbitrary phases $`\delta _1,\delta _2`$, entering in the definition of the domain of the Hamiltonian operator, represent some physical degree of freedom of the magnetic charge distribution generating the uniform field on the torus: monopole charges have, so to speak, horns. The effect is purely quantum mechanical and we empirically established that it vanishes approximately as $`\mathrm{exp}\{O(B/\mathrm{})\}`$. The mathematical roots of the result are the classic theorems of Weil (see \[Wei58\], Ch.VI, n.3, Prop.3); a thorough study of $`\theta `$-functions can be found in \[Dub81\].
Acknowledgments
I thank warmly P. Maraner and C. Destri for interesting discussions and for driving my attention to refs.\[Fub92, Dub81\].
## Appendix A A group-theoretical Wintner’s theorem
Wintner’s theorem (see \[Put67\]) states that the identity operator in a Hilbert space cannot be the commutator of two bounded operators. There is a poor’s man version of the theorem. Let $`U(a)`$ and $`V(b)`$ be unitary operators satisfying the canonical commutation relations (at the group level)
$$U(a)V(b)U(a)V(b)=e^{\overline{a}ba\overline{b}},(a,b).$$
Then $`U`$ and $`V`$ cannot be finite dimensional matrices.
Proof: just evaluate the determinant of both sides to get
$$1=\mathrm{exp}\{2iN\mathrm{}(\overline{a}b)\}\text{with }N=dim(U)\text{ }.$$
This is a contradiction, since the r.h.s. can assume any value on the unit circle. This last equation shows that we may take $`a`$ and $`b`$ in a finite subgroup and preserve the commutation relation: let $`Z_N=\{(n_1L_1+in_2L_2)/N|n_i\}`$; then the condition is satisfied precisely if $`L_1L_2=N\pi `$.
## Appendix B Dirac-Weil-Kostant-Souriau quantization condition
A general theorem (\[Hir78\], Th.21.1) relates the dimension of spaces of closed holomorphic forms on complex vector bundles to geometrical objects, namely Chern and Todd classes of the base space and of the bundle. In the simple case of a line bundle (fibre equal to $``$) over a complex two dimensional Riemann surface the theorem reduces to a simple result which has a very intuitive flavour from the point of view of Geometric Quantisation: the dimension of the physical Hilbert space coincides with the volume of phase space in units $`\mathrm{}`$ plus a constant given by half the Euler characteristic of the surface. This in turn implies that the volume of phase space must be an integer. We report here what appears to be the simplest proof, covering Dirac’s quantization condition, combining ideas from \[Alv85\] and \[SW76\]. Let us build a triangulation of the surface $``$ with vertices $`\alpha ,\beta ,\gamma ,\mathrm{}`$. Let $`𝒰_\alpha `$ denote the union of all triangles having $`\alpha `$ as vertex. By taking a sufficiently fine mesh, non empty intersections $`𝒰_\alpha 𝒰_\beta `$ consist of the union of two triangles which share the side $`\alpha \beta `$. A gauge field on $``$ is given by a closed two-form $`𝑩`$; in each local chart $`𝒰_\alpha `$ we define a potential $`𝑨_\alpha `$ such that $`𝑩=d𝑨_\alpha `$ in $`U_\alpha `$. According to Poincaré Lemma, for neighboring local charts we have
$$𝑨_\alpha 𝑨_\beta =d\chi _{\alpha \beta },$$
with differentiable transition functions $`\chi _{\alpha \beta }`$ which are antisymmetric in their indices. On triple intersections (any triangle $`𝒰_\alpha 𝒰_\beta 𝒰_\gamma `$) we have
$$𝑨_\alpha 𝑨_\beta =d\chi _{\alpha \beta },𝑨_\beta 𝑨_\gamma =d\chi _{\beta \gamma },𝑨_\gamma 𝑨_\alpha =d\chi _{\gamma \alpha }$$
It follows that $`c_{\alpha \beta \gamma }\chi _{\alpha \beta }+\chi _{\beta \gamma }+\chi _{\gamma \alpha }`$ is constant on each triangle<sup>5</sup><sup>5</sup>5It is useful to regard the relation between $`𝑨`$, $`\chi `$ and $`c`$ in terms of the coboundary operator: $`(\delta 𝑨)_{\alpha \beta }=d\chi _{\alpha \beta }`$, $`c_{\alpha \beta \gamma }=(\delta \chi )_{\alpha \beta \gamma }`$. See \[Alv85\].. Let us introduce a line bundle associated to $`𝑩`$: it is given locally by a direct product $`𝒰_\alpha \times `$ in such a way that in any overlap the complex fibres are connected by<sup>6</sup><sup>6</sup>6In the physical application the phase is $`e\chi /\mathrm{}c`$.
$$\zeta _\alpha =\zeta _\beta \mathrm{exp}\{i\chi _{\alpha \beta }\}$$
For consistency, on any triple overlap it must hold
$$\mathrm{exp}\{i\chi _{\alpha \beta }+i\chi _{\beta \gamma }+i\chi _{\gamma \alpha }\}=1$$
which implies thar $`c_{\alpha \beta \gamma }`$ must be an integer multiple of $`2\pi `$. This is usually referred as Weil theorem on holomorphic vector bundles (\[Wei58\], Prop.1, Ch.V, N.4).
The key result for our purposes is the following:
Theorem: the integral $`_{}𝑩`$ coincides with the discrete sum $`_\mathrm{\Delta }c_\mathrm{\Delta }`$ where $`\mathrm{\Delta }`$ runs over all triangles of the mesh.
Proof: The following purely algebraic identity holds (\[SW76\], p.131):
$$\begin{array}{c}_{\mathrm{\Delta }_{\alpha \beta \gamma }}𝑩=\frac{1}{3}_{\mathrm{\Delta }_{\alpha \beta \gamma }}\left(𝑨_\alpha +𝑨_\beta +𝑨_\gamma \right)\\ \hfill =\frac{1}{3}\left[(\chi _{\alpha \beta }+\chi _{\beta \gamma }+\chi _{\gamma \alpha })(\alpha )+(\chi _{\alpha \beta }+\chi _{\beta \gamma }+\chi _{\gamma \alpha })(\beta )+(\chi _{\alpha \beta }+\chi _{\beta \gamma }+\chi _{\gamma \alpha })(\gamma )\right]\\ \hfill \frac{1}{2}\left\{(\chi _{\alpha \beta }(\alpha )+\chi _{\alpha \beta }(\beta ))+(\chi _{\beta \gamma }(\beta )+\chi _{\beta \gamma }(\gamma ))+(\chi _{\gamma \alpha }(\gamma )+\chi _{\gamma \alpha }(\alpha ))\right\}\\ \hfill +\frac{1}{2}\left\{_{\alpha \beta }(𝑨_\alpha +𝑨_\beta )+_{\beta \gamma }(𝑨_\beta +𝑨_\gamma )+_{\gamma \alpha }(𝑨_\gamma +𝑨_\alpha )\right\}\end{array}$$
Notice that the terms in curly brackets average to zero when we sum over the whole triangulation, while the terms in square brackets are precisely the cocycle $`c_\mathrm{\Delta }`$, whose value is constant on the triangle. Hence we get
$$_{}𝑩=\underset{\mathrm{\Delta }}{}c_\mathrm{\Delta }$$
and as a result the flux of $`𝑩`$ is quantized.
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# Charge diffusion constant in hot and dense hadronic matter - A Hadro-molecular-dynamic calculation -
## 1 Introduction
Searching for the quark gluon plasma (QGP) is one of the hottest topics in the recent high energy nuclear/particle physics . Many kinds of phenomena have been proposed as candidates for the experimental signal to find the QGP. Charge fluctuation belongs to one of the most promising signals. Occurrence of Disoriented Chiral Condensation (DCC) at QCD phase transition would cause the large fluctuations of the ratio of the numbers of charged pions and neutral pions . Difference of the fluctuation intensity between QGP phase and hadronic phase can work as a signal of the existence of QGP . However, even if QCD phase transition takes place and the characteristic fluctuations of QGP or QCD phase transition are produced, such a fluctuation can be wiped out during hadronic era before the freeze-out. Whether the fluctuation which is caused at the phase transition or produced in the high temperature phase can survive in the hadronic era or not is the diffusion problem or the transport problem in the hot and dense hadronic matter. Macroscopic phenomenological equations, e.g., Navier-Stokes equation and diffusion equation, enable us to describe such phenomena in a simple manner. However, those kinds of phenomenological equations contain material constants, so called transport coefficients, and the dynamical calculation of the transport coefficients is a very important and difficult problem of statistical mechanics.
Interactions between hadrons are strong interaction which are believed to be described by QCD. However, in the hadronic energy region, perturbation does not work for QCD and no systematic way to treat is established. In a previous paper , we have reported a calculation of baryon number diffusion constant based on the relativistic collision event generator URASiMA (Ultra-Relativistic AA collision Simulator based on Multiple scattering Algorithm) which can reproduce hadronic spectra of nuclear collisions in the BNL-AGS and CERN-SPS. Usually collision event generators are so designed as to be suitable for the description of multiple production, and the detailed balance between the interactions is used to be paid only little attention . As we discussed in the ref. , we improved the URASiMA to recover detailed balance in the hadronic time scale and we have succeeded to establish stationary states of interacting hadrons with fixed temperatures and fixed baryon number densities . In this paper we evaluate charge diffusion constant in the hot and dense hadronic matter and investigate the diffusion of the charge fluctuation in the relativistic nuclear collisions.
## 2 Statistical ensembles
In order to prepare the statistical ensembles for the hot and dense hadronic state, we put hadrons in the box and updated with numerical code URASiMA with periodic boundary conditions. Recipes for the simulation are quite similar to the ordinary molecular dynamics. Molecular dynamics is one of the well established numerical methods in the statistical physics, however, it has been developed mainly for nonrelativistic system where particle number is conserved. The hot and dense hadronic system in which we have interest is a fully relativistic system, and the particle production and decay occur naturally. Therefore, the existence of stationary state itself is not apparent in this system.
URASiMA is originally designed as an event generator for the relativistic nuclear collisions based on the hadronic multichain model (MCM). In the high energy collision experiment, multiple production takes place and the system is thought to expand quickly. Therefore, usually, production processes play essential roles in the relativistic event generator but the reabsorption processes (reversal process of multiple production process) do not have been thought to be important, and sometimes have been neglected in the simulation code. However detailed balance of interactions is very important for statistical physics and naive application of collision event generator to the molecular dynamics in the box leads to the one-way conversion of the energy into particle production. As a result, Hagedorn type behavior appears, i.e., strange saturation of the temperature occurs .
In order to recover the detailed balance, we have improved URASiMA to contain many resonances and changed the code so as to describe some of production processes occur through production and decay of the resonance. The reversal processes of those processes have been naturally taken into account. As a result, after initial thermalization time period of about 150 fm$`/c`$, detailed balances among interactions seem to be almost kept during the simulation. It is noted that such improvement does not spoil the descriptive power for ultrarelativistic nuclear collisions . As we have already reported in ref. , slope parameter $`T`$ of the energy distribution (Table 1),
$$\frac{dN}{d^3𝒑}=\frac{dN}{4\pi EpdE}=C\mathrm{exp}(E/T),$$
became almost common value for all particles (fig. 1) and the population of particles became stationary (fig. 2). Therefore, we looked upon the system as the equilibrium state with temperature $`T`$ and fixed baryon number density.
===========
fig. 1
===========
===========
fig. 2
===========
Running URASiMA many times with the same energy and the same number of baryons, $`N_B`$, in the box with volume $`V`$, and taking equilibrium configurations, we have prepared statistical ensembles of the state with temperature $`T`$ and baryon number density $`N_B/V`$. Throughout this paper, we assumed iso-symmetry; the same number of protons and neutrons are put at initial time.
## 3 Diffusion constant in the linear response
According to Kubo’s Linear Response Theory, diffusion constant $`D`$ is obtained current (velocity) correlation . Because of the relativistic property of the hot and dense hadronic state, we should use $`𝜷=𝒑/E`$ instead of usual $`𝒗`$,
$$D=\frac{1}{3}_0^{\mathrm{}}<𝜷(t)𝜷(t+t^{})>𝑑t^{}c^2,$$
(1)
with $`c`$ being the velocity of light. In the calculation of the charge diffusion, average is taken over all charged particles. When we evaluated baryon number diffusion constant $`D_B`$ in our previous paper, we took average over only baryons. If correlation of velocities damps exponentially,
$$<𝜷(t)𝜷(t+t^{})>\mathrm{exp}(\frac{t^{}}{\tau })$$
(2)
with $`\tau `$ being relaxation time, we can rewrite eq. (1) with the simple form as,
$`D`$ $`=`$ $`{\displaystyle \frac{1}{3}}<𝜷(t)𝜷(t)>c^2\tau ,`$ (3)
$`=`$ $`{\displaystyle \frac{1}{3}}<\left({\displaystyle \frac{𝒑(t)}{E(t)}}\right)\left({\displaystyle \frac{𝒑(t)}{E(t)}}\right)>c^2\tau ,`$ (4)
with use of the relaxation time.
===========
fig. 3
===========
Figure 3 shows the correlations of the velocity of charged particles and exponential damping seems very good approximation. Diffusion constant obtained through eq. (3) is shown in fig. 4.
===========
fig. 4
===========
Diffusion constant increases with temperature and is almost independent of baryon number density. This result shows clear contrast against baryon number diffusion in our previous paper. Baryon number diffusion constant shows baryon number dependence and changes with temperature only weakly. Both calculations are almost the same but main contribution to each transport is different and we can understand fig. 4 as follows: In charge diffusion, main contribution comes from charged pions of which mass is just comparable to the temperature of the systems. On the other hand, the mass of the baryon is much larger than the temperature and baryonic system is almost nonrelativistic. The number of the baryon (baryonic charge carrier) is determined by the baryon number of the system, and even in very low temperature region, baryons must exist because of the fixed baryon number. In the lower temperature than pion mass, pion degrees of freedoms are frozen and electromagnetic charge and baryonic charge are both carried by baryons. At temperature about pion mass, pion degrees of freedom start to be melted and start to contribute to the charge transfer. In the higher temperatures, pions dominate the charge transfer. As a result, in the temperature lower than pion mass, both diffusion constants are almost the same, but in the temperature higher than pion mass, diffusion constant of charge current gradually increases with temperature.
## 4 Charge fluctuation in ultrarelativistic nuclear collisions
Using the above obtained values, let us sketch the diffusion of the charge fluctuation in the relativistic nuclear collisions. In the simplest static picture, the solution of diffusion equation,
$$\frac{}{t}f(𝒙,t)=D^2f(𝒙,t),$$
(5)
with the initial distribution,
$$f(𝒙,t_0)=\left(\sqrt{\frac{1}{2\pi R_0^2}}\right)^3\mathrm{e}^{\frac{(𝒙𝒙_0)^2}{2R_0^2}}$$
(6)
is given by
$$f(𝒙,t)=\left(\sqrt{\frac{1}{\pi (2R_0^2+4D(tt_0))}}\right)^3\mathrm{e}^{\frac{(𝒙𝒙_0)^2}{2R_0^2+4D(tt_0)}}.$$
(7)
Suppose that charge fluctuation with size $`R_0`$ produced at $`t=t_0`$, subjecting to eq. (5), the fluctuation diffuses and expands to $`R(t)=\sqrt{R_0^2+2D(tt_0)}`$. Therefore, we may regard that at such time $`t`$ that $`2D(tt_0)=R_0^2`$, charge fluctuation almost disappears. If the charge fluctuation about $`R_0=3`$ fm at $`T=T_C`$ 160 MeV is produced, it survives only 2 fm/c during hadronic matter era. But if charge fluctuation is produced at lower temperature about 120 MeV (Super Cooled DCC), diffusion constant is about 0.7 fm$`c`$ and the fluctuation with initial size of about 3 fm will survive for 6 fm/c. The length of hadronic era and the chronological change of the temperature depend on the solution of the hydrodynamical model, but the latter case is promising to be observed experimentally.
In order to take account of expansion of the hadronic fluid, we must use relativistic hydrodynamical model. In the relativistic hydrodynamical model, charged current $`J^\mu `$ is given as,
$$J^\mu =f(𝒙,t)U^\mu +J_d^\mu ,$$
(8)
with $`U^\mu `$ being local four velocity . $`f(𝒙,t)`$ is charge density on the local rest frame, $`f(𝒙,t)=U^\mu J_\mu `$. The diffusion current, $`J_d^\mu `$ is given by relativistic extension of the Fick’s law<sup>5</sup><sup>5</sup>5In the relativistic notation, $`x^0=ct`$, and $`\frac{}{x^0}=\frac{}{ct}`$.
$$J_d^\mu =\frac{D}{c}\mathrm{\Delta }_{}^{\mu }{}_{\nu }{}^{}^\nu f(𝒙,t),$$
(9)
where $`\mathrm{\Delta }_{}^{\mu }{}_{}{}^{\nu }=g^{\mu \nu }U^\mu U^\nu `$ is space-like projection operator orthogonal to $`U^\mu `$. Putting above $`J^\mu `$ into the continuity equation, $`_\mu J^\mu =0`$, we can obtain the relativistic diffusion equation,
$$U^\mu _\mu f(𝒙,t)+f(𝒙,t)_\mu U^\mu +\frac{D}{c}_\mu \mathrm{\Delta }_{}^{\mu }{}_{\nu }{}^{}^\nu f(𝒙,t)=0.$$
(10)
In the case of 1+1 dimensional scaling expansion, particular solution of the local four velocity is given as $`U^\mu =\frac{x^\mu }{c\tau }`$ with light-like coordinate, $`ct=c\tau \mathrm{cosh}\eta `$, $`z=c\tau \mathrm{sinh}\eta `$ , and diffusion equation (10) is rewritten as,
$$\frac{}{\tau }f(\eta ,\tau )+\frac{1}{\tau }f(\eta ,\tau )=\frac{D}{c^2\tau ^2}\frac{^2}{\eta ^2}f(\eta ,\tau ).$$
(11)
The solution of eq. (11) with initial condition, $`f(\eta ,\tau _0)=\sqrt{\frac{1}{2\pi R_0^2}}\mathrm{exp}(\frac{\eta ^2}{2R_0^2})f_0,`$ is given by
$$f(\eta ,\tau )=\frac{\tau _0}{\tau }\sqrt{\frac{1}{2\pi R(\tau )^2}}\mathrm{e}^{\frac{\eta ^2}{2R(\tau )^2}}f_0$$
(12)
with $`R(\tau )^2`$ being <sup>6</sup><sup>6</sup>6Note that this is the difussion in the $`\eta `$ direction and $`R`$ is the size in the $`\eta `$ coordinate. ,
$$R(\tau )^2=R_0^2+2\frac{D}{c^2\tau \tau _0}(\tau \tau _0).$$
(13)
The first factor in eq. (12), $`\frac{1}{\tau }`$, is the result of systematic expansion of the scaling solution; the region with fixed $`\mathrm{\Delta }\eta `$ expands as $`\tau \mathrm{\Delta }\eta `$. This rapid expansion also dominates the evolution of the fluctuation size, $`R(\tau )^2`$; the region with fixed $`\mathrm{\Delta }\eta `$ expands with $`\tau `$ but the diffusion effect subject to eqs. (5) or ( 10) is the expansion with only $`\sqrt{2D\tau }`$. Therefore, in $`\eta `$ coordinate, the effect of the diffusion becomes smaller and smaller with time, $`\frac{D}{c^2\tau \tau _0}`$. According to our result (fig. 4), diffusion constant in the hadronic matter is at most 2 fm$`c`$ and the diffusion term in eq. (13) becomes small enough already at $`\tau _0`$= several fm. As a result, the fluctuation in $`\eta `$ will be deformed only little by the diffusion effect in the hadronic era.
## 5 Concluding remarks
We evaluated charge diffusion constant of hot and dense hadronic matter based on the relativistic collision event generator URASiMA. Obtained charge diffusion constants increase from 0.5 fm$`c`$ to 2 fm$`c`$ in the temperature range between 80 MeV and 200 MeV. Based on the obtained diffusion constant, we made rough sketches of the diffusion of charge fluctuation in the ultrarelativistic nuclear collisions. For more improved discussion, we need to solve the hydrodynamical equation coupled with the charge current conservation equation .
Acknowledgment
The authors would like to thank prof. M. Namiki for his fruitful comments. They also would like to thank Dr. K. Homma for drawing our attention to this topic. Discussions with Prof. A. Nakamura and Prof. S. Daté were quite valuable. This work is supported by Grant-in-Aid for scientific research number 11440080 by Ministry of Education, Science, Sports and Culture, Government of Japan (Monbusho). Calculation has been done at Institute for Nonlinear Sciences and Applied Mathematics, Hiroshima University.
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# Percolation and Deconfinement in SU(2) Gauge Theory
## 1 INTRODUCTION
The study of critical phenomena has always been one of the most challenging and fascinating topics in physics. One can give many examples of systems which undergo phase transitions, from familiar cases like the boiling of water in a kettle or the paramagnetic-ferromagnetic transition of iron at the Curie temperature, to the more complicated case of the transition from hadronic matter to quark-gluon plasma which is likely to be obtained by high-energy heavy-ion experiments in the coming years. Particularly interesting are the second-order phase transitions, characterized by a continuous variation of the order parameter and a divergent correlation length at the threshold. Already in the 70’s one tried to use percolation theory in order to give a geometrical description of the critical behaviour of dynamical systems undergoing second-order phase transitions . Percolation has in fact several features in common with such systems: power law behaviour of the variables at criticality, scaling relations, universality.
The critical behaviour of the Ising model can indeed be reformulated in terms of percolation theory: magnetization sets in when suitably defined clusters of parallel spins reach the dimensions of the system . In particular, the critical exponents for percolation then become equal to the Ising exponents.
The paramagnetic-ferromagnetic transition of the Ising model is strictly related to the confinement-deconfinement transition in finite temperature $`SU(2)`$ gauge theory. This was conjectured on the basis of effective theories and confirmed by lattice studies . The order parameter for $`SU(2)`$ is the lattice average of the Polyakov loop, and it behaves like the magnetization in the Ising model; in particular, the critical exponents are the same.
These analogies are the basis of this work. We have tried to see whether it is possible to find a description of critical behaviour in terms of percolation also for the deconfinement transition in $`SU(2)`$ gauge theory. We show that in a lattice regularization which corresponds to the strong coupling limit, both in two and in three space dimensions, the percolation of Polyakov loop clusters (taken to be suitably defined areas of Polyakov loops $`L`$ of the same sign) leads to the correct deconfinement temperature and to the correct critical exponents for the deconfinement. We have achieved this result by means of computer simulations of finite temperature $`SU(2)`$ gauge theory and using standard finite size scaling tecniques to extrapolate the results to the infinite volume limit.
## 2 PERCOLATION THEORY AND THE ISING MODEL
The percolation problem is easy to formulate: just place randomly pawns on a chessboard. Regions of adjacent pawns form clusters. Percolation theory deals with the properties of these clusters when the chessboard is infinitely large. If one of the clusters spans the chessboard from one side to the opposite one, we say that the cluster percolates.
Quantitatively, one counts how many pawns belong to each cluster and calculates two quantities:
$``$ The average cluster size S, defined as:
$$S=\underset{s}{}\left(\frac{n_ss^2}{_sn_ss}\right).$$
(1)
Here $`n_s`$ is the number of clusters of size $`s`$ and the sums exclude the percolating cluster; this number indicates how big on average the clusters are which do not percolate.
$``$ The percolation strength P, defined as:
$$P=\frac{\text{size of the percolating cluster}}{\text{no. of lattice sites}}.$$
(2)
By varying the density of our pawns, a kind of phase transition occurs. We pass from a phase of non-percolation to a phase in which one of the clusters percolates. The percolation strength P is the order parameter of this transition: it is zero in the non-percolation phase and is different from zero in the percolation phase.
If we want to study percolation in dynamical systems, for example the Ising model, first of all we must define the rule to build up the clusters (for instance we can join together all nearest-neighbour spins of the same sign). Then one has to find out at which temperature we have a spanning cluster. Let us call this temperature $`T_p`$. It turns out that for $`TT_p`$ the percolation variables P and S behave in the following way:
$$P(T_pT)^{\beta _p},S|TT_p|^{\gamma _p}$$
(3)
A. Coniglio and W. Klein demonstrated that for some special definition of cluster, the onset of percolation coincides with the one of magnetization; besides, $`\beta _p`$ and $`\gamma _p`$ coincide respectively with the magnetization exponent $`\beta `$ and the susceptibility exponent $`\gamma `$. Such cluster definition had already been used by Fortuin and Kasteleyn to show that the partition function of the Ising model can be rewritten in purely geometrical terms as a sum over clusters configurations . According to the Fortuin-Kasteleyn prescription, two nearest-neighbouring spins of the same sign belong to the same cluster with a probability $`p=1exp(2\beta )`$ ($`\beta =J/kT`$, J is the Ising coupling, T the temperature). The result of Coniglio and Klein is valid in any space dimension $`d`$ ($`d2`$) and it is independent of the lattice geometry (square, triangular, honeycomb, etc.) as long as it is homogeneous.
## 3 POLYAKOV LOOP PERCOLATION IN SU(2) GAUGE THEORY
Finite temperature $`SU(N)`$ gauge theories describe the interactions of systems of mutually interacting gluons in thermal equilibrium ($`N`$ is the number of colours). Systems containing quarks together with gluons are certainly more appealing; such a scenario will be explored experimentally up to very high temperatures by means of high-energy heavy ion collisions. Nevertheless, $`SU(N)`$ gauge theories are of theoretical interest; already at this level one has a transition between confinement to deconfinement, going from a phase of bound state of gluons (glueballs) to a phase of free gluons. The order parameter of the confinement-deconfinement phase transition of $`SU(N)`$ gauge theories is the lattice average of the Polyakov loop $`<L>`$. It is related to the effective potential $`V(r)`$ of a static (mass $`\mathrm{}`$) quark-antiquark pair put in the gluons’ medium at temperature $`T`$ at a distance $`r`$ from each other, when $`r`$ is very big:
$$|<L>|^2\underset{r\mathrm{}}{lim}e^{\frac{V(r)}{T}},$$
(4)
The Polyakov loop is zero in the confined phase ($`V(r)\mathrm{}`$) and different from zero in the deconfined phase ($`V(r)`$ finite). Some time ago it was conjectured by B. Svetitsky and L. G. Yaffe that the critical behaviour of $`SU(N)`$ gauge theories has a strong relationship with the one of simple $`Z(N)`$ spin models, with which they share a common global symmetry . In particular, in case of second order phase transitions, both models would belong to the same universality class, that is they would have the same set of critical exponents. One simple test of the Svetitsky-Yaffe conjecture is provided by the $`SU(2)`$ gauge theory: numerical simulations showed that its critical exponents indeed coincide with the ones of the Ising model .
In some respect the lattice configurations of $`SU(2)`$ are similar to the ones of the Ising model. Instead of having a spin variable at each lattice site we have the value of the Polyakov loop at that site. Our aim is to see whether it is possible to build clusters of Polyakov loops such that their percolation indices (threshold, exponents) coincide with the thermal ones. In order to do that, we have to face two difficulties:
i) the Polyakov loop is not a two-valued variable like the spin in the Ising model but a continuous one; its values vary in a range that, with the normalization convention we use, is $`[1,1]`$;
ii) the $`SU(2)`$ Lagrangian is not directly a function of the Polyakov loop, therefore it is not possible to extract from it the expression of the bond probability that we need to build the clusters like in the Ising model.
In a recent work it was proved that the point i) is not a problem. If we take an Ising model with continuous instead of two-valued spins, the definition of Fortuin-Kasteleyn clusters can be generalized by taking as a bond probability between two positive (or negative) nearest-neighbour spins $`s_i`$ and $`s_j`$ the following expression:
$$p(i,j)=1exp(2\beta s_is_j)$$
(5)
The problem ii), however, is hard to overcome. In particular, it seems that there is no way to express the lattice Lagrangian of $`SU(2)`$ in terms of Polyakov loops for all lattice regularizations. Therefore we were forced to investigate $`SU(2)`$ in a special lattice regularization, which corresponds to the so called strong coupling limit. In this case it was shown that the partition function of $`SU(2)`$ can be written in a form which, apart from a factor which depends on the group measure, is the partition function of the continuous Ising model studied in . The coupling $`\kappa `$ of the spin model and the coupling $`\beta `$ of $`SU(2)`$ are related to each other by this relation:
$$\kappa (\beta /2)^{N_t}$$
(6)
($`N_t`$ fixes the temporal lattice regularization). If $`SU(2)`$ is approximately a continuous Ising model, we can try to use the general definition of clusters of . Therefore we have defined our clusters as regions of like-sign Polyakov loops connected by bonds distributed according to the bond weight
$$p(i,j)=1\mathrm{exp}(2\kappa L_iL_j),$$
(7)
($`L_i`$ and $`L_j`$ are the Polyakov loop values at the sites $`i`$ and $`j`$). For $`N_t=2`$ the approximation of is good and we chose to investigate numerically this special case in two and three space dimensions.
## 4 RESULTS
We began our percolation studies performing some test runs for different lattice sizes to check the behaviour of our percolation variables around criticality. Figure 1 shows the average cluster size $`S`$ for $`2+1SU(2)`$ in correspondence of different lattice sizes. To get the critical point of the percolation transition we used the method suggested in . For a given lattice size and a value of $`\beta `$ we counted how many times we found a percolating cluster. This number is successively divided by the total number of configurations at that $`\beta `$. We call this quantity percolation probability. This variable is directly a scaling function, analogue to the Binder cumulant in continuous thermal phase transitions. Figure 2 shows the percolation probability for the 3-dimensional case as a function of $`\beta `$ for $`24^3\times 2`$, $`30^3\times 2`$ and $`40^3\times 2`$.
The lines cross at the same point within the errors and that restricts further our $`\beta `$ range for the critical threshold. Besides, since the percolation probability is a scaling function, we could already get clear indications about the class of critical exponents of our clusters. In fact, if one knows the critical point and the exponent $`\nu `$, a rescaling of the percolation probability as a function of $`(\beta \beta _{cr})L^{1/\nu }`$ should give us the same function for each lattice size. Fig. 3 and 4 show the rescaled percolation probability for 3-dimensional $`SU(2)`$ using $`\beta _{cr}=1.8747`$ and two different values of the exponent $`\nu `$, respectively the Ising value and the random percolation one. The figures show clearly a remarkable scaling for the Ising exponent and no scaling for the random percolation exponent.
To evaluate the exponents $`\beta `$ and $`\gamma `$ we performed high-statistics simulations in the range where the percolation probability curves cross each other. The number of measurements we took for each value of the coupling varies from 50000 to 100000. We used the $`\chi ^2`$ method to determine the values of the exponents. The final results are reported in Table 1 and 2. The agreement both in two and in three dimensions is good.
## 5 CONCLUSIONS
We have shown that the confinement-deconfinement phase transition in finite temperature $`SU(2)`$ pure gauge theory can be described as percolation of some suitably defined clusters of Polyakov loops of the same sign. Our result is valid only in the strong coupling limit, and in order to make it general a different approach seems to be necessary. The use of effective theories for $`SU(2)`$ may help to solve the problem .
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